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# Estimating mixed quantum states
## I Introduction
One of the best distinctions between classical and quantum systems can be formulated using the language of state measurements. The state of a single classical system is an observable: it is in principle always possible to measure generalized coordinates and generalized momenta at a certain time $`t`$. This information is complete in the sense that it allows us to calculate the classical state at any time if the forces are known.
Quantum mechanics shows a different picture. In general, the state of a single quantum system cannot be determined since any measurement will lead to a reduction of the quantum state. As a consequence the complete reconstruction of a quantum state is only possible if we measure specific observables on an infinite ensemble of identically prepared systems. Such measurements have been intensively discussed experimentally as well as theoretically in recent years .
The interest in such questions has been renewed in the field of quantum information. In particular, it has been investigated how well one can estimate a quantum state from a finite ensemble of identical systems. For pure spin states this problem has been solved . The intriguing result is that one can learn more about the finite ensemble by performing a measurement on all quantum systems simultaneously . In fact one can optimize the readout of quantum information via simultaneous strategies. Such a joint measurement may, however, pose practical problems since we may not have all systems available at a time or cannot realize the required complicated measurement operators experimentally. Hence we have recently analyzed how close one can come to such optimal joint estimation schemes by using single adaptive quantum measurements .
Also for mixed spin states the optimal generalized measurement has been constructed for a finite ensemble . In addition it has been shown how the gain of information of generalized measurements increases when the estimated state approaches a pure quantum state.
In the present paper we shall analyze the estimation of mixed states of qubits using single adaptive measurements. As a criteria for adaption the Kullback information gain turns out to be suited. We will compare our results to the optimal measurements discussed in .
The paper is organized as follows. In Sec. II we present the quantum system we are dealing with and introduce the corresponding notation. Sec. III contains the description of our adaptive measurement scheme. The optimization strategies used in these adaptive schemes are then described in detail in Sec. IV. In Sec. V and VI we show the estimation results for the cases of estimation without and with a priori information regarding the radial distribution of quantum states inside the Bloch sphere. We conclude with Sec. VII.
## II Quantum system and measurement operator
Suppose we are given $`N`$ two-level systems (qubits) identically prepared in the mixed state $`\widehat{\rho }`$. The task is to adaptively estimate $`\widehat{\rho }`$ by using experimentally realizable single measurements on the two-level systems.
Let us first define the notation and the key elements of our analysis. Any mixed quantum state $`\widehat{\rho }`$ of a two-level system (qubit) can be written down in the Bloch-sphere representation
$$\widehat{\rho }(r,\theta ,\varphi )=\frac{1}{2}\left(\widehat{1}+\stackrel{}{r}(r,\theta ,\varphi )\stackrel{}{\widehat{\sigma }}\right)$$
(1)
with the Pauli-spin-vector $`\stackrel{}{\widehat{\sigma }}=(\widehat{\sigma }_x,\widehat{\sigma }_y,\widehat{\sigma }_z)^T`$ and the Bloch vector
$$\stackrel{}{r}(r,\theta ,\varphi )=r\left(\begin{array}{c}\mathrm{sin}\theta \mathrm{cos}\varphi \\ \mathrm{sin}\theta \mathrm{sin}\varphi \\ \mathrm{cos}\theta \end{array}\right)$$
(2)
expressed in spherical coordinates. The parameters $`r[0,1]`$, $`\theta [0,\pi ]`$, $`\varphi [0,2\pi )`$ uniquely determine the quantum state inside the Bloch sphere. Using the matrix representation of $`\stackrel{}{\widehat{\sigma }}`$ the density matrix reads
$$\widehat{\rho }(r,\theta ,\varphi )=\frac{1}{2}\left(\begin{array}{cc}1+r\mathrm{cos}\theta & r\mathrm{sin}\theta e^{i\varphi }\\ r\mathrm{sin}\theta e^{i\varphi }& 1r\mathrm{cos}\theta \end{array}\right).$$
(3)
This parametrization also allows us to represent the estimated state $`\widehat{\rho }^{(est)}`$ obtained after a certain measurement sequence. By introducing the probability density $`w(r,\theta ,\varphi )`$ we can write
$$\widehat{\rho }^{(est)}=๐Vw(r,\theta ,\varphi )\widehat{\rho }(r,\theta ,\varphi )$$
(4)
with normalization
$$๐Vw(r,\theta ,\varphi )=1.$$
(5)
The integration
$$๐V=_0^1๐rr^2_0^\pi ๐\theta \mathrm{sin}\theta _0^{2\pi }๐\varphi $$
(6)
ranges over the whole sphere.
To guarantee the experimental realizability of our proposed measurement strategy we restrict ourselves to a simple class of measurements. This class consists of von Neumann measurements, e.g., polarization or spin measurements along a certain axis $`(\theta _n,\varphi _n)`$. The tunable parameters $`\theta _n`$ and $`\varphi _n`$ define the direction of the projection on the Bloch sphere surface for the $`n`$th measurement. The corresponding projection operator $`|\theta _n,\varphi _n\theta _n,\varphi _n|`$ with the state
$$|\theta _n,\varphi _n=\mathrm{cos}\frac{\theta _n}{2}|0+\mathrm{sin}\frac{\theta _n}{2}e^{i\varphi _n}|1$$
(7)
therefore defines two measurement results. Either the system described by $`\widehat{\rho }(r,\theta ,\varphi )`$ is polarized in the direction $`(\theta _n,\varphi _n)`$ or in the opposite direction given by $`(\pi \theta _n,\pi +\varphi _n)`$. We encode the first result by the number 1 and the second by the number 0.
The two possible outcomes of the measurement occur with the probabilities
$`P_1(r,\theta ,\varphi |\theta _n,\varphi _n)`$ $`=`$ $`\theta _n,\varphi _n|\widehat{\rho }(r,\theta ,\varphi )|\theta _n,\varphi _n`$ (8)
$`=`$ $`{\displaystyle \frac{1}{2}}[1+r\mathrm{cos}\theta \mathrm{cos}\theta _n`$ (9)
$`+`$ $`r\mathrm{sin}\theta \mathrm{sin}\theta _n\mathrm{cos}(\varphi \varphi _n)],`$ (10)
$`P_0(r,\theta ,\varphi |\theta _n,\varphi _n)`$ $`=`$ $`1P_1(r,\theta ,\varphi |\theta _n,\varphi _n)`$ (11)
$`=`$ $`{\displaystyle \frac{1}{2}}[1r\mathrm{cos}\theta \mathrm{cos}\theta _n`$ (12)
$``$ $`r\mathrm{sin}\theta \mathrm{sin}\theta _n\mathrm{cos}(\varphi \varphi _n)]`$ (13)
which clearly depend on the chosen measurement direction and on the measured quantum state.
## III Adaptive measurements
We now propose an adaptive measurement strategy to improve the estimation of a quantum state $`\widehat{\rho }(R,\mathrm{\Theta },\mathrm{\Phi })`$ from a finite ensemble of $`N`$ identically prepared quantum systems. Note that the Bloch vector coordinates $`R,\mathrm{\Theta }`$ and $`\mathrm{\Phi }`$ are the same for all $`N`$ systems. Despite the fact that we restrict ourselves to simple projection measurements on single quantum systems we will show that it is possible to improve the estimation quality by using an adaptive measurement strategy. This strategy is based on an algorithm, cf. Fig. 1, which consists of five steps:
1. We take the first, $`n=1`$, of the $`N`$ quantum systems and perform a measurement with randomly chosen direction $`(\theta _1,\varphi _1)`$.
2. The $`n`$th measurement along the direction $`(\theta _n,\varphi _n)`$ yields one of the two possible outcomes that we denote by $`i=0`$ and $`i=1`$. By using this information we modify the distribution $`w_{n1}(r,\theta ,\varphi )`$ of the estimated density operator
$$\widehat{\rho }_{n1}^{(est)}=๐Vw_{n1}(r,\theta ,\varphi )\widehat{\rho }(r,\theta ,\varphi )$$
(14)
after $`n1`$ steps according to Bayesโ rule
$$w_n(r,\theta ,\varphi )=Z^1P_i(r,\theta ,\varphi |\theta _n,\varphi _n)w_{n1}(r,\theta ,\varphi ).$$
(15)
Hence with the help of the probabilities $`P_i`$, Eq. (8), we update our current knowledge about the finite ensemble as formulated by the distribution $`w_n`$. The normalization constant reads
$`Z={\displaystyle ๐VP_i(r,\theta ,\varphi |\theta _n,\varphi _n)w_{n1}(r,\theta ,\varphi )}.`$ (16)
Before we have acquired any information about the system, we assume our knowledge to be homogeneously distributed over the Bloch sphere; that is, we start from the initial distribution
$$w_0(r,\theta ,\varphi )=\frac{3}{4\pi },$$
(17)
thereby assuming that each volume element $`dV=r^2\mathrm{sin}\theta drd\theta d\varphi `$ is equally probable.
3. The updated probability distribution $`w_n(r,\theta ,\varphi )`$ describes our present knowledge about the quantum state. With its help we determine the next measuring operator, i.e., the direction $`(\theta _{n+1},\varphi _{n+1})`$ of the measurement step $`n+1`$. The new measuring operator should be designed in such a way that it allows us to gain the maximum amount of additional information about the unknown quantum state $`\widehat{\rho }`$. For this purpose we have to apply a criterion by which we quantify the notion of maximum information gain. The different criterions that we use in this context will be described in detail in the next section. This step reflects the adaptive aspect of our algorithm, because the choice of a measuring operator is based on $`\widehat{\rho }_n`$ and thereby on the history of all previous measurement outcomes.
4. Once having found the next measuring operator $`|\theta _{n+1},\varphi _{n+1}\theta _{n+1},\varphi _{n+1}|`$ we now take one of the remaining quantum systems and measure it. If we still have quantum systems left, we continue with step 2.
5. After we have used up all $`N`$ mixed qubits we arrive at the final probability distribution $`w_N(r,\theta ,\varphi )`$ which allows us to construct the corresponding estimated state $`\widehat{\rho }_N^{(est)}`$.
As the measure of our state estimation quality we will use the fidelity
$$F_N(\widehat{\rho },\widehat{\rho }_N^{(est)})=Tr^2\sqrt{\sqrt{\widehat{\rho }_N^{(est)}}\widehat{\rho }(R,\mathrm{\Theta },\mathrm{\Phi })\sqrt{\widehat{\rho }_N^{(est)}}},$$
(18)
for mixed quantum states which reduces to
$`F_N=`$ $`{\displaystyle \frac{1}{2}}`$ $`[1+\stackrel{}{r}_N^{(est)}\stackrel{}{r}(R,\mathrm{\Theta },\mathrm{\Phi })`$ (19)
$`+`$ $`\sqrt{1|\stackrel{}{r}_N^{(est)}|^2}\sqrt{1|\stackrel{}{r}(R,\mathrm{\Theta },\mathrm{\Phi })|^2}]`$ (20)
for two-level systems with Bloch vectors $`\stackrel{}{r}_N^{(est)}`$ and $`\stackrel{}{r}(R,\mathrm{\Theta },\mathrm{\Phi })`$. Note that this fidelity of course depends on the number $`N`$ of quantum systems at our disposal.
## IV Measurement Strategies
In this section we will describe the strategies that we have applied to find the direction $`(\theta _n,\varphi _n)`$ for the $`n`$th measurement by learning from the results of earlier measurements.
### A Random selection from all axes
A straightforward way to select a new measurement direction $`(\theta _n,\varphi _n)`$ is to choose the parameters $`\theta _n`$ and $`\varphi _n`$ randomly on the Bloch sphere, independent of any knowledge already acquired about the state. That is, each infinitesimal surface element $`\mathrm{sin}\theta _nd\theta _nd\varphi _n`$ occurs with the same probability $`1/4\pi `$. Clearly this strategy is not adaptive because $`(\theta _n,\varphi _n)`$ does not depend on any previous measurement outcome. Nevertheless, the estimated density operator $`\widehat{\rho }_{n1}^{(est)}`$ can still be updated after each measurement as described in step 2 of our algorithm.
The random selection implements a measurement protocol lacking any constructive strategy. Thus the results of this method will serve as a reference to which we can compare the outcomes of the optimization strategies described below.
### B Minimal measurements along three axes
In principle it is possible to determine any mixed state by measurements along three axes if an infinite number of quantum systems prepared in this state is available. Without loss of generality one can choose the Pauli operators $`\widehat{\sigma }_x`$, $`\widehat{\sigma }_y`$ and $`\widehat{\sigma }_z`$ to represent measurements along three orthogonal directions on the Bloch sphere. This set of measurement operators represents a minimal measurement or โ in other words โ corresponds to a minimal quorum.
Thus it is interesting to compare the efficiency of such a minimal measurement to our adaptive methods in the case of a finite number $`N`$ of available quantum systems. Then the measurement scheme consists of projecting an average number of $`N/3`$ systems using each of the directions $`(\pi /2,0)`$, $`(\pi /2,\pi /2)`$ and $`(0,0)`$.
### C Maximization of Kullback information gain
If we look at the measurement procedure from an information theoretic point of view then our aim will be to maximize the information that we can get in the measurement step from $`n1`$ to $`n`$. A measure for the average information gain in the $`n`$th measurement is the so called Kullback information that can be defined as
$`\overline{K}(\theta _n,\varphi _n)=`$ $`{\displaystyle \underset{i=0}{\overset{1}{}}}`$ $`p_i^{(est)}(\theta _n,\varphi _n)`$ (21)
$`\times `$ $`{\displaystyle ๐Vw_n^{(i)}(r,\theta ,\varphi )\text{ log}_2\frac{w_n^{(i)}(r,\theta ,\varphi )}{w_{n1}(r,\theta ,\varphi )}}`$ (22)
with
$$p_1^{(est)}(\theta _n,\varphi _n)=\theta _n,\varphi _n|\widehat{\rho }_{n1}^{(est)}|\theta _n,\varphi _n$$
(23)
and
$$p_0^{(est)}=1p_1^{(est)}$$
(24)
being the estimated probabilities for the outcomes $`i=0,1`$ based on our current knowledge, i.e., based on the density operator $`\widehat{\rho }_{n1}^{(est)}`$. Consequently, also the probability density
$$w_n^{(i)}=Z^1p_i^{(est)}(\theta _n,\varphi _n)w_{n1}$$
(25)
explicitly depends on the outcome $`i`$ and on the direction $`(\theta _n,\varphi _n)`$. Hence our expression, Eq. (15), for the Kullback information is a function of $`(\theta _n,\varphi _n)`$ which can be maximized. In order to see that $`\overline{K}`$ describes an estimated average information gain we rewrite it in the form
$`\overline{K}(\theta _n,\varphi _n)`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{1}{}}}p_i^{(est)}(\theta _n,\varphi _n)`$ (26)
$`\times `$ $`{\displaystyle ๐Vw_n^{(i)}(r,\theta ,\varphi )\text{ log}_2w_n^{(i)}(r,\theta ,\varphi )}`$ (27)
$``$ $`{\displaystyle ๐Vw_{n1}(r,\theta ,\varphi )\text{ log}_2w_{n1}(r,\theta ,\varphi )}`$ (28)
$`=`$ $`S_{n1}{\displaystyle \underset{i=0}{\overset{1}{}}}p_i^{(est)}(\theta _n,\varphi _n)S_n^{(i)}(\theta _n,\varphi _n).`$ (29)
The entropy $`S_{n1}`$ describes our knowledge before the measurement, whereas the entropy $`S_n^{(i)}`$ stands for the estimated entropy provided we find the result $`i`$.
Hence we can maximize the difference of entropies before and after the measurement by adjusting the parameters $`(\theta _n,\varphi _n)`$. We therefore select the measuring direction $`(\theta _n,\varphi _n)`$ that yields the maximum average information gain.
The estimated state $`\widehat{\rho }_N^{(est)}`$ is finally again determined from $`w_N`$. This strategy can be applied in the case that all measurement directions are possible as well as in the case that only the three measurement directions along $`\widehat{\sigma }_x`$, $`\widehat{\sigma }_y`$ and $`\widehat{\sigma }_z`$ are allowed. In the latter case the maximization described before is done only for the directions $`(\pi /2,0)`$, $`(\pi /2,\pi /2)`$ and $`(0,0)`$. We will discuss the resulting estimation precisions in the next section.
## V Estimation without a priori information
In this section we numerically evaluate the average fidelities for the state estimation schemes described above. Only such an average fidelity is a reasonable measure of quality of a specific estimation procedure, since we assume to have no prior information about $`\widehat{\rho }`$.
These average fidelities will depend on the number $`N`$ of identically prepared quantum systems that we have at our disposal. Thus one state estimation experiment consists of a sequence of $`N`$ measurements performed on $`N`$ identical systems in state $`\widehat{\rho }(R,\mathrm{\Theta },\mathrm{\Phi })`$ and a subsequent estimation of a mixed state $`\widehat{\rho }_N^{(est)}`$. The fidelity $`F_N`$, Eq. (14), of the state estimation is then calculated by comparing both states.
However, in order to get the average fidelity
$$F_N=F(\widehat{\rho }(R,\mathrm{\Theta },\mathrm{\Phi }),\widehat{\rho }_N^{(est)})_{\widehat{\rho }}$$
(30)
we have to perform such a single run of the (numerical) experiment over and over again for different initial states $`\widehat{\rho }(R,\mathrm{\Theta },\mathrm{\Phi })`$, i.e., for different coordinates $`(R,\mathrm{\Theta },\mathrm{\Phi })`$. Hence the initial states $`\widehat{\rho }(R,\mathrm{\Theta },\mathrm{\Phi })`$, Eq. (3), are chosen randomly from an isotropic and homogenous probability distribution $`3/(4\pi )`$ for each volume element $`R^2\mathrm{sin}\mathrm{\Theta }dRd\mathrm{\Theta }d\mathrm{\Phi }`$ of the Bloch sphere. Using this homogenous probability distribution ensures that the performance of an estimation strategy is not biased by any specific choice of initial states. For more details on the averaging procedure see Appendix A.
The average fidelity for each $`N`$ was obtained by averaging over $`10^4`$ experiments, i.e., $`10^4`$ initial states equally distributed inside the Bloch sphere. For the sake of a clear graphical presentation not the average fidelities themselves but the average errors
$$f_N=1F_N$$
(31)
are calculated for different $`N`$.
In Fig.2 the average errors $`f_N`$ are compared to the average error $`f_N^{rand}`$ of the random selection scheme by plotting the ratio
$$\gamma _N\frac{f_N}{f_N^{rand}}$$
(32)
versus $`N`$. This quantity shows the relative performance of the different schemes compared to the random selection scheme.
As we can see the scheme based on measurements along three axes improves the estimation quality by approximately 4% even if no adaptive strategy is applied. We can further decrease the estimation error if we use the Kullback information gain strategy. For this strategy the average errors are always smaller than for the non-adaptive schemes in both cases. This shows that one can indeed decrease the error of a state estimation by using adaptive algorithms. The errors decrease by approximately 3% for $`N>10`$ in the case of the all-axes scheme. And even in the case of only three possible measurement directions an optimization of the Kullback information gain yields an improvement of the estimation quality. Moreover, the resulting measurement strategy is the best strategy based on separate measurements that we found so far.
Now we analyze how our estimation schemes compare to the optimal ones that are based on collective measurements on all $`N`$ quantum systems. For pure qubits it was shown that the fidelities of single qubit measurement schemes can be always bigger than 98% of the optimal ones that use collective measurements. Can we also reach such values in the case of mixed qubits? It turns out that this is not the case. In Fig. 3 we show the ratios $`<F_N>/<F_N^{(opt)}>`$ with $`<F_N^{(opt)}>`$ being the fidelity of the optimal measurement scheme for different $`N`$. We find that for the random selection scheme we get fidelities which are in the worst case below 96% of the optimal values. Even for the best estimation scheme (measurement along three axes with Kullback information maximization) we are still about 3.5% off the optimal values. Also the convergence towards 1 is much slower for mixed qubits than in the case of pure ones . For mixed qubits we are still about 2% off the optimal fidelities at $`N=50`$ whereas for pure qubits the ratio of fidelities has already been very close to 1 at this point.
These results indicate that the advantage of collective measurements compared to single quantum system measurements increases with growing degree of mixing of the qubits. Our results therefore confirm the asymptotic findings by Gill and Massar who investigated the case of state estimation of large qubit ensembles. They showed that for large $`N`$ one can asymtotically achieve the precision of collective measurement schemes by measurements on single qubits if they are known to be in a pure state and that this is no longer true in the case of mixed qubits.
## VI Estimation with radial a priori information
Inspired by these findings we finally discuss how the precision of our estimation schemes depends on the degree of mixing of the initial states. In our case the adequate measure of precision is the estimation error $`f_N`$, Eq. (21). For a quantitative check of the relation between estimation quality and degree of mixing of the initial states we have chosen the initial states no longer according to a homogenous probability distribution inside the Bloch sphere but from a probability distribution
$$w_0(r,\theta ,\varphi ;\alpha )dV=\frac{\alpha +1}{4\pi }r^\alpha \mathrm{sin}\theta drd\theta d\varphi $$
(33)
which clearly depends on the radius $`r`$. Please note that we get the homogenous distribution for $`\alpha =2`$ again. A larger parameter $`\alpha 0`$ means a growing average radius
$$\overline{r}=\frac{\alpha +1}{\alpha +2}$$
(34)
which indicates a decreasing degree of mixing of the initial states. Thus a variation of the parameter $`\alpha `$ allows us to study the influence of the degree of mixing onto the estimation precision of our schemes in which we have of course also adapted the initial probability $`w_0`$.
However, our previously desribed estimation scheme does not converge towards the estimation scheme for pure states for $`\alpha \mathrm{}`$. The cause of this behaviour is the final readout of the estimated radius $`r^{(est)}`$. Please remember that this radius is found via the integration $`๐Vw_N(r,\theta ,\varphi )\widehat{\rho }(r,\theta ,\varphi )`$ which yields the estimated state $`\widehat{\rho }^{(est)}(r^{(est)},\theta ^{(est)},\varphi ^{(est)})`$. Thus even for solely pure initial states the estimated state will never be pure itself. That is, the readout of the estimated state does not properly account for our a priori information. For larger $`\alpha `$ the estimated radius will always be too small compared to the average radius $`\overline{r}`$, Eq. (24), of our initial distribution $`w_0`$, Eq. (23).
To overcome this non-convergence we now introduce an alternative readout scheme for the radius $`r^{(est)}`$. The estimated parameters $`\theta ^{(est)}`$ and $`\varphi ^{(est)}`$ are still obtained in the same way as before, that is, via Bloch representation of $`\widehat{\rho }^{(est)}`$. The estimated radius $`r^{(est)}`$, however, is now chosen as the average radius
$$r^{(est)}=๐Vrw_N(r,\theta ,\varphi )$$
(35)
of the final distribution $`w_N(r,\theta ,\varphi )`$. It is easy to show that โ even before the first measurement โ this estimated radius tends to 1 for large $`\alpha `$ and merges into the pure state estimation scheme for $`\alpha \mathrm{}`$.
Using this readout scheme we compare the resulting errors $`f_N`$ to the errors $`f_N^{(opt)}`$ that could be achieved by a optimal collective estimation scheme . The errors $`f_N`$ are based on a simple 3-axes estimation as described in Sec. IV.
In Fig. 4 we plotted the ratio $`f_N^{(opt)}/f_N`$ versus $`N`$. As one would expect our estimations are always worse than the optimal ones. We also find a very clear dependence of the ratio on the radial distribution described by the parameter $`\alpha `$: The smaller $`\alpha `$ the smaller is the ratio and vice versa. This means that for highly mixed states the optimal collective measurements offer a bigger advantage compared to separate measurements than for states with a small degree of mixing. Or, in other words, the bigger $`\overline{r}`$, Eq. (24), the closer can we come to the optimal limits using separate measurements. It is interesting to note that this statement is true for all $`N`$ and not only in the limit $`N\mathrm{}`$ as shown in .
## VII Conclusion
We have presented estimation methods based on separate adaptive measurements on single qubits. We have demonstrated the measurement schemes for estimating mixed quantum states of qubits. An algorithm is used to update the knowledge about the true quantum state after each measurement and to choose the best measuring operator for the next measurement. With this scheme we have been able to reduce the estimation errors compared to non-adaptive strategies. The best results can be obtained by using schemes related to Kullback information measures. Maximizing this information gain leads to considerable improvements in the estimation quality.
We have also shown that the advantage of collective measurements decreases with decreasing degree of mixing of the initial qubits for all $`N`$, thereby confirming asymptotic results for $`N\mathrm{}`$ found by Gill and Massar.
We have restricted ourselves to simple separate measurements which can be easily realized with nowadays technology. An additional advantage of our scheme is that there is no need to have all $`N`$ quantum systems available at the same time. In contrast to optimal measurement schemes, for which one needs to perform complicated collective measurements on all the systems, our schemes can also be used if the $`N`$ quantum systems can only be prepared one after the other. These features ensure the applicability of our scheme to experiments and practical state estimation problems in quantum information theory.
###### Acknowledgements.
We acknowledge support by the DFG programme โQuanten-Informationsverarbeitungโ, by the European Science Foundation QIT programme and by the IST programme โQUBITSโ of the European Commission.
## A Averaging over the Bloch sphere
In order to quantify the performance of our adaptive methods we have introduced the average fidelity $`F_N`$, Eq. (20). In this Appendix we shortly describe the averaging procedure. In principle the calculation of $`F_N`$ consists of two steps. First we have to determine the average fidelity $`\overline{F}_N(\widehat{\rho })`$ for $`N`$ identical quantum systems prepared in state $`\widehat{\rho }=\widehat{\rho }(R,\mathrm{\Theta },\mathrm{\Phi })`$ by summing over all possible measurement paths $`J`$. Given a specific adaptive method each path is uniquely determined by the initial measurement direction $`(\theta _1,\varphi _1)`$ and by the sequence of measurement results forming a string of 0โs and 1โs. That is, we have $`JJ(\theta _1,\varphi _1;\{0,1\}^N)`$. With the fidelity $`F_N(\widehat{\rho },J)`$ for each path we arrive at
$$\overline{F}_N(\widehat{\rho })\overline{F}_N(R,\mathrm{\Theta },\mathrm{\Phi })=F_N(\widehat{\rho }(R,\mathrm{\Theta },\mathrm{\Phi }),J)_J.$$
(A1)
In principle a simulation of this expression would be straightforward. We choose an initial measuring direction and perform a Monte-Carlo simulation with sufficiently many measurement sequences of length $`N`$. This should be repeated for a dense set of initial measurement directions on the Bloch sphere.
The second step of our averaging procedure consists of averaging $`\overline{F}_N(R,\mathrm{\Theta },\mathrm{\Phi })`$ over all density operators $`\widehat{\rho }=\widehat{\rho }(R,\mathrm{\Theta },\mathrm{\Phi })`$ isotropically distributed over the Bloch sphere. In this way we find the fidelity
$$F_N=\overline{F}_N(R,\mathrm{\Theta },\mathrm{\Phi })_{(R,\mathrm{\Theta },\mathrm{\Phi })}$$
(A2)
which is not biased by any specific choice of density operator and therefore measures the performance of any estimation method.
However, this averaging procedure can be simplified, if we take the following into account. The fidelity $`F_N(\widehat{\rho },J)`$ is rotationally invariant. For a density operator rotated on the Bloch sphere with unitary transformation $`U`$ we have
$$F_N(U\widehat{\rho }U^{},J)=F_N(\widehat{\rho },U^{}JU)$$
(A3)
where $`U^{}JU`$ symbolizes the corresponding rotated path. Hence instead of averaging over all paths $`J`$ in Eq. (A1) we can average over all possible density operators $`\widehat{\rho }=\widehat{\rho }(R,\mathrm{\Theta },\mathrm{\Phi })`$ for a fixed radius $`R`$, that is
$$F_N(\widehat{\rho }(R,\mathrm{\Theta },\mathrm{\Phi }),J)_J=F_N(\widehat{\rho }(R,\mathrm{\Theta },\mathrm{\Phi }),J)_{(\mathrm{\Theta },\mathrm{\Phi })}=\overline{F}_N(R)$$
(A4)
with $`J`$ chosen randomly for each setting $`(\mathrm{\Theta },\mathrm{\Phi })`$. Note that this also reduces $`\overline{F}_N`$ to a pure function of the radius $`R`$. Consequently, the final fidelity reads
$$F_N=\overline{F}_N(R)_R=F_N(\widehat{\rho }(R,\mathrm{\Theta },\mathrm{\Phi }),J)_{(R,\mathrm{\Theta },\mathrm{\Phi })}.$$
(A5)
Therefore, we numerically simulate $`F_N`$ by choosing sufficiently many points $`(R,\mathrm{\Theta },\mathrm{\Phi })`$ isotropically distributed over the Bloch sphere together with a randomly chosen path $`J`$.
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# Semiinfinite cohomology of Lie-* algebras
## 1. Introduction.
This paper is a natural extension of the previous note \[Ar2\]. Semiinfinite cohomology of Tate Lie algebra was defined in that note in terms of some duality resembling Koszul duality. The language of differential graded Lie algebroids was the main technical tool of the note.
The present note is devoted to globalization of the main construction from \[Ar2\] in the following sense. The setup in \[Ar2\] included a suitably chosen module over a Tate Lie algebra $`๐ค`$ with a fixed Lie subalgebra $`๐`$ being a c-lattice in $`๐ค`$.
The rough global analogue of this picture is as follows. Consider a compact curve $`X`$ over a field of characteristic zero. Denote by $`\mathrm{mod}D_X`$ the category of (right) D-modules on $`X`$. We fix a Lie algebra $`๐ข`$ in the category $`\mathrm{mod}D_X`$. This data can be viewed as a family of Lie algebras $`๐ค_x`$ along the curve $`X`$. Another part of the data includes a Lie subalgebra $`๐ข`$. So the problem is to define semiinfinite cohomology complex of such pair.
In fact we need some additional constraints on the pair $`๐ข`$. So the formal picture starts from a different notion of a Lie-\* algebra $``$ on $`X`$ (see the precise definition in Section 2). Roughly spaking a $`D_X`$-locally free Lie-\* algebra is a D-module incarnation of a Lie algebra in the category of vector bundles on $`X`$ with the bracket given by a differential operator. We define two types of modules over a Lie-\* algebra (see 2.2.1). The first one called a Lie-\* module is just a D-module incarnation of the module over the Lie algebra in the category of vector bundles, like above, with the action given by a differential operator. Still we will be more interested in the second type of modules over a Lie-\* algebra called chiral modules (see 2.2.1 for the definition).
So starting from a Lie-\* algebra $``$ and a chiral module $``$ we perform the main construction more or less parallel to the one from \[Ar2\]. Namely we define the Lie algebra $`๐ข=๐ข()`$ in the category $`\mathrm{mod}D_X`$ with the Lie subalgebra $`=()๐ข()`$. We show that a $``$-chiral module $``$ becomes a $`๐ข()`$-module.
Next, imitating the construction of \[Ar2\] Section 4, we define a DG Lie algebroid $`๐^{}()`$ in the category of $`D_X`$-modules over a DG $`\stackrel{!}{}`$-algebra $`^{}=^{}()`$. Koszul duality type construction provides a left DG-module $`C^{}(,)`$ over $`๐^{}()`$.
To go further one needs to pass to a central extension of $``$ called the Tate central extension and denoted by $`_{_{\mathrm{๐ณ๐บ๐๐พ}}}`$. Let $``$ be a chiral module over $`_{_{\mathrm{๐ณ๐บ๐๐พ}}}`$. It turns out that the complex of D-modules $`C^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}},)`$ by some antipode construction becomes a right module over $`๐^{}()`$.
Finally we consider the homological Chevalley complex of the DG Lie algebroid $`๐^{}()`$ in the category of $`D_X`$-modules with coefficients in $`C^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}},)`$ (see 4.2 for the definition of the homological Chevalley complex of a DG Lie algebroid). We call the obtained complex of D-modules the global standard semiinfinite complex of the Lie-\* algebra $``$ with coefficients in the chiral module $``$.
Let us say a few words about the structure of the paper. In Section 2 we collect necessary definitions and simple facts about Lie-\* algebras, Lie-\* modules, chiral modules etc. Section 3 is devoted to the construction of the Tate central extension of a Lie-\* algebra. Section 4 contains all the necessary constructions concerning DG Lie algebroids in the category of $`D_X`$-modules. In particular we present the definition of the homological Chevalley complex of of a DG Lie algebroid with coefficients in a right DG-module. Section 5 is the heart of the paper. We present the constructions of the Lie algebras $`()`$ and $`๐ข()`$ in the category of $`D_X`$-modules. Then we define the DG Lie algebroid $`๐^{}()`$ over the DG $`\stackrel{!}{}`$-algebra $`^{}()`$. Finally after overcoming the problem of necessity to pass to the Tate central extension of $``$ we present the standard semiinfinite complex $`C^\frac{\mathrm{}}{2}(_{_{\mathrm{๐ณ๐บ๐๐พ}}},)`$ for a $`_{_{\mathrm{๐ณ๐บ๐๐พ}}}`$-chiral module $``$.
Note that the paper \[BD\] contains a construction of the global BRST complex for calculating the semiinfinite cohomology of a chiral module over a Lie-\* algebra. Somehow the present paper grew out of an attempt to understand that construction avoiding the notions of chiral algebras, chiral enveloping algebras etc.The technique used in the definition in \[BD\] is quite different from ours and it is not checked that the two constructions give the same answer.
Acknowledgements. The author is happy to thank Sasha Beilinson and Dennis Gaitsgory who explained him the chiral algebra basics. The author also would like to express his deep gratitude to IAS, Princeton, USA, and IHES, Bures-sur-Yvette, France, where parts of the work on the paper were done, for hospitality and extremely stimulating working conditions.
## 2. Lie-\* algebras.
In this section we recall briefly basic notation and constructions concerning Lie-\* algebras. In our exposition we follow \[Ga\]. We will be working over a fixed smooth curve $`X`$. Denote the diagonal embedding $`XX\times X`$ by $`\mathrm{\Delta }`$. The embedding of the complementary open set $`X\times XX_\mathrm{\Delta }X\times X`$ is denoted by $`j`$.
Let $`D_X\mathrm{-}\mathrm{mod}`$ (resp. $`\mathrm{mod}\mathrm{-}D_X`$) be the category of left (resp. right) modules $``$ over the sheaf of algebras of differential operators on $`X`$ such that $``$ is quasicoherent over $`๐ช_X`$. It is known that the category $`D_X\mathrm{-}\mathrm{mod}`$ is naturally a symmetric tensor category with the tensor product given by $`_{๐ช_X}๐ฉ`$ for $`,๐ฉD_X\mathrm{-}\mathrm{mod}`$. Let $`\mathrm{\Omega }=\mathrm{\Omega }_X^1`$. Then the category $`\mathrm{mod}\mathrm{-}D_X`$ becomes a symmetric tensor category with the tensor product given by
$$\stackrel{!}{}๐ฉ=(_{๐ช_X}\mathrm{\Omega }^1)_{๐ช_X}(๐ฉ_{๐ช_X}\mathrm{\Omega }^1)_{๐ช_X}\mathrm{\Omega }.$$
### 2.1. Definition.
Recall that a Lie-\* algebra on $`X`$ is a right D-module $``$ with the map
$$\{\}:\mathrm{\Delta }_!()$$
which is antisymmetric and satisfies the Jacobi identity in the following sense. If $`abcf(x,y,z)`$ is a section of the $``$ on $`X\times X\times X`$, then the element
$$\{\{f(x,y,z)ab\}c\}+\sigma _{1,2,3}(\{\{f(z,x,y)bc\}a\})+\sigma _{1,2,3}(\{\{f(y,z,x)ca\}b\})$$
of $`\mathrm{\Delta }_{x=y=z}{}_{!}{}^{}()`$ vanishes. Here $`\sigma _{1,2,3}`$ denotes the lift of the cyclic automorphism of $`X\times X\times X`$: $`(x,y,z)(y,z,x)`$ to the D-module $`\mathrm{\Delta }_{x=y=z}{}_{!}{}^{}()`$.
#### 2.1.1.
Note that if $``$ is a Lie-\* algebra, it follows from the definition that $`DR^0()`$ is a sheaf of ordinary Lie algebras; moreover it acts on $``$ by endomorphisms of the D-module structure that are derivations of the Lie-\* structure.
In particular, for an affine subset $`UX`$, $`DR^0(U,)`$ is a Lie algebra. Thus for any point $`xX`$ the topological spaces $`DR^0(\mathrm{Spec}(\widehat{๐ช}_x),)`$ and $`DR^0(\mathrm{Spec}(\widehat{๐ฆ}_x),)`$ carry the natural structures of topological Lie algebras.
#### 2.1.2.
Our next step consists of reformulation of the Lie algebra definition in terms of Lie coalgebras in the standard tensor structure $`\stackrel{!}{}`$ on right D-modules. Lemma: Let $`_1`$, $`_2`$ be two D-modules on $`X`$ with $`_1`$ being locally free and finitely generated. Then:
* For a third D-module $``$ on $`X`$, there is a canonical isomorphism:
$$\mathrm{Hom}_{D_X^2}(_1,\mathrm{\Delta }_!(_2))\stackrel{~}{}\mathrm{Hom}_{D_X}(,_1^{}\stackrel{!}{}_2).$$
* The canonical map (from point (i)) $`\left(_1^{}\stackrel{!}{}_2\right)_1\mathrm{\Delta }_!(_2)`$ induces an isomorphism $`DR^0(_1^{}\stackrel{!}{}_2)\mathrm{Hom}_D(_1,_2)`$. โ
Below we always suppose that any Lie-\* algebra $``$ we work with is locally free and finitely generated as a D-module on $`X`$. Corollary: For a Lie-\* algebra $``$ we have canonical maps
$$\mathrm{co}\mathrm{ad}:^{}\mathrm{\Delta }_!(^{})\text{ and }\mathrm{co}\mathrm{br}:^{}^{}\stackrel{!}{}^{}.\mathit{}$$
We shall call the two maps of the above corollary โthe co-adjoint actionโ and โthe co-bracketโ, respectively.
#### 2.1.3.
Remark: In particular note that to specify a structure of a Lie-\* algebra on a D-module $``$ is the same as to provide a structure of a Lie coalgebra (in the usual tensor structure on the category of right D-modules) on the D-module $`^{}:=\underset{ยฏ}{\mathrm{Hom}}_{D_X}(,D_X\mathrm{\Omega }_X)`$. The co-bracket here is the one obtained in the previous Corollary.
### 2.2. Modules over a Lie-\* algebra.
There are two different ways to define a notion of a module over a Lie-\* algebra.
#### 2.2.1.
A Lie-\* module over a Lie-\* algebra $``$ is a (right) D-module $``$ on $`X`$ with a map
$$\rho :\mathrm{\Delta }_!()$$
such that for a section $`abm`$ of $``$ the two sections
$$\rho (\{ab\},m)\text{ and }\rho (a,\rho (b,m))\sigma _{1,2}(\rho (b,\rho (a,m)))$$
of $`\mathrm{\Delta }_{1,2,3}{}_{!}{}^{}()`$ coincide.
Like in the Lie-\* algebra case a structure of the Lie-\* module on $``$ provides a structure of the sheaf of modules over the sheaf of Lie algebras $`DR^0()`$ on $`DR^0()`$. Moreover the sheaf of Lie algebras $`DR^0()`$ acts on $``$ and contrarywize such an action recovers a Lie-\* module structure if it is given by differential operators (see \[BD\], 2.5.4). Remark: Similarly to 2.1.3 one can easily check that to specify a structure of a Lie-\* module over the Lie-\* algebra $``$ on a D-module $``$ is the same as to make $``$ a comodule over the Lie coalgebra $`^{}`$, i.e. to provide a D-module map $`\mathrm{co}\mathrm{ac}:^{}\stackrel{!}{}`$ satisfying certain coassociativity constraints.
#### 2.2.2.
A chiral module over a Lie-\* algebra is again a (right) D-module $``$ on $`X`$, but with an operation
$$\rho :j_{}j^{}()\mathrm{\Delta }_!()$$
such that every section $`f(x,y,z)abm\mathrm{\Gamma }(X\times X\times X(\mathrm{\Delta }_{x=z}\mathrm{\Delta }_{y=z}),)`$ satisfies an identity as follows:
$$\rho (\{f(x,y,z)ab\},m)=\rho (a,\rho (f(x,y,z)b,m))\sigma _{1,2}(\rho (b,\rho (f(y,x,z)a,m)),$$
as sections of $`\mathrm{\Delta }_{x=y=z}{}_{!}{}^{}()`$.
### 2.3. Naive de Rham complex for a Lie-\* module over a Lie-\* algebra.
Let $``$ be a Lie-\* module over a (locally D-free finitely generated) Lie-\* algebra $``$. Consider a complex $`C^{}(,)`$ of D-modules on X as follows: $`C^k(,)=\mathrm{\Lambda }^k(^{})\stackrel{!}{}`$ as a D-module. The differential is given by
$$d(\omega _1\mathrm{}\omega _km)=\omega _1\mathrm{}\mathrm{co}\mathrm{br}(\omega _i)\mathrm{}\omega _km+\omega _1\mathrm{}\omega _k\mathrm{co}\mathrm{ac}(m).$$
Here $`\omega _i`$ are the sections of $`^{}`$ and $`m`$ is the section of $``$. Lemma: The complex of D-modules $`C^{}(,)`$ is well-defined, i.e. $`d^2=0`$. โRemark: We call the constructed complex the naive de Rham complex for the Lie-\* module over the Lie-\* algebra. However below we construct a more sophisticated complex for a chiral module over a Lie-\* algebra. That complex is also of de Rham origin. It will be realized as de Rham complex of a certain DG Lie-\* module over a certain DG Lie-\* algebroid.
### 2.4. De Rham DG-algebra of a Lie-\* algebra.
Consider the complex $`C^{}(,)`$ for $``$ equal to the trivial Lie-\* module $`\mathrm{\Omega }`$ over $``$. Lemma: The wedge product makes the complex $`C^{}(,\mathrm{\Omega })`$ into a supercommutative DG-algebra in the tensor category of right D-modules. โ
We denote this DG-algebra by $`^{}=^{}()`$.
## 3. Tate extension of a Lie-\* algebra.
### 3.1. Matrix Lie-\* algebra.
Let $`๐ฑ`$ be a locally free finitely generated D-module on $`X`$, and let $`๐ฑ^{}:=\underset{ยฏ}{\mathrm{Hom}}_{D_X}(๐ฑ,D_X\mathrm{\Omega }_X)`$ be its (Verdier) dual D-module. Consider the $`\stackrel{!}{}`$-tensor product of $`๐ฑ`$ and $`๐ฑ^{}`$.
#### 3.1.1.
Remark: By Lemma 2.1.2(i) we have a Lie-\* pairing
$$:๐ฑ๐ฑ^{}\mathrm{\Delta }_!(\mathrm{\Omega }).$$
Lemma: The D-module $`๐ฑ\stackrel{!}{}๐ฑ^{}=:at`$ carries a natural structure of an associative-\* algebra on $`X`$.
Proof. The associative product map is defined as follows:
$$\begin{array}{c}\mathrm{as}:\left(๐ฑ_{(1)}\stackrel{!}{}๐ฑ_{(1)}^{}\right)\left(๐ฑ_{(2)}\stackrel{!}{}๐ฑ_{(2)}^{}\right)\hfill \\ \hfill \stackrel{~}{}\left((๐ฑ_{(1)}\mathrm{\Omega })\stackrel{!}{}(\mathrm{\Omega }๐ฑ_{(2)}^{})\right)\stackrel{!}{}\left((๐ฑ_{(2)}\mathrm{\Omega })\stackrel{!}{}(\mathrm{\Omega }๐ฑ_{(1)}^{})\right)\\ \hfill \stackrel{\stackrel{!}{}\mathrm{id}}{}\mathrm{\Delta }_!(\mathrm{\Omega })\stackrel{!}{}\left((๐ฑ_{(2)}\mathrm{\Omega })\stackrel{!}{}(\mathrm{\Omega }๐ฑ_{(1)}^{})\right)\stackrel{~}{}\mathrm{\Delta }_!(๐ฑ\stackrel{!}{}๐ฑ^{}).\end{array}$$
Here the indexes $`(1)`$ and $`(2)`$ denote the factors in the product, and the \*-pairing is taken between the ones in the first pair of brackets. It is left to the reader to check the associativity of the product.โ
Now we obtain the Lie-\* bracket on $`at`$ from the associative-\* product in the usual way:
$$ab\mathrm{as}(ab)\mathrm{as}(ba).$$
We call $`at`$ the matrix Lie-\* algebra of the D-module $`๐ฑ`$.
Note that the associative-\* algebra $`at`$ acts both on $`๐ฑ`$ and on $`๐ฑ^{}`$ in a canonical way and the pairing $``$ is $`at`$-invariant.
### 3.2. Tate extension of the Lie-\* algebra $`at`$.
The material of this subsection is almost word to word copied from \[BD\], 2.6. We include it in our paper just for the sake of completeness.
#### 3.2.1.
The Tate extension is a canonical central extension of Lie-\* algebras
$$0\mathrm{\Omega }_Xat^{\mathrm{}}\stackrel{\pi ^{\mathrm{}}}{}at0.$$
To define $`at^{\mathrm{}}`$ as a D-module consider the exact sequence of D-modules on $`X\times X`$
$$๐ฑ๐ฑ^{}\stackrel{๐}{}j_{}j^{}(๐ฑ๐ฑ^{})\stackrel{๐}{}\mathrm{\Delta }_!(๐ฑ\stackrel{!}{}๐ฑ^{})0.$$
Here as before $`\mathrm{\Delta }:XX\times X`$ is the diagonal embedding, $`j:U:=X\times X\mathrm{\Delta }(X)X\times X`$ is the complementary embedding, and $`\pi `$ is the canonical arrow.
Namely one has $`๐ฑ\stackrel{!}{}๐ฑ^{}=H^1\mathrm{\Delta }^!(๐ฑ๐ฑ^{})=\mathrm{Coker}\epsilon `$; explicitly, $`\pi `$ sends $`(t_2t_1)^1vv^{}j_{}j^{}๐ฑ๐ฑ^{}`$ to $`vv^{}(dt)^1\mathrm{\Delta }_{}(๐ฑ\stackrel{!}{}๐ฑ^{})\mathrm{\Delta }_!(๐ฑ\stackrel{!}{}๐ฑ^{})`$. Note that $``$ vanishes on $`\mathrm{Ker}\epsilon `$ and, pushing forward the above exact sequence by $``$, we get an extension of $`\mathrm{\Delta }_!(๐ฑ\stackrel{!}{}๐ฑ^{})`$ by $`\mathrm{\Delta }_!(\mathrm{\Omega })`$.
This extension is supported on the diagonal. Applying $`\mathrm{\Delta }^!`$ we get the Tate extension $`at^{\mathrm{}}`$. We denote the canonical morphism $`j_{}j^{}(๐ฑ๐ฑ^{})\mathrm{\Delta }_!at^{\mathrm{}}`$ as $`\mu _{Mat}`$.
#### 3.2.2.
The above construction is natural with respect to Lie-\* algebras actions. Namely, assume that a Lie-\* algebra $``$ acts on $`๐ฑ`$, $`๐ฑ^{}`$ so that $``$ is $``$-invariant. Then the $``$-action on $`at=๐ฑ\stackrel{!}{}๐ฑ^{}`$ lifts canonically to an $``$-action on $`at^{\mathrm{}}`$. To see this, consider the $`D_X^2`$-modules $`\mathrm{\Delta }^{}(๐ฑ๐ฑ^{})`$, $`\mathrm{\Delta }^{}j_{}j^{}(๐ฑ๐ฑ^{})`$. The Lie algebra $`h()`$ acts on them in the obvious manner. The morphism $`:\mathrm{\Delta }^{}(๐ฑ๐ฑ^{})\mathrm{\Delta }^{}\mathrm{\Delta }_!(\mathrm{\Omega })`$ is $`DR^0()`$-invariant. Therefore $`DR^0()`$ acts on $`at^{\mathrm{}}`$. This action is uniquely determined by property that $`\mu _{at}`$ is a morphism of $`DR^0(at)`$-modules. Evidently the action of $`DR^0()`$ is given by differential operators so we have the desired $``$-action on $`at^{\mathrm{}}`$.
In particular, the canonical $`at`$-actions on $`๐ฑ`$, $`๐ฑ^{}`$ define a $`at`$-action on $`at^{\mathrm{}}`$ that lifts the adjoint action on $`at`$. Composing this action with $`\pi ^{\mathrm{}}`$ we get an operation
$$\{\}:at^{\mathrm{}}at^{\mathrm{}}\mathrm{\Delta }_!(at^{\mathrm{}}).$$
#### 3.2.3.
Lemma: This is a Lie-\* bracket. โ
### 3.3. Tate extension of a Lie-\* algebra.
Now let $``$ be an arbitrary $`D_X`$-locally free finitely generated Lie-\* algebra. Recall that by Lemma 2.1.2 we have a canonical co-action map of the $`D_X`$-modules
$$\mathrm{co}\mathrm{ac}:\stackrel{!}{}^{}.$$
We interprete as a (D-module) map $`\mathrm{can}:at`$. Lemma: The map $`\mathrm{can}`$ is a morphism of Lie-\* algebras. โ
#### 3.3.1.
Corollary: For any $`D_X`$-locally free finitely generated Lie-\* algebra $``$ there exists a central extension in the class of Lie-\* algebras as follows
$$0\mathrm{\Omega }_X^{\mathrm{}}0.$$
This is just the inverse image of the Tate central extension of $`at`$. โ
Below we denote the Lie-\* algebra $`^{\mathrm{}}`$ by $`_{\mathrm{๐ณ๐บ๐๐พ}}`$ and cal it the Tate central extension of $``$. Remark: Note that for the complete curve $`X`$ the short exact sequence of the $`D_X`$-modules that defines the extension $`_{\mathrm{๐ณ๐บ๐๐พ}}`$ does not split, even if we forget about the Lie-\* algebra structures.
### 3.4. Local analog of the Tate extension.
Fix a point $`xX`$. Let $`\widehat{๐ช}_x`$ (resp. $`\widehat{๐ฆ}_x`$) be the completion of the local ring of the point $`x`$ (resp. of the local field of the point $`x`$).
Consider the topological Lie algebra $`DR^0(\mathrm{Spec}(\widehat{๐ฆ}_x),)`$ with the Lie subalgebra $`DR^0(\mathrm{Spec}(\widehat{๐ช}_x),)`$.
#### 3.4.1.
Lemma:
* $`DR^0(\mathrm{Spec}(\widehat{๐ฆ}_x),)`$ is a Tate Lie algebra.
* The subalgebra $`DR^0(\mathrm{Spec}(\widehat{๐ช}_x),)`$ is a c-lattice in $`DR^0(\mathrm{Spec}(\widehat{๐ฆ}_x),)`$. โ
In particular we have a one dimensional Lie algebra central extension
$$0DR^0(\mathrm{Spec}(\widehat{๐ฆ}_x),_{\mathrm{๐ณ๐บ๐๐พ}})DR^0(\mathrm{Spec}(\widehat{๐ฆ}_x),)0.$$
#### 3.4.2.
Proposition: The central extension $`DR^0(\mathrm{Spec}(\widehat{๐ฆ}_x),_{\mathrm{๐ณ๐บ๐๐พ}})`$ coincides with the extension of the Tate Lie algebra $`DR^0(\mathrm{Spec}(\widehat{๐ฆ}_x),)`$ with the help of the critical cocycle (see e.g. \[Ar2\], 4.3.3). โ
## 4. DG Lie algebroids in the category of $`D_X`$-modules.
Recall that $`\mathrm{mod}\mathrm{-}D_X`$ is a symmetric tensor category. Below we use this structure to mimick the ordinary definition of a Lie algebroid in the category of vector spaces.
### 4.1.
Definition: Let $``$ be a $`\stackrel{!}{}`$-commutative algebra in $`\mathrm{mod}\mathrm{-}D_X`$, and let $`๐`$ be a $``$-module (so we have a D-module map $`\stackrel{!}{}๐๐`$ providing the structure) carrying a Lie algebra structure in $`\mathrm{mod}\mathrm{-}D_X`$ (i.e. a Lie bracket map $`[,]:\mathrm{\Lambda }^2(๐)๐`$ satisfying the Jacobi identity is given). Moreover suppose that $`๐`$ acts on $``$ by derivations (i.e. we have a D-module map $`๐\stackrel{!}{}`$ such that $`[a,rb]=a(r)b+r[a,b]`$ is satisfied for any sections $`a,b\mathrm{\Gamma }(X,๐)`$ and $`r\mathrm{\Gamma }(X,)`$).
We call the above data the Lie algebroid in the category $`\mathrm{mod}\mathrm{-}D_X`$ over the $`\stackrel{!}{}`$-commutative algebra $``$.
From now on we assume that all the appearing $``$-algebroids are locally free as $``$-modules.
#### 4.1.1.
By definition a right module (resp. a left module) over a Lie algebroid $`(๐,)`$ on $`X`$ is a sheaf of $``$-modules $``$ with the Lie action of $`๐`$ satisfying the constraint $`(rm)a=r(ma)(a(r))m`$ (resp. $`a(rm)=r(am)+(a(r))m`$) for any sections $`a`$ of $`๐`$, $`r`$ of $``$ and $`m`$ of $``$.
#### 4.1.2.
Recall that the universal enveloping algebra for a $``$-Lie algebroid $``$ in $`\mathrm{mod}\mathrm{-}D_X`$ is defined in the same way it is done for Lie algebroids over vector spaces: we take the free algebra in the category $`\mathrm{mod}\mathrm{-}D_X`$ generated by $``$ and take its quotient by the obvious ideal of relations including the one expressing the action of $``$ on $``$ by derivations. We denote the obtained associative algebra in $`\mathrm{mod}\mathrm{-}D_X`$ by $`๐ฐ_{}()`$.
### 4.2. Homological Chevalley complex for a right module over a Lie algebroid in $`\mathrm{mod}\mathrm{-}D_X`$.
For a right $`๐`$-module $``$ consider the graded $`D_X`$-module on $`C^{}(๐,)`$ as follows:
$$C^{}(๐,)=\underset{๐}{}C^k(๐,),C^k(๐,)=_{}(\mathrm{\Lambda }_{}^k(๐)).$$
We endow the graded vector space with the differential as follows: For sections $`a_1,\mathrm{},a_p`$ of $`๐`$ and $`m`$ of $``$ we put
$$d(ma_1\mathrm{}a_p)=\underset{i}{}(1)^ima_ia_1\mathrm{}\widehat{a}_i\mathrm{}a_p$$
$$+\underset{i<j}{}(1)^{i+j}m[a_i,a_j]a_1\mathrm{}\widehat{a}_i\mathrm{}\widehat{a}_j\mathrm{}a_p.$$
Lemma:
* The differential in the complex is well defined.
* The differential satisfies $`d^2=0`$.
Proof. (i) Let us perform a calculation showing that the differential $`d:C^kC^{k+1}`$ is well defined for $`k=2`$, the general case is quite similar. We have
$$d(mra_1a_2)=m(ra_1)a_2ma_2ra_1+m[a_2,ra_1]$$
$$=m(a_2(r))a_1+mr[a_2,a_1]r(ma_2)a_1+m(ra_1)a_2$$
$$=(a_2(r))ma_1+rm[a_2,a_1](rm)a_2a_1a_2(r)ma_1+(rm)a_1a_2$$
$$=(rm)a_1a_2(rm)a_2a_1+rm[a_2,a_1]=d(rma_1a_2).$$
The general case is quite similar.
(ii) This is the usual calculation in the Chevalley complex. โRemark: We have shown above that there exists a homological Chevalley complex for a right $`๐`$-module $``$. A very similar calculation proves the existence of the cohomological Chevalley complex for a left $`๐`$-module $`๐ฉ`$, of the size $`\underset{ยฏ}{\mathrm{Hom}}_{}(\mathrm{\Lambda }_{}^{}(๐),๐ฉ)`$.
#### 4.2.1.
In fact both the homological and the cohomological versions of the Chevalley complex for a Lie algebroid appear naturally โin a coordinate-free wayโ as a result of the following construction.
Consider the tautological left $`๐`$-module $`\underset{ยฏ}{}`$. We construct its standard projective resolution in the way it is usually done for modules over Lie algebras.
Namely consider the complex
$$\mathrm{๐ฒ๐๐บ๐}^{}(๐,\underset{ยฏ}{}):=๐ฐ_{}(๐)_{}\mathrm{\Lambda }_{}^{}(๐)_{}\underset{ยฏ}{}$$
with the standard differential. Note that the differential uses the right $`CA`$-module structure on $`๐ฐ_{}(๐)`$. Lemma:
* The complex of $`D_X`$-modules $`\mathrm{๐ฒ๐๐บ๐}^{}(๐,\underset{ยฏ}{})`$ is well-defined.
* $`\mathrm{๐ฒ๐๐บ๐}^{}(๐,\underset{ยฏ}{})`$ is a complex of $`๐`$-modules.
* The homological Chevalley complex $`C^{}(๐,)`$ for a right $`๐`$-module $``$ is isomorphic to $`_๐\mathrm{๐ฒ๐๐บ๐}^{}(๐,\underset{ยฏ}{}).`$
* The cohomological Chevalley complex $`\underset{ยฏ}{\mathrm{Hom}}_{}(\mathrm{\Lambda }_{}^{}(๐),๐ฉ)`$ for a left $`๐`$-module $`๐ฉ`$ is isomorphic to $`\underset{ยฏ}{\mathrm{Hom}}_๐(\mathrm{๐ฒ๐๐บ๐}^{}(๐,\underset{ยฏ}{}),๐ฉ).\mathit{}`$
### 4.3. DG Lie algebroids in $`D_X`$-modules.
Now let $`^{}=_k^k`$ be a graded supercommutative $`\stackrel{!}{}`$-algebra on $`X`$ and let $`๐^{}=๐^k`$ be a graded super Lie algebroid over $``$. Suppose also there are an odd derivation $`d_{}`$ on $`^{}`$ of degree $`1`$ and an odd derivation $`d_๐`$ of the Lie superalgebra $`๐^{}`$ also of degree $`1`$ satisfying Leibnitz rule with respect one to another. Moreover both of them satisfy $`d^2=0`$.
#### 4.3.1.
Definition: The data $`(๐^{},^{},d_๐,d_{})`$ are called the differential graded Lie algebroid in the category $`\mathrm{mod}\mathrm{-}D_X`$ or, for short, a DG Lie algebroid on $`X`$.
The notion of a left (resp. right) DG-module over a DG algebroid in $`\mathrm{mod}\mathrm{-}D_X`$ is a natural combination of the previous definitions and we do not spell it out explicitly. The category of left (resp. right) DG-modules over a DG Lie algebroid $`๐=๐_X`$ is denoted by $`DG\mathrm{-}๐^{}\mathrm{-}\mathrm{mod}`$ (resp. $`DG\mathrm{-}\mathrm{mod}\mathrm{-}๐^{}`$).
### 4.4. Homological Chevalley complex for a DG Lie algebroid in $`\mathrm{mod}\mathrm{-}D_X`$.
Now we sort of add a second differential on the homological Chevalley complex given in 4.2. For $`^{}DG\mathrm{-}\mathrm{mod}\mathrm{-}๐^{}`$ consider the bigraded vector space $`C^{}(๐^{},M^{})`$ as follows: $`C^{}(๐,M^{})=^{}_{^{}}(\mathrm{\Lambda }_{^{}}^{}(๐^{})),`$ here the first grading comes from the number of wedges in the exterior product and
$$C^k(๐^{},^{})=\left(^{}_{^{}}(\mathrm{\Lambda }_{^{}}^{}(๐^{}))\right)^k$$
in the graded tensor product sense.
Consider the two differentials on the bigraded vector space. The first one of the grading $`(1,0)`$ is the usual Chevalley differential like in 4.2. The second differential of the grading $`(0,1)`$ is provided by the differentials on $`^{}`$ and $`\mathrm{\Lambda }_{^{}}^{}(๐^{})`$.
Consider the total grading on the bigraded space and the total differential on it.
#### 4.4.1.
Lemma: The differential $`d_1+d_2`$ is well defined and its square equals zero. โ
## 5. Semiinfinite cohomology via DG Lie algebroids in $`\mathrm{mod}\mathrm{-}D_X`$.
In this section we show that the standard complex for the computation of semiinfinite cohomology of a chiral module over a Lie-\* algebra coincides with the homological Chevalley complex of the form 4.4 for a certain DG Lie algebroid in the category $`\mathrm{mod}\mathrm{-}D_X`$ and a certain right module over it.
### 5.1. Construction of the Lie algebroid in $`\mathrm{mod}\mathrm{-}D_X`$.
Fix a Lie-\* algebra $``$ on $`X`$. As before, we suppose that it is locally free and finitely generated over $`D_X`$.
The construction will be local and we can assume that $`X`$ is affine.
Consider the completion of the $`D_{X\times X}`$-module $`\mathrm{\Omega }`$ along the diagonal. We denote this D-module by $`_{\widehat{\mathrm{\Delta }}}`$. One can view the D-module $`_{\widehat{\mathrm{\Delta }}}`$ as the restriction of $`\mathrm{\Omega }`$ to the โfamily of formal discsโ parametrized by the diagonal.
Consider also the D-module $`j_{}(๐ช_U)_{๐ช_{X\times X}}_{\widehat{\mathrm{\Delta }}}`$ denoted by $`_{\stackrel{~}{\mathrm{\Delta }}}`$. One can view the D-module $`_{\stackrel{~}{\mathrm{\Delta }}}`$ as the restriction of $`\mathrm{\Omega }`$ to the โfamily of punctured formal discsโ parametrized by the diagonal.
#### 5.1.1.
Lemma: We have the short exact sequence of the d-modules
$$0_{\widehat{\mathrm{\Delta }}}_{\stackrel{~}{\mathrm{\Delta }}}\mathrm{\Delta }_!()0\mathit{}$$
Now we take the D-module direct images $`:=p_1{}_{}{}^{}\left(_{\widehat{\mathrm{\Delta }}}\right)`$ and $`๐ข:=p_1{}_{}{}^{}\left(_{\stackrel{~}{\mathrm{\Delta }}}\right)`$ on $`X`$ (here $`p_1`$ is the projection $`(x,y)X\times XxX`$). There is an obvious map $`๐ข`$ and from the fact that $``$ is a Lie-\* algebra we infer that both $``$ and $`๐ข`$ are Lie algebras in the category of right D-modules on $`X`$.
Note that the stalk of the D-module $``$ (resp. of $`๐ข`$) at a point $`xX`$ equals $`DR^0(\widehat{๐ช}_x,)`$ (resp. $`DR^0(\widehat{๐ฆ}_x,)`$), where $`\widehat{๐ช}_x`$ (resp. $`\widehat{๐ฆ}_x`$) denotes the spectrum of the completed local ring (resp. the completed local field) at $`x`$. We abuse some notation here. Lemma: There exists a short exact sequence of D-modules on $`X`$ as follows:
$$0๐ข0.$$
Proof. Consider the short exact sequence of D-modules from the previous Lemma Now take the (D-module) direct image of the exact sequence under $`p_1`$. It is left to the reader that the sequence obtained as a result coincides with the one we need. โ
We choose the basic $`\stackrel{!}{}`$-supercommutative algebra for our algebroid to be $`^{}()=\mathrm{\Lambda }^{}(^{})`$.
#### 5.1.2.
Lemma:
* The Lie algebra in D-modules $``$ acts on $``$.
* The Lie algebra in D-modules $``$ acts on the $`\stackrel{!}{}`$-supercommutative algebra $`^{}()`$ by derivations.
Proof. Note that the second assertion of the Lemma follows from the first one since we can extend the action from the generators of $`^{}()`$ to the whole algebra by the Leibnitz rule.
Now we construct the action map for (i) explicitly. We have the following sequence of morphisms of D-modules
$$p_1{}_{}{}^{}\left(_{\widehat{\mathrm{\Delta }}}\right)\stackrel{!}{}^{}p_1{}_{}{}^{}p_{1}^{}\left(p_1{}_{}{}^{}\left(_{\widehat{\mathrm{\Delta }}}\right)\stackrel{!}{}^{}\right)p_1{}_{}{}^{}\left(\left(_{\widehat{\mathrm{\Delta }}}\right)\stackrel{!}{}p_1^{}(^{})\right)$$
$$\stackrel{~}{}p_1{}_{}{}^{}((^{})_{\widehat{\mathrm{\Delta }}})\stackrel{๐}{}p_1{}_{}{}^{}(\mathrm{\Delta }_!(^{}))\stackrel{~}{}^{}.$$
Here $`(^{})_{\widehat{\mathrm{\Delta }}}`$ denotes the completion of the D-module $`^{}`$ along the diagonal in $`X\times X`$.
Note that the D-module $`\mathrm{\Delta }_!(^{})`$ is locally finite over the ideal of the diagonal. Thus the completion of the map $`^{}\mathrm{\Delta }_!(^{})`$ is well defined.
It is left to the reader that the composition of the above maps provides the action of the Lie algebra $``$ on $`^{}`$. โCorollary: The graded D-module $`๐^{}():=\stackrel{!}{}^{}()`$ carries a natural structure of a Lie algebroid over the $`\stackrel{!}{}`$-supercommutative algebra $`^{}()`$. โ
### 5.2. Construction of the differential on $`๐^{}()`$.
Note that the differential $`d_{}`$ on the $`\stackrel{!}{}`$-supercommutative algebra $`^{}()`$ is already constructed. Moreover it remains to construct the component of the differential on $`\stackrel{!}{}^{}()=\stackrel{!}{}\mathrm{\Lambda }^{}(^{})`$ as follows: $`d_{}:\stackrel{!}{}^{}`$. After that the differential on the whole Lie algebroid is obtained from $`d_{}+d_{}`$ by the Leibnitz rule.
Now by Lemma 2.1.2(i) we rewrite the map in question as $`\mathrm{\Delta }_!()`$ or
$$\left(p_1{}_{}{}^{}\left(_{\widehat{\mathrm{\Delta }}}\right)\right)\mathrm{\Delta }_!\left(p_1{}_{}{}^{}\left(_{\widehat{\mathrm{\Delta }}}\right)\right).$$
Note that $`DR^0()`$ acts naturally on every stalk of $`p_1{}_{}{}^{}\left(_{\widehat{\mathrm{\Delta }}}\right)`$ at any point $`xX`$. Recall that the stalk equals $`DR^0(\widehat{๐ช}_x,)`$. Thus the sheaf of Lie algebras acts on it acts on $`\left(p_1{}_{}{}^{}\left(_{\widehat{\mathrm{\Delta }}}\right)\right)`$.
#### 5.2.1.
Proposition: The above action is given by differential operators, i.e. it lifts to the required morphism $`\mathrm{\Delta }_!()`$.
Proof. This is a local assertion. It is left to the reader to check it. โCorollary: The D-module $`๐^{}():=\stackrel{!}{}^{}()`$ carries a natural structure of a DG Lie algebroid over the $`\stackrel{!}{}`$-supercommutative DG algebra $`^{}()`$.
#### 5.2.2.
Remark: In particular the D-module $`๐^{}(_{\mathrm{๐ณ๐บ๐๐พ}}):=p_1{}_{}{}^{}((_{_{\mathrm{๐ณ๐บ๐๐พ}}})_{\widehat{\mathrm{\Delta }}})\stackrel{!}{}\mathrm{\Lambda }^{}(^{})`$ is a DG Lie algebroid over $`^{}()`$. This follows from the obvious fact that $`\underset{ยฏ}{\mathrm{Hom}}_{D_X}(,D_X\mathrm{\Omega }_X)=\underset{ยฏ}{\mathrm{Hom}}_{D_X}(_{\mathrm{๐ณ๐บ๐๐พ}},D_X\mathrm{\Omega }_X)`$.
### 5.3. Construction of the left $`๐^{}()`$ DG-module for a chiral $``$-module.
Note that the Lie algebra in the category of D-modules $`๐ข=p_1{}_{}{}^{}\left(_{\stackrel{~}{\mathrm{\Delta }}}\right)`$ acts naturally on an arbitrary chiral $``$-module $``$ as follows:
$$p_1{}_{}{}^{}\left(_{\stackrel{~}{\mathrm{\Delta }}}\right)\stackrel{!}{}p_1{}_{}{}^{}\left(\left(_{\stackrel{~}{\mathrm{\Delta }}}\right)\stackrel{!}{}p_1^{}()\right)p_1{}_{}{}^{}\left(\left(_{\stackrel{~}{\mathrm{\Delta }}}\right)\stackrel{!}{}(\mathrm{\Omega })\right)$$
$$p_1{}_{}{}^{}\left(()_{\stackrel{~}{\mathrm{\Delta }}}\right)\stackrel{๐}{}p_1{}_{}{}^{}\mathrm{\Delta }_{!}^{}()\stackrel{~}{}.$$
Here $`()_{\stackrel{~}{\mathrm{\Delta }}}`$ denotes the completion of the D-module $`j_{}j^{}()`$ along the diagonal in $`X\times X`$.
Note that the D-module $`\mathrm{\Delta }_!()`$ is locally finite over the ideal of the diagonal. Thus the completion of the map $`j_{}j^{}()\mathrm{\Delta }_!()`$ is well defined.
#### 5.3.1.
For a chiral $``$-module $``$ consider it as a Lie-\* module and recall its naive de Rham complex $`C^{}(,)=\mathrm{\Lambda }^{}(^{})\stackrel{!}{}`$. Lemma: The complex $`C^{}(,)`$ has a natural structure of a left DG-module over the DG Lie algebroid $`๐^{}()`$.
Proof. The statement of the Lemma follows from the existence of the $`๐ข`$-module structure on $``$ introduced in the beginning of the present subsection. โ
Here we come to the crucial point explaining the phenomenon of the Tate extension in the semiinfinite cohomology of Lie-\* algebras. What we would like to do is to consider the homological Chevalley complex of the DG Lie algebroid $`(๐^{}(),^{}(),\mathrm{})`$ with coefficients in $`C^{}(,)`$. Yet there is no naive way to do it. Somehow we have to make $`M\mathrm{\Lambda }^{}(^{})`$ into a right DG-module over our DG Lie algebroid in the category of $`D_X`$-modules. Remark: Mimicking the Tate Lie algebra case we could consider the DG Lie algebroid $`๐^{}()^{}()\mathrm{๐}`$, then construct its antipode $`\alpha `$ not commuting with the differential. However the point is that the obtained DG Lie algebroid in the category $`\mathrm{mod}\mathrm{-}D_X`$ does not come from a central extension of our Lie-\* algebra $``$. Thus there is no way to construct a left DG-module over $`(๐^{}()^{}()\mathrm{๐})^{\mathrm{opp}}`$ (with the differential twisted by the antipode) starting from a chiral module either over $``$ or over some its central extension. Instead one should act as follows.
### 5.4. Antipode map for the Tate extension of $`๐^{}()`$.
Consider the (Tate) extension of Lie algebras in D-modules on $`X`$:
$$0\mathrm{\Omega }_X_{_{\mathrm{๐ณ๐บ๐๐พ}}}0.$$
Here $`_{_{\mathrm{๐ณ๐บ๐๐พ}}}`$ denotes $`(_{_{\mathrm{๐ณ๐บ๐๐พ}}})`$. We have also the corresponding extension of DG Lie algebroids on $`X`$:
$$0\mathrm{\Omega }_X\stackrel{!}{}^{}()_{_{\mathrm{๐ณ๐บ๐๐พ}}}\stackrel{!}{}^{}()\stackrel{!}{}^{}()0.$$
Denote $`_{_{\mathrm{๐ณ๐บ๐๐พ}}}\stackrel{!}{}^{}()`$ by $`๐^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}})`$.
We will need also the universal enveloping algebras of the DG-Lie $`^{}()`$-algebroids $`๐^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}})`$ and $`๐^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}})`$. Keeping the notation from the previous section we denote these associative algebras in $`\mathrm{mod}D_X`$ by $`\stackrel{~}{๐ฐ}_{}(๐^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}}))`$ and $`๐ฐ_{}(๐^{}())`$ respectively. Let $`๐ฐ_{}(๐^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}}))`$ be the qotient of $`\stackrel{~}{๐ฐ}_{}(๐^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}}))`$ by the ideal generated by the relation $`^{}()\stackrel{!}{}(\mathrm{\Omega }_X\mathrm{๐})=^{}()`$. Here the LHS of the equality is the kernel of the DG Lie algebroid extension map while the RHS is the unit $`^{}()\stackrel{~}{๐ฐ}_{}(๐^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}}))`$.
Note that with the differentials forgotten the algebras $`๐ฐ_{}(๐^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}}))`$ and $`๐ฐ_{}(๐^{}())`$ are isomorphic.
#### 5.4.1.
We introduce an antipode map $`๐^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}})\stackrel{~}{}๐^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}})^{\mathrm{opp}}`$ as follows. Set $`\alpha (br)=br+\mathrm{co}\mathrm{ad}_b(r)\mathrm{๐}`$. Here $`b`$ is a section of $`_{_{\mathrm{๐ณ๐บ๐๐พ}}}`$, $`r`$ is a section of $`^{}()`$ and $`\mathrm{๐}`$ is the generating section of $`\mathrm{\Omega }_X`$.
Note that the antipode does not necesserily commute with the differential on $`๐^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}})`$. Lemma:
* $`\alpha `$ is well defined as a antipode of a Lie superalgebra in the category of $`D_X`$-modules $`๐^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}})`$ (with the differential forgotten).
* When restricted to any open affine subset $`\stackrel{}{X}X`$ the DG Lie superalgebra in the category of $`D_X`$-modules $`๐^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}})^{\mathrm{opp}}|_\stackrel{}{X}`$ with the differential $`\alpha d_๐^{}\alpha ^1`$ is isomorphic to the DG Lie superalgebra in the category of $`D_X`$-modules $`๐^{}(\mathrm{\Omega }_X)|_\stackrel{}{X}`$. Here $`\mathrm{\Omega }_X`$ denotes the trivial central extension.
Proof. Both statements of the Lemma follow from the corresponding local statements presented in \[Ar2\], Proposition 4.3.4. โ
Here we come to another difference with the Tate Lie algebra semiinfinite cohomology case. While previous Lemma states that when restricted to an open affine subset the complex of D-modules $`๐^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}})^{\mathrm{opp}}`$ is isomorphic to $`๐^{}()^{}()\stackrel{!}{}(\mathrm{\Omega }_X\mathrm{๐})`$, still possibly the short exact sequence of D-modules on the complete curve
$$0\mathrm{\Omega }_X\stackrel{!}{}^{}()(_{_{\mathrm{๐ณ๐บ๐๐พ}}}\stackrel{!}{}^{}())^{\mathrm{opp}}\stackrel{!}{}^{}()0$$
does not split.
That is where we use the universal enveloping algebras of our DG Lie algebroids. Extend the antipode $`\alpha `$ to $`\stackrel{~}{๐ฐ}_{}(๐^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}}))`$.
#### 5.4.2.
Proposition:
* $`\alpha `$ descends to the antipode of $`๐ฐ_{}(๐^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}}))`$.
* The DG-algebra $`๐ฐ_{}(๐^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}}))^{\mathrm{opp}}`$ with the differential twisted by the antipode $`\alpha `$ is isomorphic to the DG-algebra $`๐ฐ_{}(๐())`$.
Proof. Follows from local calculations in the Tate Lie algebra semiinfinite cohomology case (see \[Ar2\]). Corollary: Any left DG-module over the DG Lie algebroid $`๐^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}})`$ on which the center $`\mathrm{\Omega }_X\mathrm{๐}`$ acts by unity becomes a right DG-module over the DG Lie algebroid $`๐^{}()`$. โ
### 5.5. Standard semiinfinite complex for a chiral module over Tate extension of the Lie-\* algebra $``$.
For a chiral $`_{_{\mathrm{๐ณ๐บ๐๐พ}}}`$-module $``$ such that the center $`\mathrm{\Omega }_X\mathrm{๐}`$ acts on it by unity consider the naive de Rham complex of the module $`C^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}},)=\stackrel{!}{}\mathrm{\Lambda }(^{})`$ as a right DG-module over the DG-Lie algebroid in the category of $`D_X`$-modules $`(\stackrel{!}{}^{}())`$. The right DG-module structure is obtained using the antipode construction from the previous subsection.
#### 5.5.1.
Definition: We call the homological Chevalley complex of the DG Lie algebroid $`((\stackrel{!}{}^{}()),^{}(),\mathrm{})`$ with coefficients in the right DG-module $`C^{}(_{_{\mathrm{๐ณ๐บ๐๐พ}}},)`$ the standard semiinfinite complex for the chiral module $``$ over the Lie-\* algebra $`_{_{\mathrm{๐ณ๐บ๐๐พ}}}`$ and denote it by $`C^\frac{\mathrm{}}{2}(_{_{\mathrm{๐ณ๐บ๐๐พ}}},)`$. Remark: Note that as the D-module on $`X`$ the constructed complex looks as follows:
$$C^\frac{\mathrm{}}{2}(_{_{\mathrm{๐ณ๐บ๐๐พ}}},)=\mathrm{\Lambda }_{^{}}^{}(p_1{}_{}{}^{}(()_{\widehat{\mathrm{\Delta }}}\stackrel{!}{}^{}())_{^{}}\left(\stackrel{!}{}\mathrm{\Lambda }^{}(^{})\right)$$
$$=\mathrm{\Lambda }^{}(p_1{}_{}{}^{}(()_{\widehat{\mathrm{\Delta }}})\stackrel{!}{}\mathrm{\Lambda }^{}(^{})\stackrel{!}{}.$$
Here as before $`p_1`$ denotes the projection $`X\times XX,(x,y)x`$.
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# The Jet of 3C 273 observed with ROSAT HRIBased on observations with the ROSAT X-ray satellite and also on data collected with the VLA (The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc.) and at the European Southern Observatory, Chile (ESO Nโ 51-2-021).
## 1 Introduction
Jets in extragalactic radio sources play a central rรดle in our understanding of the nature of these enigmatic sources (Begelman *et al.* BBR84 (1984), Rรถser *et al.* RingbergII (1993)) as they mark the channels through which mass, energy and momentum are transported out from the nucleus into the extended radio lobes. The detailed physical conditions in the jets are still unknown since their synchrotron emission at radio frequencies provides little information about the emitting plasma. However, of the more than 100 extragalactic radio jets known (Bridle & Perley BP84 (1984)) only three are readily detectable at frequencies higher than the radio band. The two objects with the brightest optical jet emission are M 87, a radio galaxy, and 3C 273, a quasar. Due to this wide wavelength coverage we can expect that basic information about the physical conditions giving rise to the synchrotron emission can be derived, most importantly the maximum energy of the radiating particles.
We have therefore embarked on a detailed study of the jet of 3C 273 at the best angular resolutions currently available across the electromagnetic spectrum using the VLA and MERLIN at radio wavelengths, HST at optical and near-infrared and ROSAT at X-ray wavelengths (for a preliminary presentation of some of these data see Rรถser *et al.* HJRMNCDP97p (1997)). Whereas the synchrotron origin of the radio, infrared and optical emission is now firmly established (Conway & Rรถser CHJR93 (1991), Rรถser & Meisenheimer HJRM91 (1991)) it is the origin of the X-rays that is least understood. The most detailed X-ray study of this jet is due to Harris & Stern (HS87x (1987)), who carefully analysed an EINSTEIN HRI observation with 95 ksec integration time. Although they marginally detected a signal from the jet, none of their attempts to interpret the X-ray emission proved satisfactory. Our ROSAT observations were primarily aimed at verification and interpretation of the jetโs X-ray emission.
## 2 The jet of 3C 273 at radio to optical wavelengths
As indicated in Figure 1, the radio jet is detected all the way from the core out to the hot spot. The faint optical structure, however, although coinciding in position angle with the line joining radio components A (hot spot) and B (core), seems to be detectable only outwards of $`11`$โณ. It also terminates about 1โณ before the peak in the radio hot spot is reached. At the quasarโs redshift of 0.158 the projected length is $``$60 kpc<sup>1</sup><sup>1</sup>1We assume *q*<sub>0</sub> = 0.1 and *H*<sub>0</sub> = 65 km/sec/Mpc, so 1โณ corresponds to 2.67 kpc. Greenstein & Schmidt (GS64 (1964)) in their detailed study of 3C 273 and 3C 48 briefly discuss also this jet. Their spectrum of its outer end exhibited a featureless blue continuum and they assumed the optical radiation is of the same origin as the radio emission, *i.e.* synchrotron radiation from relativistic particles. This was proven by Rรถser & Meisenheimer (HJRM86 (1986), HJRM91 (1991)) using optical polarimetry. Further hints about the synchrotron emission are provided by studies of the spatially resolved continua of individual knots in the jet. Meisenheimer & Heavens (MH86 (1986)) present a simple model describing the global shape of the continuum including an exponential cut-off observed in the hot spot at the jetโs outer end reflecting the maximum energy gained by the relativistic particles. Applying this model to the other knots in the jet indicates that we also see distributions of relativistic particles truncated at some maximum energy, except for the innermost visible knot (Meisenheimer *et al.* MNHJR96p (1996)), which essentially has no cut-off at all. In view of these results the marginal detection of X-ray emission by Harris & Stern (HS87x (1987)) would naturally be associated with this innermost knot, although they place the centroid of the X-ray emission further out.
We have therefore used the ROSAT HRI with its better spatial resolution and sensitivity to verify the X-ray emission from this jet and to gain further insight into the emission mechanism.
## 3 ROSAT Observations
The data were collected in two observing cycles. In a 17.991 ksec integration (January 1992) we clearly detected the jet without any image processing (see Figure 2). As the signal-to-noise ratio (S/N) was not sufficient for detailed studies, 3C 273 was re-observed in December 1994/January 1995 for a total of 68.154 ksec (quoted times as โacceptedโ by the ROSAT standard reduction analysis).
### 3.1 Improving the resolution
Images (sky pixel size 0$`\stackrel{}{.}`$5) have been produced from the event tables using the EXSAS package provided from the ROSAT group at MPE in Garching. However, the resulting point-spread function turned out to be unsatisfactory. Whereas in the call for proposals the integral point response function of the ROSAT-XTE + HRI was quoted to have a width of 5โณ, the FWHM of the images as produced from the raw data was typically larger than 6โณ. This discrepancy is due to inaccuracies in the aspect solution, which determines the sky coordinates for each photon detected. Obviously in the standard aspect solution the spacecraft wobble with a period of 402 sec is not completely corrected for (see Figure 3).
With a count rate in the HRI of 2.8 cts/sec the signal from the quasar core itself is sufficiently strong to allow a shift-and-add procedure as follows: All integrations are divided into time bins of 10 sec duration and the centre of gravity of the photons from the quasar core are calculated for each interval. Photons were restricted to a raw amplitude between 2 and 8, typical for the quasar. The average accuracy of the centroid positions was 0$`\stackrel{}{.}`$64 in X and 0$`\stackrel{}{.}`$70 in Y-direction. The interval of 10 sec was a compromise between sufficient time resolution and positional accuracy. These offsets from the nominal centre position as a function of time directly correspond to the remnant pointing errors due to the insufficiently corrected wobble motion. They were interpolated in time by splines and the detected position of every photon was corrected as a function of its arrival time.
For comparison the same procedure was applied to eight data sets from the ROSAT archive for the white dwarf HZ 43, which is definitively a point source. To derive the radial image profile, each data set was sampled on concentric circles around the quasar respectively white dwarf with radii in steps of 0$`\stackrel{}{.}`$5. For each circle a constant was fitted to the data giving the azimuthally averaged intensity profile. The โcoreโ of these profiles was decomposed by a least-square-fitting procedure into two Gaussians, neglecting the very extended exponential component discussed in the HRI calibration report (David *et al.* DHK99 (1999)). The shape of the point-spread-function (PSF) in the raw images changed from observation to observation. But as can be seen from Table 1 the profile of the de-jittered images was constant within the error. As indicated in Table 1 and demonstrated by Figure 5 and Figure 6, the resolution could be enhanced to 4$`\stackrel{}{.}`$5 this way. A similar procedure was described by Morse (Morse94 (1994)).
From this analysis we conclude that there is no discernable difference in the core profile between quasar and white dwarf in the averaged intensity profiles, *i.e.* the quasar core is unresolved at X-rays. Only in the halo the intensity of the quasar is slightly above the normalized white dwarf profile.
### 3.2 Fitting the quasarโs point-spread-function
The X-ray signal from the jet of 3C 273 is very weak and is located in the wings of the complicated point-spread-function of the bright quasar core (total intensity ratio 400:1). Therefore the X-ray emission from the jet has to be isolated from this underlying point-spread-function background. Using the azimuthally averaged profiles derived above to analyse the point-spread-function is not sufficiently accurate to isolate the jetโs weak signal and measure its flux. A better model of the point-spread-function had to be obtained by fitting the signal sampled along concentric circles in four sections by polynomials of order 3, with section boundaries at 30, 120, 210 and 300. In this way a satisfactory flat background also interior to the optically visible jet region is achieved (Figure 7). It should be pointed out that the result of the background modelling does strongly depend on the model parameters. With the current resolution of 4$`\stackrel{}{.}`$5 and the steep intensity gradients towards the quasar core it cannot be completely excluded that the relatively complicated model absorbs a moderately extended component in the inner part of the jet ( $`r<10\mathrm{}`$).
### 3.3 The X-ray signal from the jet
The count rate of the X-ray emission from the jet was derived from the PSF-subtracted image. In a window encompassing the jet we measured 644 counts. The background in several windows of the same size and at the similar distance to the quasar was measured to be $`(28.6\pm 13.5)`$ counts. We therefore deduce a total jet flux of 610 counts in 85068 sec. The same procedure yields 253 counts for the object to the NE of the quasar. To convert this to a flux density we integrated a synchrotron power-law *f*$`{}_{\nu }{}^{}`$*$`\nu `$*<sup>ฮฑ</sup> over the effective collecting area of the ROSAT HRI as a function of frequency (effective energy 1.17 keV corresponding to $`2.83\times 10^{17}`$ Hz) including absorption due to a neutral hydrogen column of $`1.8\times 10^{20}`$ cm<sup>-2</sup> (Otterbein Otterbein92 (1992)) using the cross-sections from Morrison & McCammon (MMcC83 (1983)) for a solar abundance. To represent synchrotron emission we used spectral indices $`\alpha =1`$ and $`2`$ and thus obtained an integral flux density of $`(36.6\pm 2.8)`$ and $`(14.9\pm 1.1)`$ nJy for the jet and $`(15.2\pm 1.1)`$ and $`(6.2\pm 0.5)`$ nJy for the object in the NE of the quasar. If we assume a spectral index of $`\alpha =0`$ (resembling thermal bremsstrahlung) the corresponding values are $`(57.1\pm 4.3)`$ nJy for the jet and $`(23.7\pm 1.8)`$ nJy for the object in the NE. The error was formally derived from the scatter in the background determination.
Inspection of Figure 7 suggests that the X-ray emission from the jet is extended in the radial direction all along the visible jet. We therefore rotated the background subtracted image by 47$`\stackrel{}{.}`$8 around the quasar core to sum the jetโs signal over the 20 pixel rows encompassing it. The resulting trace along the jet is shown in Figure 8. This run of the X-ray emission along the jet was analysed in a similar manner as the optical and radio emission (Rรถser & Meisenheimer HJRM91 (1991)): Gaussian profiles with constant width of 4$`\stackrel{}{.}`$5 were placed at the positions of knots A, B, C, D and H<sup>2</sup><sup>2</sup>2Nomenclature introduced by Leliรฉvre *et al.* (LNHRS84 (1984)).. Their peak was adjusted via a least square fit to represent the trace of X-ray emission along the jet (Table 2).
The strongest X-ray signal originates from the positions of knots A and B, much weaker emission is found further out, possibly out to the hotspot H. Due to the steeply rising quasar background, and to uncertainties in background modelling mentioned above, the onset of the jetโs X-ray emission is somewhat uncertain. On the basis of our data no X-ray emission is detected from the jet inwards of knot A.
The X-ray source to the NE of the quasar does not show any optical counterpart on our deep images, neither in the radio nor in the optical. Due to its weakness we cannot say if it is extended or not in the X-rays. Its relation to 3C 273 and its nature remain unknown.
## 4 Origin of the X-ray emission
To shed light onto the possible emission mechanism(s) giving rise to the jetโs X-ray emission, the above data were combined with results from our radio-to-optical studies at 1$`\stackrel{}{.}`$3 resolution (Meisenheimer *et al.* MYHJR97 (1997), Neumann *et al.* NMHJR97 (1997), Rรถser & Meisenheimer HJRM91 (1991)) to investigate the run of the continuum across the electromagnetic spectrum for the different regions along the jet. We plot the spectra of knots A, B, C, D, and H in Figure 9, where the flux was multiplied by different factors to disentangle the individual spectra in the graph. The continua are model spectra as described by Heavens & Meisenheimer (HM87 (1987)) and Meisenheimer & Heavens (MH86 (1986)). They result from Fermi acceleration of particles in a strong shock and include the radiation losses in a finite down-stream region. The exponential cut-off reflects the maximum energy obtained by the particles and is well established for knots B to H also by our recent HST WFPC2 data at 300 nm (Jester *et al.*, in preparation). It is evident that in general the X-ray flux level is not a continuation of the radio-optical synchrotron cut-off continuum. Therefore different emission mechanisms have to be discussed for the individual knots.
### 4.1 Synchrotron radiation from the jet
Only for knot A does the extrapolation of the radio-to-optical continuum approximately meet the X-ray flux level. Extrapolation of the $`\lambda `$6cm flux with the low frequency spectral index of knot A ($`\alpha =0.83)`$ into the ROSAT range predicts an X-ray flux of 32 nJy, only a factor of roughly two above the observed level. Whereas in Meisenheimer *et al.* (MNHJR96p (1996)) a lower limit to the cut-off frequency for knot A could be set at about $`5\times 10^{16}`$ Hz, we now have to increase this value by about a factor of 50 (see Figure 9). According to standard synchrotron theory we can therefore place a new lower limit to the maximum energy of the relativistic particles in knot A of
$$\gamma _c=10^6\times \sqrt{\frac{\nu _c}{1.26\times 10^{15}Hz}\mathrm{/}\frac{B_{}}{30nT}}=3\times 10^7$$
where we have used the minimum energy field of 67 nT. For the other knots the primary synchrotron component producing the observed radio-to-optical continuum exhibits an exponential cut-off in the optical/UV-range (Meisenheimer *et al.* MNHJR96p (1996)). To maintain the synchrotron scenario also for these regions a second particle population with higher maximum energy has to be invoked<sup>3</sup><sup>3</sup>3Such a second synchrotron component was also suggested by Harris *et al.* (HHSSV99 (1999)) to account for the X-ray emission from the jet in 3C 120.. These populations are indicated in Figure 9 connecting to the ROSAT HRI points with an assumed spectral index of $`\alpha =0.82`$, the low-frequency average for knots A to D. The required density of relativistic particles in these hypothetical components decreases outwards along the jet. It is a factor of 5 below the density of the particles producing the observed radio-to-optical continuum for knot B and a factor of 15 and 200 below that in knot C and D, respectively. The small fraction of relativistic particles make it very hard to detect them at *e.g.* optical wavelengths. However, if this second population is confined to some centres of very effective acceleration on sub-arcsecond scales, we expect to see deviations from the cut-off spectrum in the optical spectral index map derived from our HST R- and U-band data at a resolution of 0$`\stackrel{}{.}`$2 (Jester *et al.* 2000, in preparation).
### 4.2 Synchrotron Self-Compton emission (SSC) from the jet
Prime sources for inverse Compton emission are compact regions with high radio photon densities in the jet, for which the hot spot H is the most likely candidate. The size of the emission region and the radio flux originating from it determine the amount of synchrotron-self-Compton emission. We have calculated the inverse Compton emission along the lines given by Blumenthal & Tucker (BT70 (1970)) as follows: The most compact component certainly is the hot spotโs acceleration region, where the optically radiating particles are produced. Its contribution to the synchrotron spectrum can be inferred from the low-frequency power-law and the high-frequency cut-off. The intermediate range with the break in the continuum (see Figure 9) is dominated by the superposition of the down-stream regions, where the relativistic particles already have lost part of their energy. Connecting the cut-off part smoothly with a power-law of index $`0.60`$ we estimate the 5 GHz-flux from the acceleration region itself to be about 0.1 Jy. From the most recent analysis of the hot spotโs spectrum by Meisenheimer *et al.* (MYHJR97 (1997)) we infer a thickness of the emission region close to the internal shock (Mach disk) of 1.4 pc, the width perpendicular to the jet (from our best-resolved radio map) is taken to be 2.2 kpc (circular cross-section assumed). For a minimum energy field of 35 nT (assumed constant over the entire hot spot in the model) we set the density in relativistic particles in this volume to reproduce the above estimated $`\lambda `$6cm flux of 0.1 Jy. Integrating this synchrotron emission over the range 10 MHz to $`5\times 10^4`$ GHz (corresponding to Lorentz-factors of 100 to $`2\times 10^5`$) produces an inverse Compton flux of about 1 nJy, to within factors of 2 to 4 what is observed (see Table 2). As this is only a rough estimate to check the order of magnitude the discrepancy could well be removed by fine-tuning the parameter assumed (filling factor, geometry ).
For the other knots synchrotron-self-Compton emission fails to meet the observed level by orders of magnitude, *e.g.* for knot A we expect an inverse Compton X-ray flux of only 0.01 nJy. As the photon density of the microwave background radiation is one order of magnitude less than the synchrotron photon density in all knots, its photons cannot account for the X-ray flux via inverse Compton scattering either.
### 4.3 Thermal Bremsstrahlung from the jet
If future high resolution observations fail to reveal locations of $`\nu _c10^{15}`$ Hz in knots B, C, and D, there remains the bremsstrahlung emission from a thermal plasma as a last resort to explain the jetโs X-ray emission outside knot A and the hot spot. A plasma at 10<sup>8</sup> K and with an electron density of 1 cm<sup>-3</sup> spread over the volume of a typical jet knot ($`0\stackrel{}{.}5\times 0\stackrel{}{.}2`$ from our HST images, Jester *et al.* 2000, *in preparation*) would produce an X-ray flux in the ROSAT window corresponding to 0.01 nJy. Even with this unrealistically high electron density this is orders of magnitude below the observed X-ray flux level. Furthermore any sufficiently dense thermal plasma would result in total depolarisation of the jetโs radio emission and in a substantial rotation measure. Both are not observed (Conway *et al.* CGPB93 (1993)). So thermal bremsstrahlung is highly unlikely to account for the observed X-ray flux.
## 5 The X-ray halo of 3C 273
For the comparison of the quasarโs profile with that of the white dwarf, the normalization constant had been optimised over the range 3โณ to 13โณ, where the profiles show excellent agreement. However, it is evident from Figure 5 that the quasar profile lies systematically above the profile of the stellar source beyond radii of about 15โณ. Although the error of individual data points is large, we regard the systematic deviation as real and attribute it to an extended X-ray halo of the quasar (3C 273 had been observed on-axis, the pointing for HZ 43 was 1โฒ off-axis). To isolate the X-ray flux from this halo we have used the scaling from Figure 5 and subtracted the azimuthally averaged HZ 43-profile from the averaged quasar profile (Figure 10). This differential profile was fitted with a King profile (Jones & Forman JF84 (1984))
$$P(r)=P(0)\left[1+\left(r/a\right)^2\right]^{3\beta +0.5}$$
with *P*(0) = 3.6, *a* = 10$`\stackrel{}{.}`$8 $`\widehat{=}`$ 28.8 kpc, *$`\beta `$* = 0.6758 determined from a least square fit to the data points between 12$`\stackrel{}{.}`$5 and 80โณ from the core. Although the formal errors for these parameters are large, we nevertheless used them for our numerical estimates, as they describe the data reasonably well. The King-profile was integrated up to give the total X-ray flux from the halo of 232 nJy corresponding to a luminosity in the HRI band of $`3.4\times 10^{43}`$ erg/sec (spectral index $`\alpha =0`$ assumed). To convert this into an estimate of the physical parameters we assumed a plasma temperature of $`2.7\times 10^7`$ K as in the M 87 halo (Bรถhringer *et al.* BBSVHT94 (1994)) and that the plasma is confined with constant density to within the core radius. From these assumptions a density of $`6\times 10^4`$ H-atoms/m<sup>3</sup> is derived for the central part of the halo. This is a factor of 10 above the upper limit derived from EINSTEIN observations by Willingale as quoted by Conway *et al.* (CDFR81 (1981)). The typical density of the thermal plasma in the cores of rich clusters of galaxies is a factor of 10 below this value (Jones & Forman JF84 (1984)) but the density derived above for the halo of 3C 273 matches the central density of $`(7\pm 2)\times 10^4`$ m<sup>-3</sup> derived for the central density derived for the intra-cluster medium around Cygnus A (Carilli *et al.* CPH94 (1994)).
## 6 Conclusions
ROSAT HRI observations confirm the X-ray emission from the jet of 3C 273 as suspected on the basis of the EINSTEIN HRI observations. They furthermore show that the X-rays originate from all along the jet, probably including the hot spot. While the jetโs radio emission is highly peaked towards the outer end, the optical emission is more or less constant along the jet. This trend with wavelength continues at X-rays in that these are strongly peaked now at the inner end (knots A/B).
Despite the considerably improved data base the problems with the interpretation of the jetโs X-ray emission (see Harris & Stern HS87x (1987)) still remain. Whereas the X-ray emission from the jet of M 87 might well be due to the synchrotron emission process (Neumann *et al.* NMHJRF97 (1997), Rรถser & Meisenheimer RingbergIV (1999)), the situation is less clear for the jet of 3C 273. An extrapolation of the radio-to-optical synchrotron continuum could only explain the X-ray emission from the innermost knot A, it fails for the rest of the jet. Any X-ray emission from the hot spot up to the level we found can be accounted for by synchrotron-self-Compton emission. The emission mechanism for knots B to D remains a mystery, as โ except for an hypothetical high-energy synchrotron component โ all three mechanisms discussed above cannot account for the X-ray emission.
The forthcoming observations of 3C 273 by CHANDRA will provide X-ray data at higher resolution and with spectral information. High spatial resolution will test if the X-ray emission does indeed coincide with the radio-optical synchrotron continuum emission, not necessarily the case for thermal bremsstrahlung emission. X-ray spectra will set important constraints on the emission mechanism mainly via the spectral slope in the X-ray band, which could directly reveal a second population of relativistic particles emitting synchrotron X-rays. Thus we can expect further insight into the mystery of X-ray emission from the jet of 3C 273 within the immediate future.
###### Acknowledgements.
We thank G. Hasinger for discussions on the ROSAT data reduction and his suggestion to try re-centring the photons. Valuable advice from C. Izzo and S. Dรถbereiner is also kindly acknowledged.
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# Untitled Document
hep-th/0005046 IASSNSโHEPโ00/37, PUPT-1929
Duality and Non-Commutative Gauge Theory
Ori J. Ganor<sup>1</sup> origa@viper.princeton.edu, Govindan Rajesh<sup>2</sup> rajesh@sns.ias.edu and Savdeep Sethi<sup>3</sup> sethi@sns.ias.edu
$``$ School of Natural Sciences, Institute for Advanced Study, Princeton, NJ 08540, USA
$``$ Department of Physics, Jadwin Hall, Princeton, NJ 08544, USA
We study the generalization of S-duality to non-commutative gauge theories. For rank one theories, we obtain the leading terms of the dual theory by Legendre transforming the Lagrangian of the non-commutative theory expressed in terms of a commutative gauge field. The dual description is weakly coupled when the original theory is strongly coupled if we appropriately scale the non-commutativity parameter. However, the dual theory appears to be non-commutative in space-time when the original theory is non-commutative in space. This suggests that locality in time for non-commutative theories is an artifact of perturbation theory.
5/00
1. Introduction
Non-commutative gauge theory provides an interesting class of examples in which to explore the effects of spatial non-locality. While it is easy to define the classical non-commutative gauge theory, it is much harder to determine whether the quantum theory exists. Since non-commutative gauge theories arise in particular string theory backgrounds, we know that these theories can be embedded consistently in string theory. The decoupling argument of Seiberg and Witten suggests that some of these theories might exist as quantum theories independent of string theory.
We are primarily interested in four-dimensional gauge theories. Our goal is to understand how S-duality \[3,,4\] generalizes to non-commutative gauge theory. The generalization is not a straightforward consequence of S-duality in type IIB string theory. To see this, let us begin by briefly recalling how S-duality of N=4 Yang-Mills arises from string theory. In the limit $`\alpha ^{}0`$, the theory on coincident D3-branes is N=4 Yang-Mills. For simplicity, we set the RR scalar $`C^{(0)}`$ to zero. The gauge theory coupling constant, $`g^2`$, is then related to the closed string coupling constant $`g_s=e^\varphi `$:
$$\frac{g^2}{4\pi }=g_s.$$
$`(1.1)`$
The conjectured $`SL(2,\text{ZZ})`$ symmetry of string theory then descends to an $`SL(2,\text{ZZ})`$ symmetry of the field theory.
To obtain non-commutative Yang-Mills, we consider a system of coincident D3-branes with NS-NS $`B`$-field non-zero along the brane. In the decoupling limit , the theory on the brane has a coupling constant related to the open string coupling constant, $`G_s`$, rather than the closed string coupling:
$$g^2=2\pi G_s.$$
$`(1.2)`$
In the decoupling limit, the closed string coupling constant goes to zero while $`G_s`$ remains finite and dependent on the $`B`$-field. In this case, S-duality of the closed string theory does not descend to a symmetry of the field theory.
For a $`U(1)`$ gauge theory, S-duality can be demonstrated directly with a purely field theoretic argument. We start with the Minkowski space action,<sup>1</sup> We use $``$ to denote the Hodge dual of a form rather than the star product.
$$S=\frac{1}{4g^2}FF,$$
$`(1.3)`$
where $`F=dA`$ is the field strength. We want to perform a Legendre transformation with respect to $`F`$. To implement the Bianchi identity,
$$dF=0,$$
we introduce a dual gauge-field $`A_D`$,
$$S=(\frac{1}{4g^2}FF+\frac{1}{2}A_DdF).$$
$`(1.4)`$
We can now treat $`F`$ as an independent variable and perform the path-integral over $`F`$. This amounts to solving the field equations for $`F`$ which gives the relation,
$$dA_D=\frac{1}{g^2}F,$$
$`(1.5)`$
and the resulting dual action,
$$S=\frac{g^2}{4}F_DF_D.$$
$`(1.6)`$
The aim of this discussion is to generalize this purely field theoretic argument to the non-commutative rank one theory. Unlike ordinary abelian gauge theory, the coupling constant cannot be scaled away even for the rank one non-commutative theory.
In the following section, we explicitly show that the non-commutative action expressed in terms of a commutative gauge-field contains only powers of $`F`$ to order $`\theta ^2`$. In particular, the gauge-field does not appear explicitly. It is not hard to argue that this must be true to all orders in $`\theta `$. This implies that we can obtain a dual description by Legendre transforming with respect to $`F`$. The resulting dual theory is classical since we neglect loops. However, to order $`\theta `$, we will see that no loops appear and the quantum and semi-classical dual descriptions agree. To order $`\theta ^2`$, loops appear and the bosonic theory needs to be regulated. At this point, the computation should be performed in the full N=4 theory.
Fortunately, our primary observations are already visible at order $`\theta `$. We find that under the duality transformation,
$$\theta \stackrel{~}{\theta }=g^2(\theta ).$$
$`(1.7)`$
That this transformation does not square to one is not so surprising since $`(S)^2`$ is not the identity operation but charge conjugation. We will also find that $`\stackrel{~}{\theta }`$ must be held fixed if the dual theory is to have a perturbative expansion in $`1/g`$. Even more interesting is the observation that if $`\theta `$ is purely spatial then $`\stackrel{~}{\theta }`$ involves a space direction and a time direction. The theory becomes non-commutative in space-time. Although we will not obtain the complete quantum dual description, it seems clear that this feature, visible at leading order in $`\theta `$, persists to higher orders. Space-time non-commutative theories are highly unusual; see for a recent discussion. Our result suggests that we cannot avoid studying these theories if we are to understand theories which perturbatively have only spatial non-commutativity.
2. The Duality Transformation
2.1. Rewriting the non-commutative Lagrangian
The non-commutative theory is defined by the action,
$$S=\frac{1}{4g^2}\widehat{F}\widehat{F}.$$
$`(2.1)`$
The change of variables given in allows us to express $`\widehat{F}`$ in terms of a commutative gauge-field $`A`$. We assume that $`\theta `$ is purely spatial. The relation takes the form,
$$\widehat{F}=F+T_\theta (A)+T_{\theta ^2}(A)+\mathrm{}.$$
$`(2.2)`$
The terms of order $`\theta `$ are given by,
$$T_\theta (A)=F\theta FA_k\theta ^{kl}_lF.$$
$`(2.3)`$
We follow the notation of where $`F\theta F=F_{ik}\theta ^{kl}F_{kj}`$. The expression for $`T_{\theta ^2}(A)`$ is found in ,
$$\begin{array}{cc}\hfill T_{\theta ^2}(A)=& F\theta F\theta F+\frac{1}{2}A_k\theta ^{kl}\left(_lA_m+F_{lm}\right)\theta ^{mn}_nF\hfill \\ & +\theta ^{kl}A_k_l\left(F\theta F\right)+\frac{1}{2}\theta ^{kl}\theta ^{mn}A_kA_m_l_nF.\hfill \end{array}$$
$`(2.4)`$
The expression for $`\widehat{F}`$ explicitly contains $`A`$. However, we can manipulate the action (2.1) so that it takes the following form,
$$S=\frac{1}{4g^2}(FF+L_\theta (F)+L_{\theta ^2}(F)+\mathrm{}).$$
$`(2.5)`$
The terms of order $`\theta `$ take the form,
$$L_\theta (F)=2\mathrm{tr}(\theta F^3)\frac{1}{2}\mathrm{tr}(\theta F)\mathrm{tr}(F^2),$$
$`(2.6)`$
where we define $`\mathrm{tr}(AB)=A_{ij}B^{ji}.`$ Since our theory is rank one, there should be no confusion with traces over group indices. It is not too hard to find an expression for $`L_{\theta ^2}(F)`$ which takes the form:
$$\begin{array}{cc}\hfill L_{\theta ^2}(F)=& 2\mathrm{tr}(\theta F\theta F^3)+\mathrm{tr}(\theta F^2\theta F^2)+\mathrm{tr}(\theta F)\mathrm{tr}(\theta F^3)\hfill \\ & \frac{1}{8}\mathrm{tr}(\theta F)^2\mathrm{tr}(F^2)+\frac{1}{4}\mathrm{tr}(\theta F\theta F)\mathrm{tr}(F^2).\hfill \end{array}$$
$`(2.7)`$
While we have explicitly demonstrated that it is possible to express (2.1) in terms of $`F`$ to order $`\theta ^2`$, it must be the case to all orders in $`\theta `$. The only gauge-invariant operator that can be constructed from $`A`$ is $`F`$. While $`\widehat{F}`$ can depend on $`A`$ explicitly, the action must be gauge-invariant under the commutative gauge-invariance. This requires that the action be expressible in terms of $`F`$ alone.
2.2. Duality at $`O(\theta )`$
Since the action can be expressed in terms of $`F`$, we can implement a duality transformation in essentially the way described in the introduction. To perform the Legendre transform, we shift the action as before
$$SS+\frac{1}{2}A_DdF.$$
$`(2.8)`$
The equation of motion for $`F`$ gives,
$$g^2F_D=F+\frac{1}{2}\frac{\delta L_\theta }{\delta F}(F)+O(\theta ^2).$$
$`(2.9)`$
To lowest order in $`\theta `$, we can solve for $`F`$ in terms of $`F_D`$:
$$F=g^2F_D\frac{1}{2}\frac{\delta L_\theta }{\delta F}|_{F=g^2F_D}+O(\theta ^2).$$
$`(2.10)`$
At order $`\theta `$, loops play no role in the duality transformation so the quantum and semi-classical dual descriptions are equivalent. Plugging (2.10) into the action (2.5) gives,
$$S=\frac{g^2}{4}(F_DF_D+2\mathrm{tr}(\stackrel{~}{\theta }F_D^3)\frac{1}{2}\mathrm{tr}(\stackrel{~}{\theta }F_D)\mathrm{tr}(F_D^2))+O(\stackrel{~}{\theta }^2).$$
$`(2.11)`$
Note that we use $`\stackrel{~}{\theta }=g^2(\theta )`$ as the new non-commutativity parameter. The factor of $`g^2`$ in $`\stackrel{~}{\theta }`$ is natural because of the following scaling argument: we can schematically expand $`\widehat{F}^2`$,
$$\widehat{F}^2F^2\left(1+\underset{n,l}{}\theta ^{n+l}()^{2l}F^n\right),$$
$`(2.12)`$
on strictly dimensional grounds. This implies that iteratively, we can express $`F`$ in schematic form:
$$Fg^2F_D\left(1+\underset{n,l}{}\theta ^{n+l}()^{2l}(g^2F_D)^n\right).$$
$`(2.13)`$
In terms of $`\stackrel{~}{\theta }`$, we see that
$$Fg^2F_D(1+\underset{n,l}{}\stackrel{~}{\theta }^{n+l}()^{2l}\left(\frac{1}{g^2}\right)^l(F_D)^n).$$
$`(2.14)`$
The action now takes the form of a derivative expansion with higher derivatives of $`F_D`$ suppressed by powers of $`g^1`$.
There are a number of observations at this point. Substituting even the lowest order expression,
$$F=g^2F_D+O(\theta ),$$
$`(2.15)`$
into (2.5) results in an infinite number of terms involving higher powers of $`\stackrel{~}{\theta }`$. While terms beyond $`O(\stackrel{~}{\theta })`$ will receive additional corrections from the $`O(\theta )`$ corrections to (2.15), it seems quite clear โ barring miraculous cancellations โ that there is no upper bound on the power of $`\stackrel{~}{\theta }`$ that appears in the dual action. This suggests that it will be difficult to quantize the theory non-perturbatively in any conventional way. We also note that the dual action to leading order in $`\stackrel{~}{\theta }`$, expressed in dual non-commutative variables, takes the form
$$S=\frac{g^2}{4}\widehat{F}_D\widehat{F}_D+O(\stackrel{~}{\theta }^2).$$
$`(2.16)`$
As is natural, we define $`\widehat{F}_D`$ with respect to a star product involving $`\stackrel{~}{\theta }`$. However, it is quite possible that the corrections to (2.16) of $`O(\stackrel{~}{\theta }^2)`$ are non-vanishing. It is not clear that the resulting dual action would then have a purely quadratic form.
Acknowledgements
It is our pleasure to thank M. R. Douglas and N. Seiberg for helpful comments, and K. Dasgupta for early participation. The work of O.J.G. is supported in part by NSF Grant No. PHY-98-02484. The work of R.G. is supported in part by NSF grant No. DMS-9627351, while that of S.S. is supported by the William Keck Foundation and by NSF grant No. PHYโ9513835.
References
relax A. Connes, M. R. Douglas and A. Schwarz, hep-th/9711162, JHEP 9802:003, 1998. relax N. Seiberg and E. Witten, hep-th/9908142, JHEP 9909:032, 1999. relax C. Montonen and D. Olive, Phys. Lett. B72 (1977) 117. relax A. Sen, hep-th/9402032, Phys. Lett. B329 (1994) 217. relax N. Seiberg, L. Susskind and N. Toumbas, hep-th/0005015. relax M. Kreuzer and J.-G. Zhou, hep-th/9912174, JHEP 0001:011, 2000.
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# A generalization of Gordonโs theorem and applications to quasiperiodic Schrรถdinger operators
## 1. Introduction
In this paper we study one-dimensional Schrรถdinger operators of the form
(1)
$$H=\frac{d^2}{dx^2}+V(x),$$
acting on $`L^2()`$, with some real-valued $`L_{\mathrm{loc}}^1`$-potential $`V`$. We will be particularly interested in potentials of the form
(2)
$$V(x)=V_1(x)+V_2(x\alpha +\theta ),$$
where we assume that $`V_1`$ and $`V_2`$ are $`1`$-periodic and locally integrable, and $`\alpha ,\theta [0,1)`$. If $`\alpha =\frac{p}{q}`$ is rational, then the potential $`V`$ is $`q`$-periodic and $`H`$ has purely absolutely continuous spectrum. If $`\alpha `$ is irrational, then the potential is quasiperiodic and the spectral theory of $`H`$ is far from trivial; compare .
We want to study the eigenvalue problem for $`H`$. More precisely, we are interested in methods that allow one to exclude the presence of eigenvalues. A notion that has proved to be useful in this context is the following. A bounded potential $`V`$ on $`(\mathrm{},\mathrm{})`$ is called a Gordon potential if there exist $`T_m`$-periodic potentials $`V^{(m)}`$ such that $`T_m\mathrm{}`$ and for every $`m`$,
$$\underset{2T_mx2T_m}{sup}|V(x)V^{(m)}(x)|Cm^{T_m}$$
for some suitable constant $`C`$. It has been shown by Gordon (see also Simon ) that $`H`$ has no eigenvalues if $`V`$ is a Gordon potential. For discrete Schrรถdinger operators, certain variants of this result have been established by Delyon and Petritis and by Sรผtล ; see for a survey of the applications of criteria in this spirit. The applications in the discrete case include in particular results for models that are generated by discontinuous functions, for example, step functions. The interest in such models stems from the theory of one-dimensional quasicrystals; compare . It is clear that in the continuum case, these functions are outside the scope of Gordonโs result. This motivates our attempt to find a more general criterion for absence of eigenvalues.
Let us call $`V`$ a generalized Gordon potential if $`VL_{\mathrm{loc},\mathrm{unif}}^1()`$, that is,
$$V_{1,\mathrm{unif}}=\underset{x}{sup}_x^{x+1}|V(x)|๐x<\mathrm{}$$
and there exist $`T_m`$-periodic potentials $`V^{(m)}`$ such that $`T_m\mathrm{}`$ and for every $`C<\mathrm{}`$, we have
(3)
$$\underset{m\mathrm{}}{lim}\mathrm{exp}(CT_m)_{T_m}^{2T_m}|V(x)V^{(m)}(x)|๐x=0.$$
Clearly, every Gordon potential is a generalized Gordon potential. Our main result is the following:
###### Theorem 1.
Suppose $`V`$ is a generalized Gordon potential. Then the operator $`H`$ in (1) has empty point spectrum.
As in the classical case , the proof gives the stronger result that for every energy $`E`$, the solutions of
(4)
$$u^{\prime \prime }(x)+V(x)u(x)=Eu(x)$$
do not tend to zero as $`|x|\mathrm{}`$, that is, $`|u(x_n)|^2+|u^{}(x_n)|^2D`$ for some constant $`D>0`$ and a sequence $`(x_n)_n`$ which obeys $`|x_n|\mathrm{}`$ as $`n\mathrm{}`$. Thus there are no $`L^2`$-solutions since $`uL^2()`$ would imply $`|u(x)|^2+|u^{}(x)|^20`$ as $`|x|\mathrm{}`$ by Harnackโs inequality (see ). Note that this uses $`VL_{\mathrm{loc},\mathrm{unif}}^1()`$, which also guarantees that the operator $`H`$ can be defined by form methods or via Sturm-Liouville theory.
Let us now discuss the application of Theorem 1 to quasiperiodic $`V`$ given by (2). Given some irrational $`\alpha [0,1)`$, we consider its continued fraction expansion
$$\alpha =\frac{1}{a_1+{\displaystyle \frac{1}{a_2+{\displaystyle \frac{1}{a_3+\mathrm{}}}}}}$$
with uniquely determined $`a_m`$ and the continued fraction approximants $`\alpha _m=p_m/q_m`$ defined by
$`p_0`$ $`=0,`$ $`p_1`$ $`=1,`$ $`p_m`$ $`=a_mp_{m1}+p_{m2},`$
$`q_0`$ $`=1,`$ $`q_1`$ $`=a_1,`$ $`q_m`$ $`=a_mq_{m1}+q_{m2};`$
compare . Recall that $`\alpha `$ is called a Liouville number if
(5)
$$|\alpha \alpha _m|Bm^{q_m}$$
for some suitable $`B`$, and that the set of Liouville numbers is a dense $`G_\delta `$-set of zero Lebesgue measure. Given $`V`$ as in (2), we consider the $`q_m`$-periodic approximants $`V^{(m)}`$ defined by
(6)
$$V^{(m)}(x)=V_1(x)+V_2(x\alpha _m+\theta ).$$
We immediately obtain the following corollary to Theorem 1.
###### Corollary 1.1.
Suppose that for every $`C`$, we have
(7)
$$\underset{m\mathrm{}}{lim}\mathrm{exp}(Cq_m)_{q_m}^{2q_m}|V_2(x\alpha +\theta )V_2(x\alpha _m+\theta )|๐x=0.$$
Then $`V`$ (as given by (2)) is a generalized Gordon potential and $`H`$ (as given by (1)) has empty point spectrum.
Note that for $`\alpha ,\theta `$ fixed, the class of functions $`V_2`$ obeying (7) is a linear space, that is, it is closed under taking finite sums and under multiplication by constants. Moreover, we shall show that condition (7) is satisfied, for example, if $`V_2`$ is a Hรถlder continuous function, a step function, or a function with power-type singularities, and $`\alpha `$ is Liouville and $`\theta `$ arbitrary.
The organization of this paper is as follows. In Section 2 we establish estimates on solutions of (4) which will imply Theorem 1. The examples for condition (7) are discussed in Section 3.
## 2. Gronwall-Type Solution Estimates and Proof of Theorem 1
In this section we study the solutions to the eigenvalue equations associated to two potentials. These two potentials will later be given by a generalized Gordon potential and one of its approximants. We assume that the solutions have the same initial conditions at $`0`$. By an a priori estimate for the equivalent first order systems, found by a standard application of Gronwallโs lemma (e.g., ), we can bound the distance of the two solutions by an integral expression involving the distance of the potentials. It is this estimate which allows us to use $`L^1`$ rather than $`L^{\mathrm{}}`$-bounds in (3). Theorem 1 follows from this bound combined with some useful properties of solutions to periodic eigenvalue equations.
Fix two potentials $`W_1L_{\mathrm{loc},\mathrm{unif}}^1()`$, $`W_2L_{\mathrm{loc}}^1()`$ and some energy $`E`$ and consider the solutions $`u_1,u_2`$ of
$$u_1^{\prime \prime }(x)+W_1(x)u_1(x)=Eu_1(x),u_2^{\prime \prime }(x)+W_2(x)u_2(x)=Eu_2(x),$$
subject to
$$u_1(0)=u_2(0),u_1^{}(0)=u_2^{}(0),|u_1(0)|^2+|u_1^{}(0)|^2=|u_2(0)|^2+|u_2^{}(0)|^2=1.$$
###### Lemma 2.1.
There exists $`C=C(W_1E_{1,\mathrm{unif}})`$ such that for every $`x`$, we have
(8)
$$\left(\begin{array}{c}u_1(x)\\ u_1^{}(x)\end{array}\right)\left(\begin{array}{c}u_2(x)\\ u_2^{}(x)\end{array}\right)C\mathrm{exp}(C|x|)_{\mathrm{min}(0,x)}^{\mathrm{max}(0,x)}|W_1(t)W_2(t)||u_2(t)|๐t.$$
Proof. We consider the case $`x0`$ (the modifications for $`x<0`$ are obvious). We have
$`\left(\begin{array}{c}u_1(x)u_2(x)\\ u_1^{}(x)u_2^{}(x)\end{array}\right)=`$ $`{\displaystyle _0^x}\left(\begin{array}{c}u_1^{}(t)u_2^{}(t)\\ (W_1(t)E)u_1(t)(W_2(t)E)u_2(t)\end{array}\right)๐t`$
$`=`$ $`{\displaystyle _0^x}\left(\begin{array}{c}0\\ (W_1(t)W_2(t))u_2(t)\end{array}\right)๐t+`$
$`+{\displaystyle _0^x}\left(\begin{array}{c}u_1^{}(t)u_2^{}(t)\\ (W_1(t)E)(u_1(t)u_2(t))\end{array}\right)๐t`$
$`=`$ $`{\displaystyle _0^x}\left(\begin{array}{c}0\\ (W_1(t)W_2(t))u_2(t)\end{array}\right)๐t+`$
$`+{\displaystyle _0^x}\left(\begin{array}{cc}0& 1\\ W_1(t)E& 0\end{array}\right)\left(\begin{array}{c}u_1(t)u_2(t)\\ u_1^{}(t)u_2^{}(t)\end{array}\right)๐t.`$
Hence
$`\left(\begin{array}{c}u_1(x)u_2(x)\\ u_1^{}(x)u_2^{}(x)\end{array}\right)`$ $`{\displaystyle _0^x}|(W_1(t)W_2(t))||u_2(t)|๐t+`$
$`+{\displaystyle _0^x}\left(\begin{array}{cc}0& 1\\ W_1(t)E& 0\end{array}\right)\left(\begin{array}{c}u_1(t)u_2(t)\\ u_1^{}(t)u_2^{}(t)\end{array}\right)๐t.`$
By Gronwallโs lemma we therefore get
$`\left(\begin{array}{c}u_1(x)\\ u_1^{}(x)\end{array}\right)\left(\begin{array}{c}u_2(x)\\ u_2^{}(x)\end{array}\right)`$ $`{\displaystyle _0^x}|(W_1(t)W_2(t))||u_2(t)|dt\times `$
$`\times \mathrm{exp}\left({\displaystyle _0^x}\left(\begin{array}{cc}0& 1\\ W_1(t)E& 0\end{array}\right)๐t\right).`$
Choosing $`C`$ suitably, the assertion of the lemma follows. $`\mathrm{}`$
We see that we can control the difference of the solutions in terms of an integral condition involving the difference of the potentials. The other key ingredient in the proof of Theorem 1 is the fact that for periodic potentials, we have some knowledge about the norm of the solution vector $`(u(x),u^{}(x))^T`$ at certain points $`x`$. This is made explicit in the following lemma which is essentially well known (particularly in the discrete case ).
###### Lemma 2.2.
Suppose $`W`$ is $`p`$-periodic and $`E`$ is some arbitrary energy. Then every solution of
(9)
$$u^{\prime \prime }(x)+W(x)u(x)=Eu(x),$$
normalized in the sense that
(10)
$$|u(0)|^2+|u^{}(0)|^2=1,$$
obeys the estimate
$$\mathrm{max}(\left(\begin{array}{c}u(p)\\ u^{}(p)\end{array}\right),\left(\begin{array}{c}u(p)\\ u^{}(p)\end{array}\right),\left(\begin{array}{c}u(2p)\\ u^{}(2p)\end{array}\right))\frac{1}{2}.$$
Proof. This follows by the same reasoning as in the discrete case; compare . For the readerโs convenience, we sketch the argument briefly. Consider the solutions $`u`$ of (9). For $`x,y`$, $`x<y`$, the mapping
(11)
$$M(x,y):\left(\begin{array}{c}u(x)\\ u^{}(x)\end{array}\right)\left(\begin{array}{c}u(y)\\ u^{}(y)\end{array}\right)$$
is clearly linear and depends only on the energy $`E`$ and the potential on the interval $`(x,y)`$. Thus, since $`W`$ is $`p`$-periodic, we have
(12)
$$M(p,0)=M(0,p)=M(p,2p)=:M.$$
Moreover, by the Cayley-Hamilton theorem, we have
(13)
$$M^2\mathrm{tr}(M)M+I=0.$$
If $`|\mathrm{tr}(M)|1`$, we apply this equation to $`(u(0),u^{}(0))^T`$ obeying (10) and obtain, using (12),
$$\mathrm{max}(\left(\begin{array}{c}u(p)\\ u^{}(p)\end{array}\right),\left(\begin{array}{c}u(2p)\\ u^{}(2p)\end{array}\right))\frac{1}{2},$$
since $`(u(0),u^{}(0))^T`$ has norm one. If $`|\mathrm{tr}(M)|>1`$, we apply (13) along with (12) to $`(u(p),u^{}(p))^T`$ and obtain
$$\mathrm{max}(\left(\begin{array}{c}u(p)\\ u^{}(p)\end{array}\right),\left(\begin{array}{c}u(p)\\ u^{}(p)\end{array}\right))\frac{1}{2},$$
again since the vector $`(u(0),u^{}(0))^T`$ has norm one. Put together, we obtain the claimed result. $`\mathrm{}`$
We are now in a position to give the
Proof of Theorem 1. Let $`V`$ be a generalized Gordon potential and let $`V^{(m)}`$ be the $`T_m`$-periodic approximants obeying (3). Fix some $`m`$ and apply Lemma 2.1 with $`W_1=V`$ and $`W_2=V^{(m)}`$. We obtain
(14)
$$\left(\begin{array}{c}u(x)\\ u^{}(x)\end{array}\right)\left(\begin{array}{c}u_m(x)\\ u_m^{}(x)\end{array}\right)C_1\mathrm{exp}(C_1|x|)_{\mathrm{min}(0,x)}^{\mathrm{max}(0,x)}|V(t)V^{(m)}(t)||u_m(t)|๐t,$$
where $`u`$ (resp., $`u_m`$) solves $`u^{\prime \prime }(x)+V(x)u(x)=Eu(x)`$ (resp., $`u_m^{\prime \prime }(x)+V^{(m)}(x)u_m(x)=Eu_m(x)`$) and $`u,u_m`$ are both normalized at the origin and obey the same boundary condition there. We conclude from (3) that $`V^{(m)}_{1,\mathrm{unif}}`$ is bounded in $`m`$. Thus a second application of Lemma 2.1 with $`W_1=V^{(m)}`$ and $`W_2=0`$, noting that the constant in (8) only depends on the $`L_{\mathrm{loc},\mathrm{unif}}^1`$-norm of $`W_1E`$, leads to
$$\left(\begin{array}{c}u_m(x)\\ u_m^{}(x)\end{array}\right)\left(\begin{array}{c}u_0(x)\\ u_0^{}(x)\end{array}\right)C_2\mathrm{exp}(C_2|x|)_{\mathrm{min}(0,x)}^{\mathrm{max}(0,x)}|V^{(m)}(t)||u_0(t)|๐t,$$
where $`C_2`$ does not depend on $`m`$ and $`u_0`$ is a normalized solution of $`u_0^{\prime \prime }=Eu_0`$. Noting that $`u_0`$ is exponentially bounded, this gives
$$|u_m(x)|C_3\mathrm{exp}(C_3|x|)$$
with $`C_3`$ independent of $`m`$. This and (14) yield
$$\left(\begin{array}{c}u(x)\\ u^{}(x)\end{array}\right)\left(\begin{array}{c}u_m(x)\\ u_m^{}(x)\end{array}\right)C\mathrm{exp}(C|x|)_{\mathrm{min}(0,x)}^{\mathrm{max}(0,x)}|V(t)V^{(m)}(t)|๐t.$$
By (3) we find some $`m_0`$ such that for $`mm_0`$, we have
$$\left(\begin{array}{c}u(x)\\ u^{}(x)\end{array}\right)\left(\begin{array}{c}u_m(x)\\ u_m^{}(x)\end{array}\right)\frac{1}{4}$$
for every $`x`$ with $`T_mx2T_m`$. Combining this with Lemma 2.2, we can conclude the proof. $`\mathrm{}`$
## 3. Examples of Generalized Gordon Potentials
In this section we give examples of functions $`V_2`$ that obey condition (7) for Liouville frequencies $`\alpha `$ and hence induce quasiperiodic functions $`V`$ by (2) which are generalized Gordon potentials. These will include Hรถlder continuous functions, step functions, functions with local singularities, and linear combinations thereof.
Let us observe the following:
###### Proposition 3.1.
For fixed $`\alpha ,\theta `$, the class of functions $`V_2`$ obeying (7) is a linear space, that is, it is closed under taking finite sums and under multiplication by constants.
Proof. This is obvious. $`\mathrm{}`$
Define for some $`1`$-periodic function $`f`$,
$$\mathrm{osc}_{f,\epsilon }(x)=\underset{y,z(x\epsilon ,x+\epsilon )}{sup}|f(y)f(z)|.$$
Then we have the following proposition.
###### Proposition 3.2.
(a) If there are $`0<\delta ,D<\mathrm{}`$ such that
(15)
$$_0^1\mathrm{osc}_{V_2,\epsilon }(x)๐xD\epsilon ^\delta $$
for all sufficiently small $`\epsilon >0`$, then for every Liouville number $`\alpha [0,1)`$ and every $`\theta [0,1)`$, condition (7) is satisfied.
(b) Condition (15) holds for all Hรถlder continuous functions and for all step functions.
Proof. (a) Fix some $`C`$. Then by (5) and (15), we have
$`\underset{m\mathrm{}}{lim\; sup}\mathrm{exp}(Cq_m){\displaystyle _{q_m}^{2q_m}}|`$ $`V_2(x\alpha +\theta )V_2(x\alpha _m+\theta )|dx`$
$`\underset{m\mathrm{}}{lim\; sup}\mathrm{exp}(Cq_m){\displaystyle \frac{3q_m\alpha +1}{\alpha }}{\displaystyle _0^1}\mathrm{osc}_{V_2,2q_m|\alpha \alpha _m|}(x)๐x`$
$`\underset{m\mathrm{}}{lim\; sup}\mathrm{exp}(Cq_m){\displaystyle \frac{3q_m\alpha +1}{\alpha }}D(2q_m|\alpha \alpha _m|)^\delta `$
$`\underset{m\mathrm{}}{lim\; sup}\mathrm{exp}(Cq_m){\displaystyle \frac{3q_m\alpha +1}{\alpha }}D2^\delta q_m^\delta B^\delta m^{\delta q_m}`$
$`=0.`$
(b) This is straightforward. $`\mathrm{}`$
The class for which (7) was established in Proposition 3.2 contains only bounded potentials. We finally provide an example which shows that the use of generalized Gordon potentials allows one to exclude eigenvalues for some unbounded quasiperiodic potentials. We will exhibit some $`V_2`$ that has an integrable power-like singularity and which satisfies (7), and therefore $`H`$ defined by (1) and (2) has empty point spectrum. Note that by linearity this also gives examples with negative singularities and multiple singularities with different values for $`\gamma `$.
###### Proposition 3.3.
Let $`0<\gamma <1`$ and $`V_2(x)`$ be the $`1`$-periodic potential which for $`1/2x1/2`$ is given by $`V_2(x)=|x|^\gamma `$. Then for every Liouville number $`\alpha [0,1)`$ and every $`\theta [0,1)`$, condition (7) is satisfied.
Proof. For simplicity, we will only establish (7) for $`\theta =0`$. The calculations for general $`\theta `$ are similar but slightly more tedious. Start by writing
(16)
$$_{q_m}^{2q_m}|V_2(\alpha x)V_2(\alpha _mx)|๐x=\frac{q_m}{p_m}\underset{n=p_m}{\overset{2p_m}{}}_n^{n+1}\left|V_2\left(\frac{\alpha q_m}{p_m}y\right)V_2(y)\right|๐y$$
and
(17)
$$_n^{n+1}\left|V_2\left(\frac{\alpha q_m}{p_m}y\right)V_2(y)\right|๐y=_0^1\left|V_2\left(y+\left(\frac{\alpha q_m}{p_m}1\right)(y+n)\right)V_2(y)\right|๐y.$$
We have $`|\frac{\alpha q_m}{p_m}1||y+n|2p_m|\frac{\alpha q_m}{p_m}1|=:\delta <1/4`$ for $`m`$ sufficiently large and can estimate
(18) $`{\displaystyle _0^{1/2}}\left|V_2\left(y+\left({\displaystyle \frac{\alpha q_m}{p_m}}1\right)(y+n)\right)V_2(y)\right|๐y`$
$`C\delta +C\delta ^{1\gamma }+\left|{\displaystyle _0^{1/2}}\left(V_2\left(y+\left({\displaystyle \frac{\alpha q_m}{p_m}}1\right)(y+n)\right)V_2(y)\right)๐y\right|,`$
where the $`\delta ^{1\gamma }`$ term arises from the singularity of $`V_2`$ at $`0`$, and the monotonicity of $`V_2`$ in $`[0,1/2]`$ was used to take the absolute value outside the integral. The integral on the right can be calculated explicitly, which eventually leads to an estimate $`C(p_m|\frac{\alpha q_m}{p_m}1|)^{1\gamma }`$ for its absolute value and thus also for (18).
In a similar way we get the same estimate for the integral from $`1/2`$ to $`1`$ on the right hand side of (17). Inserting into (16) we finally find
$$_{q_m}^{2q_m}\left|V_2(\alpha x)V_2(\alpha _mx)\right|๐xCp_m^{2\gamma }\left|\frac{\alpha q_m}{p_m}1\right|^{1\gamma }.$$
In view of (5) this suffices to imply (7). $`\mathrm{}`$
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# Resolution of singularities of null cones1991 Mathematics Subject Classification. Primary 14L35, 22G.
## Introduction
The purpose of this paper is to present canonical resolutions of singularities of certain cone varieties arising naturally in invariant theory. Examples of such cone varieties are provided by the so-called null cones which are the zero locus of the $`O_n`$ (or $`Sp_n`$) invariant polynomials on $`M_{n,m}`$, where $`M_{n,m}`$ denotes the space of $`n\times m`$ complex matrices (cf. e.g. \[H\]). Another example is the null cone for general linear groups. Our construction resembles the celebrated Springer resolution of the nilpotent variety in a complex semisimple Lie algebra which plays a vital role in representation theory and geometry (cf. \[CG\]).
It turns out that the cone varieties studied here include as special cases the variety $`Z`$ constructed in \[KP1, KP2\] associated to two-column partitions. The Kraft-Procesi variety $`Z`$ is a complete intersection whose quotient under certain group coincide with the closure of a nilpotent conjugacy class in classical Lie algebras. It is a very interesting problem to construct a canonical resolution of singularities of $`Z`$ in general. We establish relations between our resolutions of singularities and those for closure of conjugacy classes associated to two-column partitions. In the general linear Lie algebra case, we find more than one canonical resolution of singularities of null cones and of closure of conjugacy classes.
We mention in passing that the present work grows out of an attempt to generalize the geometric construction of $`(gl_n,gl_m)`$-duality \[W\] in the spirit of \[CG\] to other Howe duality \[H\].
The paper is divided into two sections. In Sect. 1 we present resolutions of singularities of null cones in the $`O_n`$ and $`Sp_n`$ cases, and establish connections with resolutions of singularities of closure of conjugacy classes. In Sect. 2 we treat the analog of Sect. 1 for the general linear Lie algebras.
## 1 Null cones associated to $`O_n`$ and $`Sp_n`$
Denote by $`GL_n,O_n`$, and $`Sp_{2n}`$ the complex general linear group, orthogonal group, and the symplectic group respectively. Denote by $`M_{n,m}`$ the set of all complex $`n\times m`$ matrices.
### 1.1 The orthogonal setup
Let $`V`$ be a vector space of complex dimension $`n`$ with a non-degenerate symmetric bilinear form. We identify $`V`$ with $`C^n`$ endowed with the standard symmetric bilinear form
$$(u,v)=\underset{i=1}{\overset{n}{}}u_iv_i,$$
where $`u,v`$ denote the $`n`$-tuples $`(u_1,\mathrm{},u_n)`$ and $`(v_1,\mathrm{},v_n)`$.
For convenience, we can identify $`M_{n,m}`$ with the direct sum $`(C^n)^m`$ of $`m`$ copies of $`C^n`$. We may write a typical element of $`(C^n)^m`$ as an $`m`$-tuple:
$$A=(v_1,v_2,\mathrm{},v_m),v_jC^n.$$
We define a set of quadratic polynomials $`\stackrel{~}{\xi }_{ij}=\stackrel{~}{\xi }_{ji},1i,jm`$ on $`(C^n)^m`$ by
$$\stackrel{~}{\xi }_{ij}(v)=(v_i,v_j)=\underset{a=1}{\overset{n}{}}v_{ai}v_{aj}.$$
###### Remark 1
The First Fundamental Theorem for $`O_n`$ (cf. \[H\]) says that the polynomials $`\stackrel{~}{\xi }_{ij}=\stackrel{~}{\xi }_{ji},1i,jm,`$ generate the algebra of $`O_n`$ invariant polynomials on $`(C^n)^m`$.
We identify the space $`S^2(C^m)`$ of second symmetric tensors as the space of symmetric $`m\times m`$ matrices. Define a map
$`Q:M_{n,m}`$ $``$ $`S^2(C^m)`$
$`T`$ $``$ $`T^tT.`$
Here $`T^t`$ denotes the transpose of $`T`$. By identifying $`M_{n,m}`$ with $`(C^n)^m`$, we can equivalently define $`Q(T)`$ as the $`m\times m`$ symmetric matrix whose $`(i,j)`$-th entry is equal to $`\stackrel{~}{\xi }_{ij}`$.
### 1.2 The symplectic setup
Let $`V`$ be a vector space of complex dimension $`n`$ with a non-degenerate symplectic (i.e. anti-symmetric) bilinear form. It is well known that $`n`$ has to be an even number, say $`2p`$. We identify $`V`$ with $`C^n`$ endowed with the following standard symplectic bilinear form
$$(v,v^{})=\underset{j=1}{\overset{n/2}{}}(x_jy_j^{}y_jx_j^{}).$$
Here $`v=(x,y)`$ and $`v^{}=(x^{},y^{})`$ are elements of $`C^{2p}`$, expressed as $`2`$-tuples of elements of $`C^p`$.
We identify $`M_{n,m}`$ with the direct sum $`(C^n)^m`$ of $`m`$ copies of $`C^n`$ as before. We may write a typical element of $`(C^n)^m`$ as an $`m`$-tuple:
$$A=(v_1,v_2,\mathrm{},v_m),v_jC^n.$$
We define a set of quadratic polynomials $`\stackrel{ห}{\xi }_{ij}=\stackrel{ห}{\xi }_{ji},1i,jm`$ on $`(C^n)^m`$ by
$$\stackrel{~}{\xi }_{ij}(v)=(v_i,v_j).$$
###### Remark 2
The First Fundamental Theorem for $`Sp_n`$ (cf. \[H\]) says that the polynomials $`\stackrel{ห}{\xi }_{ij}=\stackrel{ห}{\xi }_{ji},1i,jm,`$ generate the algebra of $`Sp_n`$ invariant polynomials on $`(C^n)^m`$.
We identify the space $`\mathrm{\Lambda }^2(C^m)`$ of second anti-symmetric tensors in $`C^m`$ as the space of anti-symmetric $`m\times m`$ matrices. Denote by
$$J=\left[\begin{array}{cc}\hfill 0& \hfill I_p\\ \hfill I_p& \hfill 0\end{array}\right],$$
where $`I_p`$ denotes the identity $`p\times p`$ matrix. Define a map
$`Q_{sp}:M_{n,m}`$ $``$ $`\mathrm{\Lambda }^2(C^m)`$
$`T`$ $``$ $`T^tJT.`$
We also write $`Q_{sp}`$ as $`Q`$ when no ambiguity arises. By identifying $`M_{n,m}`$ with $`(C^n)^m`$, we can equivalently define $`Q(T)`$ as the $`m\times m`$ anti-symmetric matrix whose $`(i,j)`$-th entry equal to $`\stackrel{ห}{\xi }_{ij}`$.
### 1.3 A resolution of singularities of $`๐ฉ๐๐ฌ`$
Let $`V`$ be a vector space of complex dimension $`n`$ with a nondegenerate bilinear form, either symmetric or anti-symmetric. Let $`G(V)`$ be the isometry group of the form, which is $`O_n`$ in symmetric case and $`Sp_n`$ in anti-symmetric case. Let $`g(V)`$ be its Lie algebra.
Given $`A`$ in $`S^2(C^m)`$ (resp. $`\mathrm{\Lambda }^2(C^m)`$), its inverse image $`Q^1(A)`$ under the map $`Q`$ is referred to as the fiber of $`Q`$ over $`A`$. Of particular interest is the fiber over $`0`$ (cf. \[H\]), which we will refer to as the null cone (since it is obviously preserved by scalar dilations) and denote by $`๐ฉ๐๐ฌ`$. Equivalently, the variety $`๐ฉ๐๐ฌ`$ is the set of $`n\times m`$ matrices on which all $`G(V)`$ invariant polynomials on $`M_{n,m}`$ take value $`0`$. Often we will think of $`M_{n,m}`$ as the space $`\text{Hom}(C^m,V)`$ of all linear maps from $`C^m`$ to $`V`$ (or $`C^n`$). An element in $`๐ฉ๐๐ฌ`$ is called a null mapping. We easily have
$`Q(gT)`$ $`=`$ $`Q(T),TM_{n,m},gG(V),`$ (1)
$`Q(Th)`$ $`=`$ $`h^tQ(T)h,hGL_m.`$ (2)
###### Remark 3
It follows from Eqs. (1) and (2) that $`G(V)`$ acts on $`๐ฉ๐๐ฌ`$ and this action commutes with the action of $`GL_m`$. The space of regular functions on $`๐ฉ๐๐ฌ`$, under the induced action of $`G(V)\times GL_m`$, has a beautiful decomposition into $`G(V)\times GL_m`$-modules (cf. \[H\]).
A subspace of $`V`$ is called isotropic if any two vectors in the subspace are orthogonal to each other with respect to the corresponding bilinear form. We observe that $`TM_{n,m}=\text{Hom}(C^m,V)`$ is a null mapping if and only if the image of $`T`$, denoted by $`\mathrm{}T`$, form an isotropic subspace of $`C^n`$. Denote by $`J_r(V)`$ the set of all $`r`$-dimensional isotropic subspaces of $`V`$. The set $`J_r(V)`$ is nonempty if and only if $`rn/2`$. It is well known that $`G(V)`$ acts on $`J_r(V)`$ transitively (cf. \[H\], Appendix 3). The dimension of $`J_r(V)`$ $`(rn/2)`$ can be calculated to be $`r(2n3r1)/2`$.
Note that the null cone $`๐ฉ๐๐ฌ`$ is a singular variety defined in terms of a finite set of quadratic equations. The first goal of this paper is to present a canonical resolution of singularities of $`๐ฉ๐๐ฌ`$. Our construction is reminiscent of the Springer resolution of singularities of the nilpotent cone in a complex semisimple Lie algebra (cf. e.g. \[CG\]).
Set $`r=\mathrm{min}(m,[n/2])`$ from now on, where $`[n/2]`$ denotes the integer closest to and no larger than $`n/2`$. We introduce the following variety
$$\stackrel{~}{๐ฉ๐๐ฌ}=\{(T,U)๐ฉ๐๐ฌ\times J_r(V)\mathrm{}TU\}.$$
We have the following projection maps:
$$\begin{array}{ccc}& \stackrel{~}{๐ฉ๐๐ฌ}& \\ \hfill \mu _0& & \pi _0\hfill \\ \hfill ๐ฉ๐๐ฌ& & J_r(V)\hfill \end{array}$$
The diagram is $`O_n\times GL_m`$-equivariant, where $`GL_m`$ acts trivially on $`\stackrel{~}{๐ฉ๐๐ฌ}`$ and $`J_r(V)`$. $`\mu _0`$ is proper since $`\stackrel{~}{๐ฉ๐๐ฌ}๐ฉ๐๐ฌ\times J_r(V)`$ and $`J_r(V)`$ is compact. It is easy to see that $`\pi _0`$ and $`\mu _0`$ are both surjective.
Denote by $`T_J`$ the tautological bundle over $`J_r(V)`$:
$$T_J=\{(u,U)V\times J_r(V)uU\}.$$
Denote by $`\underset{ยฏ}{C^r}`$ the rank $`r`$ trivial bundle over $`J_r(V)`$. Given $`UJ_r(V)`$, the fiber of $`\pi _0`$ over $`U`$ is canonically identified with $`\text{Hom}(C^m,U)`$. Thus we have
$`dim\stackrel{~}{๐ฉ๐๐ฌ}`$ $`=`$ $`mr+dimJ_r(V)`$
$`=`$ $`mr+r(2n3r1)/2`$
$`=`$ $`r(2m+2n3r1)/2.`$
This proves the following characterization of $`\stackrel{~}{๐ฉ๐๐ฌ}`$.
###### Proposition 1
The variety $`\stackrel{~}{๐ฉ๐๐ฌ}`$ is isomorphic to the tensor product $`Taut\underset{ยฏ}{C}^m`$ of vector bundles $`Taut`$ and $`\underset{ยฏ}{C}^m`$ over the variety $`J_r(V)`$ such that the following diagram commutes:
$$\begin{array}{ccc}\hfill \stackrel{~}{๐ฉ๐๐ฌ}& & T_J\underset{ยฏ}{C}^m\hfill \\ \hfill \pi _0& & \hfill \\ & J_r(V)& \end{array}$$
In particular, $`\stackrel{~}{๐ฉ๐๐ฌ}`$ is a smooth variety of dimension $`r(2m+2n3r1)/2`$, where $`r=\mathrm{min}(m,[n/2])`$.
###### Theorem 1
The map $`\mu _0:\stackrel{~}{๐ฉ๐๐ฌ}๐ฉ๐๐ฌ`$ is a resolution of singularities.
Proof. It is clear that the set $`๐ฉ๐๐ฌ_0`$ of all null mappings of maximal rank which is equal to $`r=\mathrm{min}(m,[n/2])`$ is a Zariski-open subvariety of $`๐ฉ๐๐ฌ`$. Given $`T๐ฉ๐๐ฌ_0`$ there exists a unique $`UJ_r(V)`$ containing $`\mathrm{}T`$, namely $`\mathrm{}T`$ itself. Thus the map $`\mu _0:\stackrel{~}{๐ฉ๐๐ฌ}๐ฉ๐๐ฌ`$ over an open set $`\mu _0^1(๐ฉ๐๐ฌ_0)`$ is one-to-one. Together with the smoothness of $`\stackrel{~}{๐ฉ๐๐ฌ}`$ provided by Proposition 1, we conclude the proof. $`\mathrm{}`$
###### Remark 4
In the case $`n2m`$ and so $`r=m`$, we easily see that the map $`Q`$ maps surjectively to the space $`Sym_m`$ of $`m\times m`$ symmetric matrices, and it is submersive at any point $`T`$ in $`๐ฉ๐๐ฌ_0`$. It follows that
$`dim๐ฉ๐๐ฌ`$ $`=`$ $`dimM_{n,m}dimSym_m`$
$`=`$ $`nm(m^2+m)/2=dim\stackrel{~}{๐ฉ๐๐ฌ}.`$
This of course was also implied by Theorem 1.
It is known that $`J_r(V)`$ is disconnected if and only if we are in the orthogonal case and $`n=2m`$ (which implies $`r=m`$) (cf. \[H\], Appendix 2). Let us assume that we are in such a case first of all, and so $`J_r(V)`$ has two smooth connected components (cf. \[H\]). It follows by Proposition 1 that $`\stackrel{~}{๐ฉ๐๐ฌ}`$ also has two smooth connected components, denoted by $`\stackrel{~}{๐ฉ๐๐ฌ}^+`$ and $`\stackrel{~}{๐ฉ๐๐ฌ}^{}`$ respectively. The null cone $`๐ฉ๐๐ฌ`$ also consists of two irreducible components $`๐ฉ๐๐ฌ^\pm `$, which are the image of the two connected components of $`\stackrel{~}{๐ฉ๐๐ฌ}`$ respectively.
Now assume we are in the symplectic case, or in the orthogonal case but $`n2m`$. Then $`J_r(V)`$ is connected smooth and so is $`\stackrel{~}{๐ฉ๐๐ฌ}`$. Thus $`๐ฉ๐๐ฌ`$ is irreducible as the image of the surjective map $`\mu _0`$ of the irreducible variety $`\stackrel{~}{๐ฉ๐๐ฌ}`$.
### 1.4 Relations with closure of conjugacy classes
In this subsection,we always assume that $`n2m`$, and in addition $`m`$ is even in the orthogonal case. We need to reformulate the definition of $`๐ฉ๐๐ฌ`$.
Let $`W`$ be a vector space of complex dimension $`m`$ with a non-degenerate bilinear form of type opposite to the one on $`V`$. We identify $`\text{Hom}(W,V)=M_{nm}`$. Given any $`T\text{Hom}(W,V)`$ the adjoint $`T^{}`$ is defined by
$$(Tw,v)_V=(w,T^{}v)_W,wW,vV.$$
Here we use the subscripts to indicate to which vector space a bilinear form belong.
###### Remark 5
In the setup of Subsect. 1.1 and 1.2 and $`W=C^m`$, we easily check by definition that if the bilinear form on $`V`$ is anti-symmetric then $`T^{}=JT^t`$; if the bilinear form on $`V`$ is symmetric, then $`T^{}=T^tJ`$.
Consider the diagram
$$\begin{array}{ccc}\hfill \text{Hom}(W,V)& \stackrel{\stackrel{~}{Q}}{}& g(W)\hfill \\ \hfill R& & \\ \hfill g(V)& & \end{array}$$
where $`\stackrel{~}{Q}`$ is given by $`TT^{}T`$ and $`R`$ by $`TTT^{}`$. It is easy to check that the images of $`\stackrel{~}{Q}`$ and $`R`$ lie in $`g(W)`$ and $`g(V)`$ respectively. It follows from definitions and Remark 5 that $`\stackrel{~}{Q}^1(0)=Q^1(0)๐ฉ๐๐ฌ`$.
Denote by $`๐_\lambda `$ the conjugacy class of the group $`G(V)`$ associated to the partition $`\lambda =(2^m,1^{n2m})`$, which is the intersection of $`g(V)`$ with the conjugacy class of $`gl(V)`$ associated to $`\lambda `$. The closure $`\overline{๐}_\lambda `$ is indeed the variety of endomorphisms in $`g(V)`$ of rank no greater than $`m`$ and whose image is an isotropic subpace of $`V`$.
The variety $`๐ฉ๐๐ฌ=\stackrel{~}{Q}^1(0)`$ appears as a special case of the variety $`Z`$ in \[KP2\]. Recall that a quotient of a $`G`$-variety $`M`$ by the group $`G`$ is by definition the spectrum of the algebra of $`G`$-invariant regular functions on $`M`$. A special case of a theorem of Kraft and Procesi relevant to our considerations can be formulated as follows.
###### Theorem 2
The map $`R`$ maps $`๐ฉ๐๐ฌ`$ surjectively to the closure $`\overline{๐}_\lambda `$ of the conjugacy class $`๐_\lambda `$ associated to the partition $`\lambda =(2^m,1^{n2m})`$. Furthermore $`R`$ can be identified with the quotient map by $`G(W)`$ from $`๐ฉ๐๐ฌ`$ to $`\overline{๐}_\lambda `$.
Define the variety
$`\stackrel{~}{๐}_\lambda =\{(g,U)\overline{๐}_\lambda \times J_m(V)|\mathrm{}gU\},`$
and denote by $`p_0`$ the (surjective) projection from $`\stackrel{~}{๐}_\lambda `$ to the first factor $`\overline{๐}_\lambda `$. We can identify $`\stackrel{~}{๐}_\lambda `$ with a vector bundle over $`J_m(V)`$, whose fiber over $`UJ_m(V)`$ is the vector space
$$F_0=\{gg(V)|\mathrm{}gU\}.$$
We remark that for $`gF_0`$ we have $`U\mathrm{ker}g`$ and thus $`g^2=0gl(V)`$. Since $`p_0:\stackrel{~}{๐}_\lambda \overline{๐}_\lambda `$ is one-to-one over the open subset $`p_0^1(๐_\lambda )`$ consisting of the endomorphisms of maximal rank $`m`$, we have proved that
###### Proposition 2
The map $`p_0:\stackrel{~}{๐}_\lambda ๐_\lambda `$ is a resolution of singularities.
The quotient map $`R`$ by $`G(W)`$ induces a natural quotient map by $`G(W)`$ from $`\stackrel{~}{๐ฉ๐๐ฌ}`$ to $`\stackrel{~}{๐}_\lambda `$. We have the following commutative diagram:
$$\begin{array}{ccc}\stackrel{~}{๐ฉ๐๐ฌ}\hfill & \stackrel{\mu _0}{}& ๐ฉ๐๐ฌ\hfill \\ /G(W)\hfill & & /G(W)\hfill \\ \stackrel{~}{๐}_\lambda \hfill & \stackrel{p_0}{}& \overline{๐}_\lambda \hfill \end{array}$$
###### Remark 6
In the case $`n=2m`$ and $`G(V)=O_n`$, the $`O_n`$-conjugacy class $`๐_\lambda `$ associated to $`\lambda =(2^m)`$ is very even, and so $`\stackrel{~}{๐}_\lambda `$ splits into two $`SO_n`$-conjugacy class $`\stackrel{~}{๐}_\lambda ^\pm `$ (cf. \[KP2\]). Then the above diagram can be refined as follow:
$$\begin{array}{ccc}\hfill \stackrel{~}{๐ฉ๐๐ฌ}^\pm & \stackrel{\mu _0}{}& ๐ฉ๐๐ฌ^\pm \hfill \\ \hfill & & \hfill \\ \hfill \stackrel{~}{๐}_\lambda ^\pm & \stackrel{p_0}{}& ๐_\lambda ^\pm \hfill \end{array}$$
where $`\stackrel{~}{๐}_\lambda ^\pm `$ is defined by
$$\stackrel{~}{๐}_\lambda =\{(g,U)\overline{๐}_\lambda \times J_m^\pm (V)|\mathrm{}gU\},$$
and $`J_m^\pm (V)`$ are the two connected components of $`J_m(V)`$.
## 2 Null cones associated to $`GL_n`$
Let $`V`$ be a vector space of complex dimension $`n`$. Define a map
$$\phi :\text{Hom}(V,C^s)\times \text{Hom}(C^m,V)\text{Hom}(C^m,C^s)=M_{s,m}$$
by $`\phi (A,B)=AB=(\xi _{ij})_{1is,1jm}`$. Hence $`\xi _{ij}`$ defines a polynomial function on $`\text{Hom}(V,C^s)\times \text{Hom}(C^m,V)V^m(V^{})^s`$. Denote by $`๐ฉ=\phi ^1(0).`$ The group $`GL_s\times GL(V)\times GL_m`$ acts on $`๐ฉ`$ by
$$(g_1,g_2,g_3).(A,B)=(g_1Ag_2^1,g_2Bg_3^1).$$
###### Remark 7
The First Fundamental Theorem for $`GL(V)`$ says that the $`\xi _{ij},1is,1jm`$ generate the algebra of $`GL(V)`$ invariant polynomials on $`V^m(V^{})^s`$.
From now on we restrict ourselves to the case $`ns+m`$.
### 2.1 A resolution of singularities of $`๐ฉ`$
We introduce a variety
$`\stackrel{~}{๐ฉ}`$ $`=`$ $`\{(A,B,U_1,U_2)๐ฉ\times (m,ns,V)\mathrm{}BU_1U_2\mathrm{ker}A\}`$
$``$ $`๐ฉ\times (m,ns,V),`$
where $`\mathrm{ker}A`$ denotes the kernel of $`A`$, and $`(m,ns,V)`$ denotes the generalized flag variety
$$(m,ns,V)=\{(U_1,U_2)U_1U_2V,dimU_1=m,dimU_2=ns\}.$$
The group $`GL_s\times GL(V)\times GL_m`$ acts on $`\stackrel{~}{๐ฉ}`$ by letting
$$(g_1,g_2,g_3).(A,B,U_1,U_2)=(g_1Ag_2^1,g_2Bg_3^1,g_2U_1,g_2U_2).$$
We have the following projection maps:
$$\begin{array}{ccc}& \stackrel{~}{๐ฉ}& \\ \mu \hfill & & \pi \hfill \\ ๐ฉ\hfill & & (m,ns,V)\hfill \end{array}$$
The diagram above is $`GL_s\times GL(V)\times GL_m`$-equivariant, where $`G(V)`$ acts naturally on the generalized flag variety $`(m,ns,V)`$ while $`GL_s`$ and $`GL_m`$ act on $`(m,ns,V)`$ trivially. It is easy to see that $`\pi `$ and $`\mu `$ are surjective thanks to the assumption $`ns+m`$, and that $`\mu `$ is a proper map due to the compactness of $`(m,ns,V)`$. The fiber of $`\pi `$ over a point $`(U_1,U_2)(m,ns,V)`$ can be identified with the vector space
$$F=\text{Hom}(V/U_2,C^s)\times \text{Hom}(C^m,U_1).$$
(3)
In other words, $`\stackrel{~}{๐ฉ}`$ can be identified with a vector bundle $`๐ฆ`$ over the generalized flag variety $`(m,ns,V)`$:
$$๐ฆ:=\underset{ยฏ}{C^m}\text{Taut}\underset{ยฏ}{C^s}๐ฌ^{},$$
where $`\underset{ยฏ}{C^s},\underset{ยฏ}{C^m}`$ are trivial bundles of rank $`s`$ and $`m`$ respectively, and Taut, $`๐ฌ^{}`$ are respectively the tautological bundle and the dual quotient bundle:
Taut $`=`$ $`\{(u,U_1,U_2)vU_1\}V\times (m,ns,V)`$
$`๐ฌ^{}`$ $`=`$ $`\{(v,U_1,U_2)v(V/U_2)^{}\}V^{}\times (m,ns,V).`$
Summarizing, we have proved that
###### Proposition 3
We have the following isomorphism of vector bundles over the generalized flag variety $`(m,ns,V)`$:
$$\begin{array}{ccc}\stackrel{~}{๐ฉ}\hfill & & ๐ฆ\hfill \\ \pi \hfill & & \hfill \\ & (m,ns,V)& \end{array}$$
In particular, $`\stackrel{~}{๐ฉ}`$ is a smooth variety of dimension $`sn+mnsm`$.
Proof. It remains to calculate that
$`dim\stackrel{~}{๐ฉ}`$ $`=`$ $`dim(m,ns,V)+dimF`$
$`=`$ $`{\displaystyle \frac{1}{2}}(n^2(s^2+(nsm)^2+m^2)+(s^2+m^2)`$
$`=`$ $`sn+mnsm,`$
where the equality $`dim(m,ns,V)=\frac{1}{2}(n^2(s^2+(nsm)^2+m^2)`$ follows readily from the description of the generalized flag variety $`(m,ns,V)`$ in terms of the quotient of $`GL(V)`$ by an appropriate parabolic subgroup. $`\mathrm{}`$
###### Theorem 3
$`\mu :\stackrel{~}{๐ฉ}๐ฉ`$ is a resolution of singularities.
Proof. The subset $`๐ฉ_0`$ of $`๐ฉ`$ consisting of pairs of maximal rank matrices $`(A,B)`$ is a Zariski-open set in $`๐ฉ`$. A pair $`(A,B)๐ฉ_0`$ uniquely determines a pair $`(U_1,U_2)(m,ns,V)`$ such that $`(A,B,U_1,U_2)\stackrel{~}{๐ฉ}`$, namely $`U_1=\mathrm{}B,U_2=\mathrm{ker}A`$. This shows that $`\mu `$ maps $`\mu ^1(๐ฉ_0)\stackrel{~}{๐ฉ}`$ bijectively to $`๐ฉ_0`$. Together with the smoothness of $`\stackrel{~}{๐ฉ}`$ proved in Proposition 3, we have concluded the proof. $`\mathrm{}`$
###### Remark 8
Indeed it is easy to see that $`๐ฉ_0`$ is a single $`GL(V)`$-orbit since $`ns+m`$. An easy calculation shows that $`\phi `$ is submersive over any point in $`๐ฉ_0`$. So we have
$`dim๐ฉ`$ $`=`$ $`dim\text{Hom}(V,C^s)+dim\text{Hom}(C^m,V)dimM_{s,m}`$
$`=`$ $`sn+mnsm=dim\stackrel{~}{๐ฉ}`$
by Proposition 3. This fact, of course was implied by Theorem 3.
When $`s=m`$ and thus $`n2m`$, the space $`๐ฉ`$ appears as a special case of the variety $`Z`$ studied in \[KP1\]. It is shown \[KP1\] that the quotient of $`๐ฉ`$ by the diagonal action of $`GL_m`$ is the variety of $`n\times n`$ matrices of square $`0`$ and rank at most $`m`$, which is the closure $`\overline{๐ช}_\lambda `$ of the conjugacy class $`๐ช_\lambda gl(V)`$ corresponding to the partition $`\lambda =(2^m,1^{n2m})`$. More explicitly the quotient map from $`๐ฉ`$ to $`\overline{๐ช}_\lambda `$ is given by $`(A,B)BA`$.
Denote by
$$\stackrel{~}{๐ช}_\lambda =\{(g,U_1,U_2)\mathrm{}gU_1U_2\mathrm{ker}g\}\overline{๐ช}_\lambda \times (m,nm,V),$$
and by $`p`$ the natural surjective projection from $`\stackrel{~}{๐ช}_\lambda `$ to $`\overline{๐ช}_\lambda `$. By a similar argument which leads to Proposition 3, we have the following identification
$$\begin{array}{ccc}\stackrel{~}{๐ช}_\lambda \hfill & & TautQ^{}\hfill \\ \hfill & & \hfill \\ & (m,nm,V)& \end{array}$$
which implies that $`\stackrel{~}{๐ช}_\lambda `$ is smooth. Since the natural surjective projection $`p`$ from $`\stackrel{~}{๐ช}_\lambda `$ to (the first factor) $`\overline{๐ช}_\lambda `$ is one-to-one over the open set $`p^1(๐ช_\lambda )`$, it is a resolution of singularities. Thus we have established
###### Proposition 4
The map $`p:\stackrel{~}{๐ช}_\lambda \overline{๐ช}_\lambda `$ is a resolution of singularities.
###### Remark 9
This resolution $`p:\stackrel{~}{๐ช}_\lambda \overline{๐ช}_\lambda `$ differs from the classical one in terms of cotangent bundle of a grassmannian (compare with Proposition 6).
The quotient map from $`๐ฉ`$ to $`\overline{๐ช}_\lambda `$ induces a quotient map by $`GL_m`$ from $`\stackrel{~}{๐ฉ}`$ to $`\stackrel{~}{๐ช}_\lambda `$ which makes the following digram
$`\begin{array}{ccc}\stackrel{~}{๐ฉ}\hfill & \stackrel{\mu }{}& ๐ฉ\hfill \\ /GL_m\hfill & & /GL_m\hfill \\ \stackrel{~}{๐ช}_\lambda \hfill & \stackrel{p}{}& \overline{๐ช}_\lambda \hfill \end{array}`$ (7)
commutative.
### 2.2 A second resolution of singularities of $`๐ฉ`$
We introduce a variety
$`\stackrel{~}{๐ฉ}_1`$ $`=`$ $`\{(A,B,U)๐ฉ\times Gr(m,V)\mathrm{}BU\mathrm{ker}A\}`$
$``$ $`๐ฉ\times Gr(m,V),`$
where $`Gr(m,V)`$ denotes the grassmannian of $`m`$-dimensional subspaces of $`V`$. The group $`GL_s\times GL(V)\times GL_m`$ acts on $`\stackrel{~}{๐ฉ}_1`$ by letting
$$(g_1,g_2,g_3).(A,B,U)=(g_1Ag_2^1,g_2Bg_3^1,g_2U).$$
We have the following projection maps:
$$\begin{array}{ccc}& \stackrel{~}{๐ฉ}_1& \\ \mu _1\hfill & & \pi _1\hfill \\ ๐ฉ\hfill & & Gr(m,V)\hfill \end{array}$$
The diagram above is $`GL_s\times GL(V)\times GL_t`$-equivariant, where $`G(V)`$ acts naturally on $`Gr(m,V)`$ while $`GL_s,GL_t`$ act trivially on $`Gr(m,V)`$. It is easy to see that $`\pi _1`$ and $`\mu _1`$ are surjective and that $`\mu _1`$ is a proper map. The fiber of $`\pi _1`$ over a point $`UGr(m,V)`$ can be identified with the vector space
$$F_1=\text{Hom}(V/U,C^s)\times \text{Hom}(C^m,U).$$
In other words, $`\stackrel{~}{๐ฉ}_1`$ can be identified with a vector bundle $`๐ฆ_1`$ over the $`Gr(m,V)`$:
$$๐ฆ_1:=\underset{ยฏ}{C^m}\text{Taut}_1\underset{ยฏ}{C^s}๐ฌ_1^{},$$
where $`\underset{ยฏ}{C^s},\underset{ยฏ}{C^m}`$ are trivial bundles of rank $`s`$ and $`t`$ respectively, and $`\text{Taut}_1`$, $`๐ฌ_1^{}`$ are respectively the following tautological bundle and the dual quotient bundle:
$`\text{Taut}_1`$ $`=`$ $`\{(u,U)vU\}V\times Gr(m,V)`$
$`๐ฌ_1^{}`$ $`=`$ $`\{(v,U)v(V/U)^{}\}V^{}\times Gr(m,V).`$
Summarizing, we have proved that
###### Proposition 5
We have the following isomorphism of vector bundles over the generalized flag variety $`Gr(m,V)`$:
$$\begin{array}{ccc}\stackrel{~}{๐ฉ}_1\hfill & & ๐ฆ_1\hfill \\ \pi _1\hfill & & \hfill \\ & Gr(m,V)& \end{array}$$
In particular, $`\stackrel{~}{๐ฉ}_1`$ is a smooth variety of dimension $`sn+mnsm`$.
Proof. It remains to calculate that
$`dim\stackrel{~}{๐ฉ}`$ $`=`$ $`dimGr(m,V)+dimF_1`$
$`=`$ $`m(nm)+m^2(nm)s`$
$`=`$ $`sn+mnsm.`$
$`\mathrm{}`$
Using Proposition 5, the following theorem can now be proved in the same way as Theorem 3.
###### Theorem 4
The map $`\mu _1:\stackrel{~}{๐ฉ}_1๐ฉ`$ is a resolution of singularities.
Now we restrict ourselves to the case $`s=m`$ and thus $`n2m`$. Recall that the Kraft-Procesi quotient map from $`๐ฉ`$ to $`\overline{๐ช}_\lambda `$ is given by $`(A,B)BA`$. Denote by
$$\stackrel{~}{๐ช}_\lambda ^1=\{(g,U)\mathrm{}gU\mathrm{ker}g\}\overline{๐ช}_\lambda \times Gr(m,V),$$
and by $`p_1`$ the surjective projection from $`\stackrel{~}{๐ช}_\lambda ^1`$ to the first factor $`\overline{๐ช}_\lambda `$. We easily have the following identification
$$\begin{array}{ccc}\stackrel{~}{๐ช}_\lambda ^1\hfill & & T^{}Gr(m,V)\hfill \\ \hfill & & \hfill \\ & Gr(m,V)& \end{array}$$
where $`T^{}Gr(m,V)`$ denotes the cotangent bundle over the grassmannian.
Noting that the fiber of the projection $`p_1`$ over $`g๐ช_\lambda `$ consists of a single point. Since $`\stackrel{~}{๐ช}_\lambda ^1`$ is smooth, the projection $`p_1`$ from $`\stackrel{~}{๐ช}_\lambda ^1`$ to $`\overline{๐ช}_\lambda `$ is a resolution of singularities. Thus we have established the following classical result.
###### Proposition 6
The map $`p_1:\stackrel{~}{๐ช}_\lambda ^1\overline{๐ช}_\lambda `$ is a resolution of singularities.
The quotient map from $`๐ฉ`$ to $`\overline{๐ช}_\lambda `$ induces a quotient map by $`GL_m`$ from $`\stackrel{~}{๐ฉ}_1`$ to $`\stackrel{~}{๐ช}_\lambda ^1`$ which makes the following digram
$`\begin{array}{ccc}\stackrel{~}{๐ฉ}_1\hfill & \stackrel{\mu _1}{}& ๐ฉ\hfill \\ /GL_m\hfill & & /GL_m\hfill \\ \stackrel{~}{๐ช}_\lambda ^1\hfill & \stackrel{p_1}{}& \overline{๐ช}_\lambda \hfill \end{array}`$
commutative.
### 2.3 A third resolution of singularities of $`๐ฉ`$
We introduce the variety
$`\stackrel{~}{๐ฉ}_2`$ $`=`$ $`\{(A,B,U)๐ฉ\times Gr(ns,V)\mathrm{}BU\mathrm{ker}A\}`$
$``$ $`๐ฉ\times Gr(ns,V),`$
together with the surjective projection $`\mu _2`$ from $`\stackrel{~}{๐ฉ}_2`$ to the first factor $`๐ฉ`$. The variety $`\stackrel{~}{๐ฉ}_2`$ is a vector bundle over $`Gr(ns,V)`$ of total dimension $`sn+mnsm`$.
###### Theorem 5
The map $`\mu _2:\stackrel{~}{๐ฉ}_2๐ฉ`$ is a resolution of singularities.
Proofs of all the statements concerning $`\stackrel{~}{๐ฉ}_2`$ and $`\stackrel{~}{๐ช}_\lambda ^2`$ are similar to those for $`\stackrel{~}{๐ฉ}_1`$ and $`\stackrel{~}{๐ช}_\lambda ^1`$ in the previous subsection which we omit.
Now we restrict again to the case $`s=m`$ and $`n2m`$. Denote by
$$\stackrel{~}{๐ช}_\lambda ^2=\{(g,U)\mathrm{}gU\mathrm{ker}g\}\overline{๐ช}_\lambda \times Gr(nm,V).$$
We again can identify $`\stackrel{~}{๐ช}_\lambda ^2`$ as the cotangent bundle $`T^{}Gr(nm,V)`$. We can show that the natural surjective projection $`p_2`$ from $`\stackrel{~}{๐ช}_\lambda ^2`$ to $`\overline{๐ช}_\lambda `$ is a resolution of singularities.
The quotient map from $`๐ฉ`$ to $`\overline{๐ช}_\lambda `$ induces a quotient map by $`GL_m`$ from $`\stackrel{~}{๐ฉ}_2`$ to $`\stackrel{~}{๐ช}_\lambda ^2`$ which makes the following digram
$`\begin{array}{ccc}\stackrel{~}{๐ฉ}_2\hfill & \stackrel{\mu _2}{}& ๐ฉ\hfill \\ /GL_m\hfill & & /GL_m\hfill \\ \stackrel{~}{๐ช}_\lambda ^2\hfill & \stackrel{p_2}{}& \overline{๐ช}_\lambda \hfill \end{array}`$
commutative.
The relation among the three resolutions of $`๐ฉ`$ is shown by the following diagram:
$$\begin{array}{ccc}& \stackrel{~}{๐ฉ}& \\ q_1\hfill & & q_2\hfill \\ \stackrel{~}{๐ฉ}_1\hfill & & \stackrel{~}{๐ฉ}_2\hfill \\ \mu _1\hfill & & \mu _2\hfill \\ & ๐ฉ& \end{array}$$
where the morphisms $`q_1`$ and $`q_2`$ are defined by sending $`(A,B,U_1,U_2)`$ to $`(A,B,U_1)`$ and $`(A,B,U_2)`$ respectively.
Acknowledgment This paper was initiated when I was in Yale University in 1998. It is a pleasure to thank Roger Howe for stimulating discussions and encouragement.
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# FUSE Observations of Interstellar Gas Towards the LMC Star Sk -67 05
## 1. Introduction
The analysis of interstellar absorption lines provides fundamental information about the content and physical conditions of the interstellar medium (ISM). Absorption line spectroscopy can be used to study the ISM in our Galaxy, in nearby systems such as the Magellanic Clouds, and in the intergalactic medium out to the most distant QSOs. These studies provide an opportunity to compare elemental abundances and physical conditions in regions with differing chemical histories.
As part of a general program to investigate the interstellar medium of the Milky Way and nearby galaxies, we have used the Far Ultraviolet Spectroscopic Explorer (FUSE) satellite (Moos et al. 2000) to observe the star Sk -67 05 (HD 268605) in the Large Magellanic Cloud (LMC). The first observations of Galactic halo gas with IUE were along sightlines to stars in the LMC (Savage & de Boer 1979; 1981). In this Letter we discuss the first FUSE observations of an LMC star revealing O VI absorption in both the Milky Way and LMC. We also discuss the measurements of H<sub>2</sub>, P II, Si II, Ar I, and Fe II absorption along this sightline.
## 2. Observations and Data Processing
Sk -67 05 ($`l`$ = 278$`\stackrel{}{\mathrm{.}}`$89, $`b`$ = -36$`\stackrel{}{\mathrm{.}}`$32) is a B0 Ia star (Smith Neubig & Bruhweiler 1999) located near the western edge of the LMC. It lies in a diffuse H II region (Chu et al. 1994) with relatively low diffuse X-ray emission compared to regions closer to the center of the LMC (Snowden & Petre 1994). Ardeberg et al. (1972) give B V= -0.12, and for (B V)<sub>0</sub> = -0.23 (Binney & Merrifield 1998), we find E(B V) = 0.11. Using the Galactic gas-to-dust correlation (Diplas & Savage 1994) we infer $`N`$(H I) $`5.4\times 10^{20}\mathrm{cm}^2`$ along this sightline.
This star was observed during the In-Orbit Checkout phase of the mission (Sahnow et al. 2000) at various times between 1999 August 20 and 1999 October 19. As these observations were part of tests designed primarily to align the four optical channels in the instrument, the star was stepped across the $`30\mathrm{}\times 30\mathrm{}`$ aperture during the observations. However, the data were taken in time-tagged photon-address mode, so that corrections could be made for this image motion. The analysis here uses LiF1 (990 - 1080 ร
) data only; the LiF2 channel has lower sensitivity and poorer flatfield characteristics, and during this period no flatfield corrections were available. The SiC (905 - 1100 ร
) channels were generally not aligned. The instrument was still in its preflight focus configuration, and the spectral resolution was $`\lambda /\mathrm{\Delta }\lambda 15,000`$. The wavelength scale was established with a pre-flight dispersion solution. However, because of the stepping manner in which the data were obtained, additional zero point adjustments were required. Relative velocities are generally accurate to $``$10 km s<sup>-1</sup>over limited spectral ranges, and the absolute scale was set by comparison with IUE spectra. The data presented represent a total of approximately 33 ksec of on-target exposure time for the LiF1A and 35 ksec for the LiF1B spectral regions, approximately equally split between day and night. Additional details of the data processing for this data set can be found elsewhere (Massa et al. 2000).
This star exhibits variable stellar wind features. However, as shown in Figure 3 of Massa et al. (2000), the variablility is minimal at the wavelengths of interest here, and does not affect our conclusions.
## 3. Interstellar Absorption Features
Figure 1 shows the composite spectrum from the LiF1 channel. Several metal lines that appear in the ISM of both the Milky Way and LMC are identified. Most other lines are due to H<sub>2</sub> absorption in the LMC. The strong emission line is terrestrial Ly$`\beta \lambda 1025.72`$ airglow from the daytime exposures.
In this analysis we have used only the 1031.93 ร
member of the O VI doublet since the 1037.62 ร
line is blended with C II 1037.02 ร
and several H<sub>2</sub> lines. Figure 2 shows the FUSE O VI $`\lambda 1031.93`$ absorption line and a high-dispersion IUE spectrum of the C IV $`\lambda 1550.77`$. The H<sub>2</sub> (6-0) P(3) line originating in the LMC appears at +88 km s<sup>-1</sup>, and must be modelled and removed to get an accurate estimate of the O VI column density.
We have analyzed the (2-0) to (9-0) H<sub>2</sub> Lyman bands, using the method described in Shull et al. (2000) to determine the equivalent widths, $`b`$-values, and column densities for the Milky Way and LMC components of H<sub>2</sub>. Since the typical spacing between the rotational-vibrational lines within each Lyman band is very close to the spacing between LMC and the Milky Way absorption, the measurement of the Milky Way H<sub>2</sub> depends on a careful decomposition of blended lines. Because of this difficulty, we used only the (4-0) and (2-0) Lyman bands to measure the $`J=0,1`$ levels from the Milky Way gas.
The results of the H<sub>2</sub> model, discussed below, give $`\mathrm{log}`$ N(H<sub>2</sub>) = 15.28 for the (6-0) P(3) in the LMC. This was convolved with the instrumental resolution of 25 km s<sup>-1</sup> and divided out of the original spectrum to remove the effects of the H<sub>2</sub> absorption. The resulting O VI profile is shown as a light line in Fig. 2.
To calculate the O VI absorption we have considered two possible continuum placements, designated โhighโ and โlow,โ which are displayed in Fig. 2 as long-dashed and dash-dotted lines, respectively. The arrow at +180 km s<sup>-1</sup> denotes the velocity we have adopted as separating the Milky Way and LMC components of O VI. Table 1 gives O VI equivalent widths derived using both continua.
Figure 3 shows the spectra of several important metal lines and a molecular hydrogen line along the sightline to Sk -67 05. For comparison, the IUE spectrum of Si II $`\lambda `$1808.01 is shown. The individual absorption lines have separately been shifted in velocity to align the Milky Way components with the corresponding feature in the IUE spectrum. We established the velocity scale for the O VI region by shifting the Si II $`\lambda `$1020.70 line to match the Si II $`\lambda `$1808.01 IUE velocity scale, which sets the velocities of the nearby H<sub>2</sub> Lyman (7-0) P(2) $`\lambda `$1016.46 and P(3) $`\lambda `$1019.50 lines. The (6-0) P(3) $`\lambda `$1031.19 and R(4) $`\lambda `$1032.35 lines were then used to establish the O VI velocity, as shown in Fig. 2. The measured equivalent widths of several interstellar metal lines are also given in Table 1.
The adopted column densities for several ionic species, as well as several rotational states of H<sub>2</sub> along this sightline, are given in Table 2. These were calculated by fitting to a single-component Doppler-broadened curve of growth for Fe II and H<sub>2</sub>, and by using the apparent column density method (Savage & Sembach 1991) for the other species listed. The two values given for the O VI column are based on the high and low continuum placements. We have added a systematic error of 0.04 dex in quadrature with the statistical error for O VI to account for errors in continuum placement and the velocity interval over which the apparent column density is integrated.
## 4. Discussion
Since 114 eV are required to convert O<sup>+4</sup> to O<sup>+5</sup>, O VI is almost never produced by photoionization from starlight. Thus, it is a sensitive tracer of hot ($`300,000`$ K) collisionally ionized gas in the interstellar medium. Adopting the high continuum placement shown in Fig. 2, which we believe is more appropriate, the O VI column density is $`\mathrm{log}N(\text{}\text{VI})=14.40\pm 0.02`$ and $`13.89\pm 0.03`$ in the Milky Way and LMC, respectively. For the MW gas, $`\mathrm{log}[N(\text{}\text{VI})`$sin$`|b|]`$ = 14.17, which agrees well with the median value of 14.21 along 11 Galactic halo sightlines studied by Savage et al. (2000). Here $`b`$ = -36$`\stackrel{}{\mathrm{.}}`$32 is the galactic latitude.
Using $`\mathrm{log}N(\text{}\text{IV})=14.41`$ in the Milky Way from Wakker et al. (1998) and assuming an error of 0.05 dex, we find $`N(\text{}\text{IV})/N(\text{}\text{VI})=1.00\pm 0.16`$. If the low continuum is adopted, this ratio is $`1.23\pm 0.20`$. Either value is greater than, but consistent with, the halo value $`N(\text{}\text{IV})/N(\text{}\text{VI})0.6`$ determined from FUSE observations of four extragalactic sightlines (Savage et al. 2000), as well as $`N(\text{}\text{IV})/N(\text{}\text{VI})0.9`$ from Copernicus and IUE observations of stars in the lower halo (Spitzer, 1996).
Note that the widths of O VI components in both the MW and LMC are much broader than the thermal line width, which is $``$30 km s<sup>-1</sup> (FWHM) for gas at 300,000 K. The substantial difference in the widths of the C IV and O VI lines may be due to the different scale heights in the Galactic halo of these ions (Savage et al. 2000).
Sk -67 05 was the only star for which Wakker et al. (1998) did not detect C IV in their study of the LMC halo. We recalculated their upper limit to $`N(\text{}\text{IV})`$ assuming the velocity range observed in the O VI gas ($`+180`$ to $`+322`$ km s<sup>-1</sup>) and find $`\mathrm{log}N(\text{}\text{IV})<13.5(3\sigma )`$. Adopting the high continuum case we find $`N(\text{}\text{IV})/N(\text{}\text{VI})<0.4(3\sigma )`$ for the LMC material along the Sk -67 05 sightline.
The H<sub>2</sub> column density varies greatly between the LMC and Milky Way. For the LMC we find $`\mathrm{log}N(\mathrm{H}_2)=19.50\pm 0.08`$, summed over the $`J=0`$ to 5 states (Table 2), and a $`b`$-value of $`5\pm 2`$ km s<sup>-1</sup>. This column density is significantly higher than that measured along other LMC sightlines (de Boer et al. 1998; Shull et al. 2000). By comparison, FUSE observations of the metal deficient galaxy I Zw 18, have yielded only an upper limit, $`\mathrm{log}N(\mathrm{H}_2)15`$ (Vidal-Madjar et al. 2000).
The distribution of H<sub>2</sub> rotational states in the LMC is best fitted by a two-component excitation temperature (Fig. 4) with $`T_{01}=59\pm 5`$ K for $`J=0,1`$ and $`T_{ex}=800\pm 330`$ K for $`J=3,4,5`$. This is similar to the temperature distribution seen along the sightline to star LH10:3120 in the LMC using ORFEUS (de Boer et al. 1998). This indicates that the gas is not in thermal equilibrium, and the higher $`J`$ states are fluorescently pumped by UV radiation. This excited gas may therefore exist in the outer, optically thin regions of the cloud. In the interior regions the molecular fraction may increase due to self-shielding, which is expected when E(B V) $``$ 0.1 (Savage et al. 1977). The $`J=2`$ point falls almost exactly on the extension of the $`T_{01}`$ line, suggesting that the density of this gas is rather high.
For the Milky Way gas we find an H<sub>2</sub> column density of $`\mathrm{log}N(\mathrm{H}_2)=14.95\pm 0.08`$, summed over the $`J=0`$ to 3 states (Table 2), and a rotational temperature of $`T_{01}=99_{20}^{+30}`$ K (Fig. 4). This temperature is similar to other values measured with FUSE (Shull et al. 2000) and to Copernicus measurements of Galactic stars (Savage et al. 1977). There is an indication of a two-component temperature distribution for the Milky Way gas as well. However, because of the relatively small number of measurable lines and severe blending, we are unable to accurately determine an excitation temperature for the higher-$`J`$ MW material.
Absorption from both Milky Way and LMC material is clearly seen in all of the low ions present in the LiF1 data (e.g., Figs. 1 and 3). Blending of atomic, ionic, and H<sub>2</sub> absorption along this sightline limits the number of species for which we can derive accurate equivalent widths and column densities. The depleted species Fe II has several well-observed transitions in our dataset (Table 1). We have constructed a single-component Doppler-broadened curve of growth for the Fe II lines observed in both the Milky Way and the LMC. The best fit yields $`\mathrm{log}N(\text{Fe II})=14.8\pm 0.1`$ and $`14.8\pm 0.1`$ with $`b=14.2\pm 0.5`$ and $`12.1_{1.6}^{+2.1}`$ km s<sup>-1</sup> for the Milky Way and LMC, respectively. Due to potential uncertainties in some of the $`f`$-values we have adopted a conservative error of 0.1 dex in these column densities. Our data permit good measurements of the non-depleted species P II $`\lambda `$1152.82, and we derive lower limits to the Milky Way and LMC column densities (Table 2) using the apparent column density method of Savage & Sembach (1991). If the $`b`$-values for P II are similar to those derived for Fe II, the data suggest only very moderate saturation corrections of $`0.1`$ dex.
For the sightline through the halo of the Milky Way, we derive a gas-phase abundance $`[\mathrm{Fe}/\mathrm{P}]\mathrm{log}N(\text{Fe II})\mathrm{log}N(\text{P II})\mathrm{log}(\mathrm{Fe}/\mathrm{P})_{}0.6`$, assuming a solar system ratio of $`\mathrm{log}(\mathrm{Fe}/\mathrm{P})_{}=+1.92`$ (Anders & Grevesse 1989). For the LMC gas along this direction we find $`[\mathrm{Fe}/\mathrm{P}]0.7`$. This suggests significant incorporation of iron into dust grains in both galaxies. Differences in the relative abundance of singly- and doubly-ionized iron and phosphorous could affect this measurement (Sembach et al. 2000), but the magnitude of this effect is likely to be much too small to account for the gas-phase deficiency of iron along this sightline. The derived values of \[Fe/P\] are similar to the values of $`[\mathrm{Fe}/\mathrm{Zn}]=1.06`$ and $`0.91`$ for the Milky Way halo and LMC absorption towards SN 1987A (Welty et al. 1999). Future FUSE observations of a large number of LMC sightlines will allow us to study the distribution of gas-phase abundances and infer the composition of interstellar dust in this environment.
## 5. Summary
We have reported equivalent widths and column densities of O VI, H<sub>2</sub>, P II, Si II, Ar I, and Fe II along the line of sight to Sk -67 05 in the LMC using FUSE data. The principal results of this study are:
1. In the halo of the Milky Way toward the LMC $`N(\text{}\text{IV})/N(\text{OVI})=1.00\pm 0.16`$, a value somewhat greater than but consistent with other FUSE observations through the halo (Savage et al. 2000). In the LMC, where only an upper limit on $`N(\text{}\text{IV})`$ is available, we find $`N(\text{}\text{IV})/N(\text{}\text{VI})<0.4(3\sigma )`$. This is consistent with the lower ratio seen in the disk of the Milky Way compared to the halo (Savage et al. 2000; Spitzer 1996).
2. The LMC is rich in H<sub>2</sub> along this sightline, $`\mathrm{log}N(\mathrm{H}_2)=19.50\pm 0.08`$. A two-component temperature fit gives $`T_{01}=59\pm 5`$ K for and $`T_{ex}=800\pm 330`$ K.
3. The gas-phase abundances of $`[\mathrm{Fe}/\mathrm{P}]0.6`$ and $`0.7`$ in the Milky Way and LMC suggests significant depletion in both locations relative to solar abundances.
We thank D. Massa for providing the software to co-add the spectra used in this analysis. This work is based on data obtained for the Guaranteed Time Team by the NASA-CNES-CSA FUSE mission operated by the Johns Hopkins University. Financial support to U. S. participants has been provided by NASA contract NAS5-32985
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# 1 Introduction
## 1 Introduction
The possibility of our $`3+1`$-dimensional world being a cosmic defect (brane) in a higher-dimensional theory has recently attracted much interest. In particular, the observed hierarchy between the electroweak and the gravitational scales was considered in this context in .
In this paper, we show how a class of non-supersymmetric string vacua can naturally lead to an exponential hierarchy. In general we consider $`p+1`$-dimensional cosmic defects embedded in a $`D`$-dimensional spacetime. These types of models have been considered before in the literature . Here we show explicitly that such cosmic defects emerge as spacetime varying string vacua. The exponential hierarchy between the electroweak and gravitational scales arises from non-trivial warp factors in the metric and naturally is generated by the string coupling of $`O(1)`$. (The rรดle of warp factors in string theory and their relationship to cosmic brane models have been studied in .)
Our solutions resemble the stringy cosmic strings <sup>5</sup><sup>5</sup>5Historically these solutions were studied in four dimensions in which the defects correspond to string like objects. In the remainder of this article we will use the notation cosmic brane as this is more appropriate for the current context. of . The general framework, described in section 2, is that of a higher-dimensional string theory compactified on a Calabi-Yau (complex) $`n`$-fold, $`M_n`$, some moduli of which are allowed to vary over part of the noncompact space. In the uncompactified Type IIB theory, the rรดle of space-dependent moduli is played by the dilaton-axion system, very much like in Vafaโs description of F-theory . The cosmic defect (brane) appears as a singularity of the induced spacetime metric, with its characteristics governed by the energy momentum tensor of the moduli. (In addition, a naked singularity is located at a finite proper distance from the core of the brane.) However, our solutions are non-BPS. This points out to possible stability problems whose detailed study we defer to a future work .
The paper is organized as follows: In section 2, we recall the construction of spacetime varying string vacua and generalize the Ansรคtze of Refs. for the metric; we find non-trivial solutions in which the moduli are non-holomorphic functions over the noncompact space. While this breaks supersymmetry, the exponential warp factors in the metric induce a large hierarchy between the Planck scales in the higher- and lower-dimensional spacetime. In section 3, we study our spacetime-varying dilaton-axion solution when further compactified on a $`K3`$ or $`T^4`$ to four dimensions, and the generation of an exponential hierarchy. We briefly consider the case of compactifying string theory on a spacetime-varying $`K3`$, which also gives rise to $`3+1`$-dimensional vacua with exponential hierarchy. Finally, in section 4, we discuss various possible generalizations of our present work.
## 2 Spacetime Varying Vacua and Warp Factors
Let us consider compactifications of string theory in which the โinternalโ space (a Calabi-Yau $`n`$-fold $`M_n`$) varies over the โobservableโ spacetime. The parameters of the โinternalโ space then become spacetime variable moduli fields $`\varphi ^\alpha `$. The effective action describing the coupling of moduli to gravity of the observable spacetime can be derived by dimensionally reducing the higher dimensional Einstein-Hilbert action . In this procedure one retains the dependence of the Ricci scalar in the gravitational action only on the moduli $`\varphi ^\alpha `$. Then, the relevant part of the low-energy effective $`D`$-dimensional action of the moduli of the Calabi-Yau $`n`$-fold, $`M_n`$, coupled to gravity reads <sup>6</sup><sup>6</sup>6As a concrete example, the reader may think of a Type II string theory compactified on a $`T^2`$, in which case $`D=8`$ and $`n=1`$, or on a $`K3`$ or $`T^4`$, when $`D=6`$ and $`n=2`$, etc.
$$S_{\mathrm{eff}}=d^Dx\sqrt{g}(\frac{1}{2}R\frac{1}{2}G_{\alpha \overline{\beta }}g^{\mu \nu }_\mu \varphi ^\alpha _\nu \varphi ^{\overline{\beta }}+\mathrm{}),$$
(1)
where $`\mu ,\nu =0,\mathrm{},D1`$. We neglect higher derivative terms and set the other fields in the theory to zero as in .
We will restrict the moduli to depend on $`x_i`$, $`i=D2,D1`$, so that $`_a\varphi =0`$, $`a=0,\mathrm{},D3`$. The equations of motion are
$$g^{ij}\left(_i_j\varphi ^\alpha +\mathrm{\Gamma }_{\beta \gamma }^\alpha (\varphi ,\overline{\varphi })_i\varphi ^\beta _j\varphi ^\gamma \right)=0,$$
(2)
and
$$R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R=T_{\mu \nu }(\varphi ,\overline{\varphi }),$$
(3)
where the energy-momentum tensor of the moduli is
$$T_{\mu \nu }=G_{\alpha \overline{\beta }}\left(_\mu \varphi ^\alpha _\nu \varphi ^{\overline{\beta }}\frac{1}{2}g_{\mu \nu }g^{\rho \sigma }_\rho \varphi ^\alpha _\sigma \varphi ^{\overline{\beta }}\right).$$
(4)
It is useful to define $`z(x_{D2}+ix_{D1})`$ and rewrite the effective action as
$$S_{\mathrm{eff}}=\frac{1}{2}d^Dx\sqrt{g}R\frac{1}{2}d^{D2}xE$$
(5)
where $`d^{D2}x`$ refers to the integration measure over the first $`D2`$ coordinates, $`x_0,\mathrm{},x_{D3}`$. $`E`$, which we later interpret as the energy density (tension) of the cosmic brane, is given by
$$Ed^2z\sqrt{g}G_{\alpha \overline{\beta }}g^{z\overline{z}}\left(_z\varphi ^\alpha _{\overline{z}}\varphi ^{\overline{\beta }}+_{\overline{z}}\varphi ^\alpha _z\varphi ^{\overline{\beta }}\right).$$
(6)
In order to solve the coupled equations of motion for the moduli and gravity, Eqs. (2) and (3), we start with the following Ansatz for the metric
$$ds^2=e^{2A(z,\overline{z})}\eta _{ab}dx^adx^b+e^{2B(z,\overline{z})}dzd\overline{z}.$$
(7)
This type of Ansatz has appeared recently in various field theory , supergravity inspired scenarios and in the context of string theory . The authors in considered the possibility of having a superpotential and hence a potential for the scalar fields. The scalars in our effective action are assumed to be true moduli, with no (super)potential.
Using this Ansatz, Eqs. (2) and (3) produce the following . The โ$`zz`$โ Einstein equation leads to ($`_z`$ and $`\overline{}_{\overline{z}}`$)
$$(D2)\left[2AB(A)^2^2A\right]=2G_{\alpha \overline{\beta }}\varphi ^\alpha \varphi ^{\overline{\beta }}$$
(8)
and the โ$`z\overline{z}`$โ Einstein equation becomes
$$(D2)\left[\overline{}A+(D2)A\overline{}A\right]=0.$$
(9)
Note that the โ$`\overline{z}\overline{z}`$โ equation is obtained straightforwardly from (8) by replacing $``$ by $`\overline{}`$. We consider the case when $`D>2`$. Then, Eq. (9) lets us simplify the โ$`ab`$โ Einstein equation into
$$(D3)\overline{}A+2\overline{}B=G_{\alpha \overline{\beta }}(\varphi ^\alpha \overline{}\varphi ^{\overline{\beta }}+\overline{}\varphi ^\alpha \varphi ^{\overline{\beta }}).$$
(10)
Finally, the equation for the moduli reads
$$(D2)\left[2A\overline{}\varphi ^\alpha +2\varphi ^\alpha \overline{}A\right]+4\overline{}\varphi ^\alpha +4G^{\alpha \overline{\gamma }}\frac{G_{\delta \overline{\gamma }}}{\varphi ^\beta }\varphi ^\beta \overline{}\varphi ^\delta =0.$$
(11)
Let us start with the well-known supersymmetric solution by requiring that $`A=0`$ and $`\overline{}\varphi ^\alpha =0`$ . Because of holomorphicity, we can simplify the energy density (6). In the case of a variable Calabi-Yau $`n`$-fold $`M_n`$
$$G_{\alpha \overline{\beta }}_\alpha _\beta K,K\mathrm{log}(i^n\mathrm{\Omega }\overline{\mathrm{\Omega }}),$$
(12)
where $`\mathrm{\Omega }`$ is the appropriate holomorphic $`n`$-form on $`M_n`$. Therefore, Eq. (6) simplifies to
$$E=i\overline{}\mathrm{log}(i^n\mathrm{\Omega }\overline{\mathrm{\Omega }}).$$
(13)
Note that $`\sqrt{g}g^{z\overline{z}}=1`$ because $`A=0`$. Finally, one can show that the tension is equal to the deficit angle, $`\delta \theta `$, caused by the cosmic brane . As expected this saturates the BPS bound,
$$E\delta \theta .$$
(14)
In order to explicitly solve the equations of motion we consider $`M_1=T^2`$, fix the Kรคhler structure and focus on the complex structure modulus, $`\tau `$, for which the metric on the moduli space is
$$G_{\tau \overline{\tau }}=\frac{1}{(\tau \overline{\tau })^2}.$$
(15)
This is the case considered in . There, the authors studied a Type II string theory compactified on either a fixed $`K3`$ or $`T^4`$, followed by a further $`T_z^2`$ compactification from six to four dimensions. The subscript $`z`$ indicates that the torus varies over two ($`x_{2,3}`$) of the remaining four coordinates, $`x_i,i=0,\mathrm{},3`$. When the moduli of $`T_z^2`$ depend holomorphically on $`z=x_2+ix_3`$ cosmic brane solutions with finite energy density which saturate the BPS-bound (14), emerge at particular points in the $`z`$-plane. Since $`\tau =\tau (z)`$ Eq. (11) for the modulus $`\tau `$ is immediately satisfied. The only remaining equation, (10), reads
$$\overline{}(2B)=\overline{}\mathrm{log}(\tau \overline{\tau })$$
(16)
which can be explicitly solved and gives rise to the stringy cosmic branes .
We are however interested in other solutions for which $`A,B0`$, but still $`\overline{}\varphi ^\alpha =0`$. Although the $`\varphi ^\alpha `$ are holomorphic, this does not necessarily guarantee that the solution is supersymmetric . (For example, the general analysis of the supersymmetry transformations in $`D=10`$ type IIB supergravity implies that supersymmetry is broken if the modulus $`\tau `$, the axion-dilaton system, is non-holomorphic, or the warp factor $`A0`$ .) The holomorphicity of $`\varphi ^\alpha `$ simplifies the equations of motion considerably. In particular, from Eq. (11) we have that moreover $`\varphi ^\alpha =`$const, and one finds the general solution:
$`e^{2A}`$ $`=`$ $`\left({\displaystyle \frac{a(z)+\overline{a}(\overline{z})}{2a_0}}\right)^{\frac{2}{D2}},`$ (17)
$`e^{2B}`$ $`=`$ $`|b_0|^2\left({\displaystyle \frac{2a_0}{a(z)+\overline{a}(\overline{z}))}}\right)^{\frac{D3}{D2}}|a(z)|^2,`$ (18)
where $`a_0`$ and $`b_0`$ are integration constants and $`a(z)`$ is an arbitrary holomorphic function. Because $`\varphi ^\alpha =`$const the geometry of the moduli has decoupled. We therefore turn to non-supersymmetric solutions with non-holomorphic moduli. In order to simplify the problem we look for an axially symmetric solution in the plane perpendicular to the cosmic brane. The metric (7) then becomes
$$ds^2=e^{2A(\rho )}\eta _{ab}dx^adx^b+\mathrm{}^2e^{2\stackrel{~}{B}(\rho )}\left(d\rho ^2+d\theta ^2\right),$$
(19)
where $`z=r_0e^{\rho _0\rho +i\theta }`$ and $`\stackrel{~}{B}(\rho )=B(z,\overline{z})\rho `$; the constant parameter $`\mathrm{}=r_0\mathrm{exp}(\rho _0)`$ has dimensions of length. Note that this change of variables corresponds to choosing $`a(z)=\mathrm{log}(z/\mathrm{})`$ and $`b_0=\mathrm{}\stackrel{~}{b}_0`$ in Eqs. (17) and (18). The Einstein equations (8)-(10) then become
$`2\left[\left(\genfrac{}{}{0pt}{}{D2}{2}\right)(A^{})^2+(D3)A^{\prime \prime }+\stackrel{~}{B}^{\prime \prime }\right]`$ $`=`$ $`G_{\alpha \overline{\beta }}(_\rho \varphi ^\alpha _\rho \varphi ^{\overline{\beta }}+_\theta \varphi ^\alpha _\theta \varphi ^{\overline{\beta }}),`$ (20)
$`2(D2)\left[\frac{(D3)}{2}(A^{})^2+A^{}\stackrel{~}{B}^{}\right]`$ $`=`$ $`+G_{\alpha \overline{\beta }}(_\rho \varphi ^\alpha _\rho \varphi ^{\overline{\beta }}_\theta \varphi ^\alpha _\theta \varphi ^{\overline{\beta }}),`$ (21)
$`2(D2)\left[\frac{(D1)}{2}(A^{})^2+A^{\prime \prime }A^{}\stackrel{~}{B}^{}\right]`$ $`=`$ $`G_{\alpha \overline{\beta }}(_\rho \varphi ^\alpha _\rho \varphi ^{\overline{\beta }}_\theta \varphi ^\alpha _\theta \varphi ^{\overline{\beta }}),`$ (22)
while the equation for the moduli (11) reads
$$2(D2)A^{}_\rho \varphi ^\alpha +2(_\rho ^2+_\theta ^2)\varphi ^\alpha +2G^{\alpha \overline{\gamma }}\frac{G_{\delta \overline{\gamma }}}{\varphi ^\beta }(_\rho \varphi ^\beta _\rho \varphi ^\delta +_\theta \varphi ^\beta _\theta \varphi ^\delta )=0.$$
(23)
We have used the fact that the two warp factors, $`A,\stackrel{~}{B}`$, in Eq. (19) depend only on $`\rho `$, and abbreviated $`A^{}_\rho A`$, etc.
To solve these equations we will consider a single modulus scenario, $`\varphi ^\alpha =\tau `$ in which $`\tau =a+ie^\varphi `$, i.e., the axion-dilaton system of the $`D=10`$ Type IIB string theory. The holomorphic solution $`\tau =\tau (z)`$ is that of $`D7`$-branes . A more general solution can be obtained along the lines of our earlier analysis. In terms of the $`\rho `$ and $`\theta `$ variables we have for non-holomorphic $`\tau `$ and $`A,\stackrel{~}{B}0`$ the following solution,
$`e^{2A(\rho )}`$ $`=`$ $`({\displaystyle \frac{\rho }{\rho _0}})^{\frac{2}{D2}}`$ (24)
$`e^{2\stackrel{~}{B}(\rho )}`$ $`=`$ $`({\displaystyle \frac{\rho _0}{\rho }})^{\frac{D3}{D2}}e^{(\rho _0^2\rho ^2)/\rho _0}.`$ (25)
This is obtained by requiring that $`_\rho \tau =0`$ in (23<sup>7</sup><sup>7</sup>7If $`_\theta \tau =0`$ one obtains a different solution without the factor $`\mathrm{exp}((\rho _0^2\rho ^2)/\rho _0)`$ essentially reproducing (18) but with a different power of $`\rho `$ . The modulus turns out to be functionally the same as $`\tau _{II}`$ in (28) below..
We find two solutions for $`\tau `$, depending on whether $`\tau `$ is purely imaginary or not <sup>8</sup><sup>8</sup>8The functional form of our solutions for the dilaton-axion system resembles somewhat the non-supersymmetric electric $`D7`$-branes solution of IIB supergravity . The actual form of the metric describing our solution in ten dimensions is also similar to the recently found classical background of the non-supersymmetric $`USP(32)`$ Type I string theory .. Also, note that Eqs. (20)โ(23) are unchanged if $`\varphi ^\alpha =\tau `$ is rescaled by a positive multiplicative (normalization) constant. Restricting to $`\tau =\overline{\tau }`$, we find the simple particular solution
$$\tau _I(\theta )=a_0+ig_0^1e^{\frac{2\sqrt{2}\theta }{\sqrt{\rho _0}}},\theta [\pi ,\pi ].$$
(26)
Note that when $`a_00`$ $`\tau _I`$ satisfies (23) because of the form the metric $`G_{\tau \overline{\tau }}`$ takes (15). The constants $`a_0`$ and $`g_0`$ are chosen such that rotations through the defining domain of $`\theta `$ induce an $`SL(2,\text{ZZ})`$ action on this solution:
$$\tau _I(+\pi )=\frac{a\tau (\pi ))+b}{\tau _I(\pi )+d}.$$
(27)
With $`\tau \overline{\tau }`$, we find the following particular solution
$$\tau _{II}(\theta )=\left(\frac{n}{2}\mathrm{coth}[\frac{2\sqrt{2}\pi }{\sqrt{\rho _0}}]\right)\frac{\pm \mathrm{sinh}[\frac{2\sqrt{2}\theta }{\sqrt{\rho _0}}]+i}{\mathrm{cosh}[\frac{2\sqrt{2}\theta }{\sqrt{\rho _0}}]},\theta [\pi ,\pi ],$$
(28)
where both the axion and the dilaton have a more complicated $`\theta `$-dependence. Rotations through the defining domain of $`\theta `$ induce a monodromy action on this solution:
$$\tau _{II}(+\pi )=\tau _{II}(\pi )\pm n.$$
(29)
Finally, before turning to a more detailed analysis of these solutions, we compute the RR-charge. For (26) the charge is obviously zero as the axion is a constant. This is an indication that this particular cosmic brane is not a $`D`$-brane <sup>9</sup><sup>9</sup>9It may well turn out, however, that we instead have a $`D`$$`\overline{D}`$-brane system. We thank N. Itzhaki for pointing this out to us.. Still, there are certain $`D7`$-brane configurations whose RR-charge are zero . For (28) we note that when $`\rho _0>>8\pi ^2`$ we can write the axion as $`a\pm n\theta /(2\pi )`$ from which we immediately read off the RR-charge as $`Q=\pm n`$. (In fact, it is easy to show that this holds for all values of $`\rho _0`$.) In this limit we also note that to leading order $`\mathrm{exp}(\varphi )=1`$. For a fixed $`\mathrm{}`$, i.e., a fixed stress-energy tensor, $`\rho _0\mathrm{}`$ corresponds to taking $`r_00`$ and hence the core to zero size (see below for identification of the core in our solution). We thus obtain a situation which is more familiar from that of a $`D7`$-brane, except for the existence of the singularities at $`\rho =0`$ and $`\rho =\mathrm{}`$. One may hope that string corrections may render these singularities harmless in such a way that the relation between our solutions and the standard $`D7`$-branes can be made clearer (see also the discussion in section 3). We will defer a more detailed study of these matters for future work .
## 3 Exponential Hierarchy
In the model considered above, where the dilaton and axion vary over $`z=x_8+ix_9`$, the metric
$$ds^2=(\frac{\rho }{\rho _0})^{\frac{2}{(D2)}}\eta _{ab}dx^adx^b+\mathrm{}^2(\frac{\rho _0}{\rho })^{\frac{(D3)}{(D2)}}e^{(\rho _0^2\rho ^2)/\rho _0}\left(d\rho ^2+d\theta ^2\right),$$
(30)
is identical to that of the global cosmic brane solution studied by the authors of Ref. . They considered a theory with a complex scalar coupled to gravity in which the global $`U(1)`$ is spontaneously broken. For the particular case of $`D=6`$ this gives rise to a cosmic 3-brane which naturally can produce an exponential hierarchy between the electroweak and gravitational scales. By compactifying $`x_4,\mathrm{},x_7`$ on a small fixed K3 or a $`T^4`$ in (30), we have thus derived this type of solution from string theory.
Let us now study our solution in more detail following . First, in hindsight it is not too surprising that our solution corresponds to a global cosmic brane. Since the Einstein tensor only depends on $`\rho `$ while $`T_{\mu \nu }`$ is independent of $`\rho `$ they have to be constant in order to satisfy (20)-(22). This is exactly the starting point of the model considered in . The location of the brane is at $`\rho =\mathrm{}`$, which corresponds to $`z=0`$. (Recall the change of variables in making the ansatz for the metric (19), $`|z|=r_0e^{\rho _0\rho }`$.) However, unlike the supersymmetric case, in which the $`D7`$-brane is a delta-function source for the stress-energy tensor, the global cosmic brane has finite extent. Roughly speaking the size is that of the location of the core, i.e., $`|z|r_0`$ or $`\rho \rho _0`$. As we will see $`\rho _0`$ (or equivalently $`r_0`$) plays an important role in relating the string coupling to the size of the exponential hierarchy.
Eq. (30) has two (naked) curvature singularities, one at the center of the core, $`z=0`$ or $`\rho =\mathrm{}`$, and one at $`|z|=\mathrm{}`$ or $`\rho =0`$. In the latter case the spacetime ends on this naked singularity which is located at a finite proper distance from the core of the brane. According to Ref. , one can obtain a different solution by adding a small negative bulk cosmological constant. (This would correspond in our case to a negative correction to the flat potential on the moduli space.) At small distances the solution is that of since the negative contribution to the stress-energy tensor is negligible. However, the effect of the cosmological constant is important at large distances such that the spacetime becomes smooth. In the authors argue that the singularity at $`z=0`$ is unphysical since the solution does not satisfy the properties of a global cosmic brane; the stress-energy tensor is zero at the center of the core. To completely solve for the global cosmic brane one would have to match the above solution (30) to the metric inside the core and to the metric beyond $`|z|=\mathrm{}`$. Strictly speaking, we cannot claim to have a global cosmic brane solution from string theory until this is done. For now, however, we will not address this issue. We are hopeful that when string corrections are taken into account the singularities will be smoothed out <sup>10</sup><sup>10</sup>10One possibility is the enhanรงon mechanism . In our case, this would correspond to having zero stress-energy tensor inside the core..
We now compute the tension. For both of our two particular solutions, (26) and (28), the integrand in Eq. (6) is $`\frac{2}{\rho _0^2}\rho `$. Thus, upon integration from the core to the boundary of spacetime, $`0\rho \rho _0`$, assuming that the contribution from inside the core is finite, small and positive, we find $`E=2\pi +O(e^{\rho _0})`$. From (30), after changing back to the $`z`$-variable, it is easy to see that the deficit angle, $`\delta \theta =2\pi `$. Since $`E>\delta \theta `$ from Eq. (14) the solution is not BPS. Recall the monodromy outside the core given by Eqs. (27) and (29). These monodromies would in the supersymmetric case correspond to $`E_n`$ and $`\widehat{E}_n`$ (the affine extension of $`E_n`$) configurations for $`0n8`$ <sup>11</sup><sup>11</sup>11From the monodromy alone we cannot distinguish between $`K_{\tau _{I,II}}`$ (the monodromy matrix for $`\tau _{I,II}`$) and $`K_{\tau _{I,II}}`$. Hence, there are certain subtleties as far as identifying the precise $`D7`$-brane configuration. For now we will ignore this issue. which consist of $`n+2`$ and $`n+3`$ $`D7`$-branes respectively . However, these configurations cannot account for the deficit angle $`\delta \theta =2\pi `$. On the other hand, a combination of $`n+2`$ (or $`n+3`$) $`D7`$-branes placed at one, and $`10n`$ (or $`9n`$) $`D7`$-branes placed at the other singularity can produce the deficit angle of $`2\pi `$. In fact, the computation of the monodromy around the respective singularities tells us that the charges of these two collections of $`D7`$ branes should be opposite to each other. This indicates the possibility that the naked singularities could be resolved by replacing them with the above $`D7`$-branes, subject to the appropriate boundary conditions. In particular, the metric (30) would be matched to the metric of the respective $`D7`$-brane configurations at $`|z|=r_0`$ and $`|z|=\mathrm{}`$.
As noted above, the solution is by construction non-supersymmetric, and thus it is not in general protected against string corrections. However, the particularly interesting feature of this solution is that the ratio of the $`D2`$\- and $`D`$-dimensional Planck scales can be naturally large, as in Ref. . Let us therefore briefly review the analysis of . Let $`v`$ be a new coordinate, $`vv(\rho )`$:
$$v=\mathrm{}_0^\rho e^{(\rho _0^2\varrho ^2)/2\rho _0}(\frac{\rho _0}{\varrho })^{\frac{(D1)}{2(D2)}}๐\varrho =\mathrm{}e^{\rho _0/2}\left(\frac{\rho _0}{2}\right)^{\frac{3D5}{4(D2)}}\gamma (\frac{D3}{4(D2)},\frac{\rho ^2}{2\rho _0}),$$
(31)
where $`\gamma (\alpha ,\xi )\stackrel{\mathrm{def}}{=}_0^\xi \mathrm{dt}\mathrm{e}^\mathrm{t}\mathrm{t}^{\alpha 1}`$ is the incomplete (little) gamma-function. Note that we restrict the integration to $`0<\rho <\rho _0`$, i.e., where our solution is valid. Then the metric for our solution reads
$$ds^2=(\frac{\rho }{\rho _0})^{\frac{2}{(D2)}}\eta _{ab}dx^adx^b+(\frac{\rho }{\rho _0})^{\frac{2}{(D2)}}dv^2+\mathrm{}^2(\frac{\rho _0}{\rho })^{\frac{(D3)}{(D2)}}e^{(\rho _0^2\rho ^2)/\rho _0}d\theta ^2.$$
(32)
To examine the gravitational effects of the cosmic brane one considers metric perturbations of the form
$$ds^2=(\frac{\rho }{\rho _0})^{\frac{2}{(D2)}}(\eta _{ab}+h_{ab})dx^adx^b+(\frac{\rho }{\rho _0})^{\frac{2}{(D2)}}dv^2+\mathrm{}^2(\frac{\rho _0}{\rho })^{\frac{(D3)}{(D2)}}e^{(\rho _0^2\rho ^2)/\rho _0}d\theta ^2.$$
(33)
After imposing the gauge condition $`_ah^{ab}=h_a^a=0`$ (see e.g. ), we look for solutions of the form
$$h_{ab}=\epsilon _{ab}e^{ip_ax^a}e^{in\theta }\frac{\phi (v)}{\psi (v)},$$
(34)
where the polarization tensor $`\epsilon _{ab}=`$const.,
$$\psi (v(\rho ))=\mathrm{}^{1/2}e^{(\rho _0^2\rho ^2)/4\rho _0}(\frac{\rho }{\rho _0})^{\frac{D3}{4(D2)}},$$
(35)
and $`\eta _{ab}p^ap^b=m^2`$, as in Ref. .
The linearized Einstein equation for this type of fluctuations (34) reduces to
$$[\frac{d^2}{dv^2}+\frac{1}{\psi }\frac{d^2\psi }{dv^2}+n^2e^{(\rho ^2\rho _0^2)/\rho _0}(\frac{\rho }{\rho _0})^{\frac{D1}{D2}}]\phi (v)=m^2\phi (v).$$
(36)
As shown in Ref. , there is one normalizable zero mode ($`n=m=0`$), which corresponds to $`\varphi (v)=\psi (v)`$, and is well behaved near the singularity. (There is another zero mode which diverges near the singularity, which as argued by the authors of , can be eliminated by choosing appropriate boundary conditions. These boundary conditions prevent all conserved quantities from โleaking outโ through the boundary.) As pointed out in (see also ) the operator on the left-hand side of Eq. (36) for $`n=0`$ is positive semidefinite, and hence $`m^20`$.
Following , the existence of a normalizable zero mode allows us to obtain the relation between the four- and six-dimensional Planck scales <sup>12</sup><sup>12</sup>12The relation between the six and the ten-dimensional Planck scales is that of an ordinary Kaลuลผa-Klein compactification; $`M_6^4=M_{10}^8\text{Vol}(M_2)`$, where $`M_2`$ is either K3 or $`T^4`$.
$$M_4^2=M_6^4\psi ^2(v)๐v๐\theta =\pi \mathrm{}^2e^{\rho _0}\rho _0^{5/8}\mathrm{\Gamma }(\frac{3}{8})M_6^4.$$
(37)
If the upper limit of integration is $`\rho _0`$, $`\mathrm{\Gamma }(\frac{3}{8})`$ should be replaced with $`\gamma (\frac{3}{8},\rho _0)`$; however, the difference quickly becomes negligible if $`\rho _0>1`$. The exponential dependence of $`M_4/M_6`$ on $`\rho _0`$ produces an exponential hierarchy, except for small $`\rho _0`$, where (37) reverts to a power-law. Our non-supersymmetric 3-brane solution, derived explicitly from string theory, hence naturally provides for an exponential hierarchy. In addition, the exponential hierarchy is naturally related to a string coupling of $`O(1)`$. To see this, recall that in our solutions (26) and (28), the dilaton is of the form $`\varphi =2\sqrt{2}\theta /\sqrt{\rho _0}`$. Unless $`\rho _0`$ is very small this means that $`\varphi \pm O(1)`$. In particular, if $`\rho _0=8\pi ^2`$ then $`M_4/M_610^{18}(\mathrm{}M_6)`$.
The Fourier $`\theta `$-expansion in Eq. (34) is now easy to recognize as the Kaลuลผa-Klein mode expansion, especially since it is the $`n=0`$ mode that turns out to correspond to the four-dimensional graviton. The spectrum must be discrete since the volume of the transverse space is finite, and hence the mass gap is $`\delta m^2M_6^4/M_2^2`$. (For a more detailed analysis, see appendix A.) Note that although $`\delta m^21`$ (in units of $`M_6^2`$) the corrections to Newtonโs law are of Yukawa potential type and thus exponentially suppressed.
We now turn to the case of a varying K3 compactified string theory. (A varying $`T^4=T^2\times T^2`$ is easy to deal with by iterating the analysis from the previous section.) In this situation, the non-compact spacetime is $`D=6`$ dimensional. As before, the moduli of the K3 depend on $`z=x_5+ix_6`$ and its conjugate, $`\overline{z}`$. Rather than studying the full $`20`$-dimensional moduli space, we will focus on a one-parameter family relevant for discussing $`A_k`$ singularization of a $`K3`$ surface. In appendix B, the metric is worked out as a function of $`t`$, the deformation parameter in the defining equation for an $`A_k`$ singularity,
$$XY+U^{k+1}=t(z,\overline{z}).$$
(38)
One can show that, for $`|t|<(R^{})^2`$ (see appendix B), we have
$$\begin{array}{ccc}\hfill G_{t\overline{t}}& =& _t_{\overline{t}}(\mathrm{log}(K(t;2))\hfill \\ & =& \left(2R^2|t|^{\frac{2k}{k+1}}K(t;2)\right)^1+\left(4R^4|t|^{\frac{2(k1)}{k+1}}K^2(t;2)\right)^1.\hfill \end{array}$$
(39)
This is to be compared with the metric on the moduli space in the vicinity of a nodal $`T^2`$, $`XY=t(z,\overline{z})`$, (see appendix B):
$$\begin{array}{ccc}\hfill G_{t\overline{t}}& =& _t_{\overline{t}}\left(\mathrm{log}(\mathrm{log}(t)\mathrm{log}(\overline{t}))\right)\hfill \\ & =& [t\overline{t}\mathrm{log}^2(t\overline{t})]^1.\hfill \end{array}$$
(40)
As in the case of the $`T^2`$, we then proceed to obtain the equations of motion from the action. We get essentially the same solution, except that the Kรคhler metric $`G_{\alpha \overline{\beta }}`$ in (11) is replaced with the one in Eq. (39). Near the location of the singularity, it is enough to keep the leading term. If we consider the solution for which $`_\rho t=0`$, where $`\rho =\mathrm{log}r`$, $`z=re^{i\theta }`$ as before, and $`t=\overline{t}`$, one finds in the limit $`t0`$
$$^2t\frac{k}{(k+1)}\frac{1}{t}(t)^2=0.$$
(41)
This can be compared with the case for the torus
$$^2t\frac{1}{t}(t)^2=0.$$
(42)
The particular solution of the type (26) for Eq. (42) is of the form $`t=\stackrel{~}{c}e^{\stackrel{~}{w}\theta }`$, while that of (41) is $`t=\widehat{c}(\theta /(k+1))^{k+1}`$. Note that the latter approaches the former when $`k\mathrm{}`$. Furthermore, when taking $`t0`$ we are restricting $`t`$ to be real because of the way that $`t`$ is related to $`\tau `$, as $`t\mathrm{exp}(2\pi \tau _2)`$. The other solution can be dealt with in a similar fashion. It is worth pointing out that in this case one would expect gauge and matter degrees of freedom on the cosmic brane from wrapped $`D`$-branes on the shrinking cycles<sup>13</sup><sup>13</sup>13We thank S. Kachru for reminding us of this..
## 4 Discussion
To summarize, we have shown in this paper that exponential hierarchy can naturally arise from non-supersymmetric spacetime varying string vacua. The emergence of an exponential hierarchy is naturally related to having the string coupling of $`O(1)`$ <sup>14</sup><sup>14</sup>14However, we can in principle access all of the coupling space by taking $`\rho _0`$ very small in which case the dilaton $`\varphi `$ varies over all of the real line as $`\theta `$ goes from $`\pi `$ to $`\pi `$.. The four dimensional Planck scale is dynamically determined by the Planck scale in the bulk and the string coupling. In the $`D=10`$ Type IIB theory, the rรดle of space-dependent moduli is played by the dilaton-axion system. This solution, when further compactified on a $`K3`$ or $`T^4`$, leads to a four dimensional world previously considered only as a phenomenological solution. Note that the emergence of exponential hierarchy without demanding supersymmetry in this particular brane world scenario is akin to the similar phenomenon in technicolor.
There are many realizations of our scheme besides the straightforward compactification on $`K3`$ or $`T^4`$. For example, iterating non-trivially the above analysis, one might compactify $`x_8,x_9`$ on a $`T^2`$ and have its complex structure modulus, $`\tau `$, vary over $`z=x_6+ix_7`$, while its complexified Kรคhler class, $`\varrho `$, varies over $`w=x_4+ix_5`$. Each modulus then gives rise to a 5-brane intersecting in a cosmic 3-brane. Another possibility is to have three intersecting 7-branes in $`D=10`$ type IIB theory, repeating the earlier discussion for the individual 7-branes. Each brane depends on a different transverse complex plane and hence they intersect in a 3-brane <sup>15</sup><sup>15</sup>15It is interesting to note that because the branes carry RR-charge open strings stretch from one 7-brane to another. Since there are three branes this would give rise to three different types of matter multiplets, which in principle could account for the three generations observed in nature..
In all of our scenarios the issue of stability is clearly very important. One possibility is that the non-supersymmetric solutions we have found are unstable and that they decay to a supersymmetric set of $`(p,q)`$ $`D7`$-branes. It is natural to expect that certain properties would be preserved in this decay such as the $`SL(2,\text{ZZ})`$ monodromy transformation. DeWolfe et al have computed the monodromy for all possible types of $`(p,q)`$ $`D7`$-brane configurations. As discussed in section 3, the $`E_n`$ ($`\widehat{E}_n`$) and $`H_{8n}`$ ($`A_{8n}`$) singularities, given by sets of $`n+2`$ ($`n+3`$) and $`10n`$ ($`9n`$) isolated $`D7`$-branes, have monodromies, which coincide with the $`SL(2,\text{ZZ})`$ transformations of our $`\tau _I`$ ($`\tau _{II}`$) solution. This could be taken as an indication that our solutions are longlived excitations around these supersymmetric vacua. One could imagine that the naked singularities in our metric can be resolved in the following manner. After imposing appropriate boundary conditions we patch up the spacetime metric of the supersymmetric solutions with the non-singular part of our metric. The solution obtained in this way would hopefully have the essential features of our solution, i.e., exponential hierarchy. We hope to return to some of these issues in the future .
Notice that the warp factors in the expression for the metric that describes our solution depend only on one extra dimension. It seems natural to ask whether there exists a holographic renormalization group interpretation of the bulk equations of motion that describe this brane-world solution.
Finally, it would be interesting to see whether the recent discussion about the possible relaxation mechanisms for the cosmological constant applies to this particular class of string vacua. It would be also interesting to understand whether the conjectured relationship between holography and the cosmological constant problem can be made more specific in this situation.
Acknowledgments: We thank V. Balasubramanian, O. Bergman, R. Corrado, M. Dine, E. Gimon, J. Gomis, P. Hoลava, G. Horowitz, T. Imbo, N. Itzhaki, C. Johnson, S. Kachru, P. Mayr, A. Peet, J. Polchinski, E. Silverstein, S. Sethi, K. Sfetsos and S. Thomas for useful discussions. The work of P. B. was supported in part by the National Science Foundation under grant number PHY94-07194. P.B. would like to thank the Caltech/USC Center for Theoretical Physics, as well as Stanford University and University of Durham for their hospitality in the final stages of this project. T. H. wishes to thank the US Department of Energy for their generous support under grant number DE-FG02-94ER-40854 and the Institute for Theoretical Physics at Santa Barbara, where part of this work was done with the support from the National Science Foundation, under the Grant No. PHY94-07194. The work of D. M. was supported in part by the US Department of Energy under grant number DE-FG03-84ER40168. D. M. would like to thank Howard University and the University of Illinois at Chicago for their hospitality while this work was in progress.
## Appendix A Potential for massive gravitational modes
In this appendix we analyze the potential, $`\psi ^1d^2\psi /dv^2`$, in the linearized Einstein equation for the gravitational fluctuations (36).
Note that for $`\rho 0`$, we have the power series expansion
$$v=\mathrm{}\frac{1}{2}e^{\frac{\rho _0}{2}}\rho _0^{\frac{D1}{2(D2)}}\rho ^{\frac{D3}{2(D2)}}\underset{k=0}{\overset{\mathrm{}}{}}\frac{\left(\rho ^2/2\rho _0\right)^k}{k!\left(\frac{D3}{2(D2)}+k\right)},\rho 0,$$
(43)
so that $`\psi (\mathrm{}v/\rho _0)^{1/2}`$. While an asymptotic formula for $`v=v(\rho )`$ may be derived for $`\rho \mathrm{}`$, it turns out to be less reliable as it involves a formally divergent series.
Since the transcendental change of variables (31) is not invertible in closed form, we are unable to give a $`\psi `$ as an explicit function of $`v`$ in closed form. However, using Mathematicaโs ParametricPlot, we can plot $`\psi `$ vs. $`v`$, and the result is shown in Fig. 1.
Clearly, as a function of $`v`$, $`\psi (v)`$ is well approximated by
$$\psi (v;c)\stackrel{\mathrm{def}}{=}\sqrt{\mathrm{v}}(\frac{\mathrm{v}_{\mathrm{}}\mathrm{v}}{\mathrm{}})^{1/\mathrm{c}},$$
(44)
for a suitable $`c`$.
With this approximation, we obtain
$$\frac{1}{\psi }\frac{d^2\psi }{dv^2}\frac{1/4}{v^2}\frac{1/c}{v(v_{\mathrm{}}v)}\frac{(c1)/c^2}{(v_{\mathrm{}}v)^2},$$
(45)
which is plotted in Fig. 2.
Note that these terms may be thought of as the superposition of two attractive potentials, one at $`v=0`$ and the other at $`v=v_{\mathrm{}}`$. The attractive potential at $`v=0`$ (and so $`\rho =0`$, i.e., $`|z|=\mathrm{}`$) corresponds to a naked singularity a finite distance from the core at $`\rho =\rho _0`$ (or $`|z|=r_0`$), while the one at $`v_{\mathrm{}}`$ (i.e., $`\rho =\mathrm{}`$) corresponds to the other singularity at the origin of the $`z`$-plane.
The approximate potential may well define an exactly soluble quantum mechanical problem. For now, however, let us discuss a further simplification. Firstly, we neglect the middle term in Eq. (45), since it is dominated by the first and the last term, when $`v0`$ and $`vv_{\mathrm{}}`$, respectively. Secondly, except for the โchargeโ of the potential terms, $`\frac{1}{4}`$ for the first term and $`\frac{c1}{c^2}`$ for the last, the two terms represent the same type of divergence. Thus we analyze โthe first half of the potentialโ, i.e., $`V(v)=\frac{1}{4v^2}`$, and look for solutions of the Schrรถdinger equation
$$\left[\frac{d^2}{dv^2}\frac{1}{4v^2}\right]\phi ^{(0)}=m^2\phi ^{(0)},$$
(46)
of the form $`\phi ^{(0)}(v)=\sqrt{v}f(v)e^{bv}`$. Upon substitution, we find
$$\left[\frac{1}{4v^2}\frac{f^{}}{vf}+\frac{b}{v}\frac{f^{\prime \prime }}{f}+\frac{2bf^{}}{f}b^2\frac{1}{4v^2}m^2\right]\phi ^{(0)}=0.$$
(47)
This suggests the identification $`b=\sqrt{m^2}`$: for negative โenergyโ, $`m^2<0`$, $`b`$ is real and the solutions would exponentially decay; for positive โenergyโ, $`m^2>0`$, $`b`$ is imaginary and the solutions would have a plane wave factor. The resulting equation for $`f(v)`$ then is
$$vf^{\prime \prime }+(12bv)f^{}bf=0,$$
(48)
which is easily solved in terms of a power series, $`f=_{k=0}^{\mathrm{}}c_kv^{k+s}`$. Direct substitution yields
$`0`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}c_k(k+s)[(k+s1)+1]v^{k+s1}b{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}c_k[2(k+s)+1]v^{k+s1},`$ (49)
$`=`$ $`c_0(s)^2v^{s1}+{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}\left\{c_{k+1}(k+s+1)^2c_kb[2(k+s)+1]\right\}v^{k+s}.`$
The vanishing of the coefficient of the first term implies $`s=0`$, and the vanishing of the coefficients for the remaining terms imply the recursion relation:
$$c_{k+1}=c_kb\frac{2(k+s)+1}{(k+s+1)^2}.$$
(50)
Since
$$\frac{c_{k+1}v^{k+1}}{c_kv^k}=\frac{2k+1}{(k+1)^2}bv\frac{2bv}{k}<1,\text{when}k\mathrm{},$$
(51)
the series defined by the recursion relation (50) converges. Moreover, since large powers of $`v`$ dominate when $`v`$ is large, we have that $`f(v)e^{2bv}`$ when $`v`$ is large, since the power series of $`e^{2bv}`$ has the limiting ratio (51). Note that the asymptotically exponential dependence ensures that $`v`$ equaling a small fraction of $`v_{\mathrm{}}`$ is already sufficiently large for the asymptotic estimate (see Fig. 2). Thus, although the present solution is derived for the simplified toy potential in Eq. (46), it will turn out to be useful both for the full toy model, with the potential (45), and also for the real case, with the potential $`\psi ^{\prime \prime }/\psi `$.
As there is no choice of $`b=\sqrt{m^2}`$ for which the series would terminate, we conclude that $`\phi ^{(0)}(v)=\sqrt{v}f(v)e^{bv}\sqrt{v}e^{+bv}`$ for large $`v`$. For $`m^2<0`$, when $`b`$ is real and positive, all such solutions are unphysical as they are unnormalizable. (Recall that on retaining only the simple potential as shown in (46), we also let $`0v\mathrm{}`$.) For $`m^2>0`$, however, $`b`$ is imaginary, and $`\phi ^{(0)}(v)=\sqrt{v}f(v)e^{i\beta v}\sqrt{v}e^{i\beta v}`$, where $`\beta =|b|`$. Note in particular that $`\phi ^{(0)}(0)0`$.
Since the potential (45) is qualitatively symmetric with respect to $`\{v=0\}\{v=v_{\mathrm{}}\}`$, we expect to find appropriately โreflectedโ solutions around the attractive potential at $`v=v_{\mathrm{}}`$. Based on the result in the previous case, we also expect $`\phi ^{(\mathrm{})}(v)=(v_{\mathrm{}}v)^{1/c^{}}g(v)e^{b(v_{\mathrm{}}v)}`$, and in particular that $`\phi ^{(\mathrm{})}(v_{\mathrm{}})=0`$.
Then, the solutions for the full toy potential (45), and also for our actual problem (36) with $`n=0`$, should be well approximated by the linear superpositions $`\phi ^\pm \stackrel{\mathrm{def}}{=}\frac{1}{\sqrt{2}}(\phi ^{(0)}\pm \phi ^{(\mathrm{})})`$. Easily, $`\phi ^\pm (0)0\phi ^\pm (v_{\mathrm{}})`$โas if there existed impenetrable walls at $`v=0`$ and $`v=v_{\mathrm{}}`$. That is, the quantum mechanical problem with the toy potential (45), and also our actual problem (36) with $`n=0`$, behave qualitatively as if the wave-functions $`\phi (v)`$ were confined in a (smoothed) infinite potential well. Therefore, we expect a discrete โenergyโ spectrum $`m^20`$, and a mass gap between $`\psi `$, the $`n=0=m`$ solution of Eq. (36) and the lowest lying state with $`m^2>0`$.
## Appendix B Kรคhler Potential Singularities
In this appendix we study the singularities of the Weil-Petersson-Zamolodchikov Kรคhler potential on the complex structure moduli spaces for several cases of interest. This follows and extends the results of Ref. , using the hallmark of Calabi-Yau $`n`$-folds, $`M_n`$: their covariantly constant and nowhere vanishing $`(n,0)`$-form. Variations of this $`(n,0)`$-form describe the complex structure moduli space, in which a codimension-1 subspace (the discriminant locus) corresponds to singular $`n`$-folds.
### B.1 The potential near singularities
Following Ref. , the holomorphic $`(n,0)`$-form may be written as
$$\mathrm{\Omega }_t=\frac{\underset{ij}{}dx_i}{\varphi ,_j(x;t)},i,j=1,\mathrm{},n+1.$$
(52)
Here, $`x_i`$ are local coordinates on the $`n+1`$-fold, $`๐ด`$, in which the Calabi-Yau $`n`$-fold, $`M_n`$, is embedded as the hypersurface $`f(x;t)=0`$; the $`t`$ are local coordinates on the (complex structure) moduli space. Generalizations of this to complete intersections of hypersurfaces and other elements of cohomology are straightforward .
The integral $`K(t;n)\stackrel{\mathrm{def}}{=}\mathrm{i}^\mathrm{n}_{\mathrm{M}_\mathrm{n}}\overline{\mathrm{\Omega }}_\mathrm{t}\mathrm{\Omega }_\mathrm{t}`$ turns out to play a double rรดle in the models we consider. Its logarithm is the Kรคhler potential for the Weil-Petersson-Zamolodchikov metric on the complex structure moduli space . It is also the relative conformal factor in the spacetime metric for the line element transversal to the cosmic $`(72n)`$-branes ; the peculiar factor, $`i^n`$, ensures the Hermiticity of $`K(t;n)`$, as necessary for its latter rรดle.
As a toy model, compactify the spacetime coordinates $`x_8,x_9`$ on a $`T^2`$, the complex structure of which is permitted to vary over the space-like $`(x_6,x_7)`$-surface. To this end, $`M_1=T^2`$ may be defined as a cubic hypersurface $`f(y;t)=0`$ in $`๐ด=\mathrm{IP}^2`$, the coefficients of which are functions of the moduli, $`t`$, which in turn are functions of $`z=x_8+ix_9`$ and its conjugate:
$$\begin{array}{ccc}\hfill f(y;t)& \stackrel{\mathrm{def}}{=}& y_1^3+y_2^3+y_3^33t(z,\overline{z})y_1y_2y_3=0,\hfill \\ & & (y_1,y_2,y_3)(ลy_1,ลy_2,ลy_3).\hfill \end{array}$$
(53)
This torus becomes singular ($`d_yf=0`$) at $`t_n=e^{2n\pi i/3}`$, $`n=0,1,2`$, and $`t_{\mathrm{}}=\mathrm{}`$. At each of these points in the $`t`$-plane, and the corresponding points of the $`z`$-plane, this highly symmetric torus has three singular points, each of which may be described as a node, i.e., an $`A_1`$ singularity in the classification of Ref. . That is, first change variables in the $`t`$-plane so that a singularity of Eq. (53) is moved to $`t=0`$. Then, a holomorphic (but nonlinear) change of local coordinates on $`\mathrm{IP}^2`$ turns the defining equation, $`f(y;t)=0`$, into $`XY=t`$, up to higher order terms<sup>16</sup><sup>16</sup>16For example, let $`t1/3t`$, work in the coordinate patch where $`y_30`$, divide through by $`y_3^3`$, define $`X=y_1/y_3`$ and $`Y=y_2/y_3`$ and neglect $`X^3,Y^3`$ as compared to 1.. For nonzero $`t`$, the solution of this equation is a rotational hyperboloid. As $`t0`$, this pinches into a two-sheeted cone with the vertexโthe singularityโat $`X=0=Y`$. The holomorphic $`(1,0)`$-form may here be written as
$$\mathrm{\Omega }_t=f(Y)\frac{dY}{Y},\overline{}_Yf(Y)=0,$$
(54)
this expression being valid within the neighborhood $`|X|,|Y|<R`$. Furthermore, since $`X=t/Y`$, we have
$$R>|X|=\frac{|t|}{|Y|},|Y|>\frac{|t|}{R},$$
(55)
so that $`\frac{|t|}{R}<|Y|<R`$, and therefore also $`|t|<R^2`$.
The integral $`i_{M_1}\overline{\mathrm{\Omega }}_t\mathrm{\Omega }_t`$, as a function of $`t`$, may then be estimated by dividing the integral into the contributions from each neighborhood in which a singularity develops, and the remaining part which is regular in $`t`$. The singularities all being equal, locally, we calculate the contribution from the one at $`t=0`$ and multiply by the number of them, $`N`$. Therefore:
$$K(t;1)\stackrel{\mathrm{def}}{=}\mathrm{i}_{\mathrm{M}_1}\overline{\mathrm{\Omega }}_\mathrm{t}\mathrm{\Omega }_\mathrm{t}=\mathrm{V}_0(\mathrm{t})+\mathrm{iN}_{|\mathrm{t}|/\mathrm{R}<|\mathrm{Y}|<\mathrm{R}}\left|\mathrm{f}(\mathrm{Y})\right|^2\frac{\mathrm{dY}\mathrm{d}\overline{\mathrm{Y}}}{|\mathrm{Y}|^2}.$$
(56)
Following Ref. , we expand $`f(Y)`$ in a Taylor series, and note that all the positive powers integrate to regular functions of $`t`$. We may therefore write
$$i_{M_1}\overline{\mathrm{\Omega }}_t\mathrm{\Omega }_t=V(t)+2N|f_0|^2_{|t|/R<|Y|<R}\frac{d^2Y}{|Y|^2},$$
(57)
where all the higher order contributions from $`[f(Y)f_0]`$ are absorbed in the regular function $`V(t)`$. Writing $`Y=re^{i\theta }`$, we easily obtain:
$$_{|t|/R<|Y|<R}\frac{d^2Y}{|Y|^2}=_0^{2\pi }๐\theta _{|t|/R}^R\frac{rdr}{r^2}=2\pi \mathrm{log}\left(\frac{R^2}{|t|}\right)$$
(58)
Upon the transformation $`t=R^2j(\tau )R^2e^{2\pi i\tau }`$ near $`t=0`$, i.e., near $`\mathrm{}m(\tau )=\mathrm{}`$,
$$K(t;1)\stackrel{\mathrm{def}}{=}\mathrm{i}_{\mathrm{M}_1=\mathrm{T}^2}\overline{\mathrm{\Omega }}_\mathrm{t}\mathrm{\Omega }_\mathrm{t}\mathrm{}\mathrm{m}(\tau ),$$
(59)
which is indeed the standard Kรคhler potential in the Teichmรผller theory <sup>17</sup><sup>17</sup>17Recall that for the Calabi-Yau 1-fold, the 2-torus $`M_1=T^2`$, the space of complex structures may be parameterized as the $`|\tau |1`$ portion of the $`\frac{1}{2}\mathrm{}e(\tau )+\frac{1}{2}`$ strip in the upper half-plane $`\mathrm{}m(\tau )>0`$..
Before we turn to higher dimensional cases, it is useful to consider the action in terms of the local coordinate $`tj(\tau )`$. This will facilitate comparison, as we will momentarily see, with the metric for the $`A_k`$ singularities. The Kรคhler metric in the $`t`$ coordinate near $`t=0`$ is given by
$$\begin{array}{ccc}\hfill G_{t\overline{t}}& =& _t_{\overline{t}}\left(\mathrm{log}[\mathrm{log}(t)\mathrm{log}(\overline{t})]\right)\hfill \\ & =& \left(|t|\mathrm{log}|t|^2\right)^2.\hfill \end{array}$$
(60)
Recall that the metric in terms of the $`\tau `$ coordinate is given by $`G_{\tau \overline{\tau }}=[\mathrm{}m(\tau )]^2`$.
We now repeat the calculation of $`K(t;2)\stackrel{\mathrm{def}}{=}_{\mathrm{K3}}\mathrm{\Omega }_\mathrm{t}\overline{\mathrm{\Omega }}_\mathrm{t}`$, for the Calabi-Yau (complex) 2-folds, the K3 surfaces. The first marked distinction is that there can now be many more types of singularities . Herein, we consider the $`A_k`$ type, for which the defining equation may be brought into the general form<sup>18</sup><sup>18</sup>18In Refs. , the modulus $`t(z,\overline{z})`$ has been chosen to be a holomorphicโmoreover linearโfunction of $`z`$. This need not be so in general, and indeed our main solutions are non-holomorphic.
$$XY+U^{k+1}=t(z,\overline{z}),$$
(61)
within the neighborhood $`|X|,|Y|,|U|^{\frac{k+1}{2}}<R`$. Solving for $`X`$, we have that
$$R>|X|=\frac{|tU^{k+1}|}{|Y|},|Y|>\frac{|tU^{k+1}|}{R},$$
(62)
which gives a lower bound for $`|Y|`$. The upper bound, $`|Y|<R`$, then implies through Eq. (62)
$$\frac{|tU^{k+1}|}{R}<|Y|<R,|tU^{k+1}|<R^2,$$
(63)
which is not guaranteed by $`|U|^{\frac{k+1}{2}}<R`$ when $`t0`$. We thus take $`|U|^{\frac{k+1}{2}}<R^{}`$, where $`R^{}`$ is sufficiently smaller than $`R`$. Therefore, with $`N`$ isolated singular points of the $`A_k`$ type,
$$K(t;2)=V(t)+2N|f_0|^2_{|U|^{\frac{k+1}{2}}<R^{}}d^2U_{|tU^{k+1}|/R<|Y|<R}\frac{d^2Y}{|Y|^2}.$$
(64)
The $`Y`$-integral straightforwardly gives $`2\pi \mathrm{log}\left(\frac{R^2}{|U^{k+1}t|}\right)`$, and we are left with the $`U`$-integral:
$$K(t;2)=V(t)+4\pi N|f_0|^2_{|U|^{\frac{k+1}{2}}<R^{}}d^2U\mathrm{log}\left(\frac{R^2}{|U^{k+1}t|}\right).$$
(65)
The result of the $`U`$-integration depends on whether $`|U|^{k+1}<|t|`$ or $`|U|^{k+1}>|t|`$. In the former case, we write
$$|U|^{k+1}<|t|:\mathrm{log}(|U^{k+1}t|)=\mathrm{log}|t|+\mathrm{log}|1\frac{U^{k+1}}{t}|,$$
(66)
while in the latter case we write
$$|U|^{k+1}>|t|:\mathrm{log}(|U^{k+1}t|)=(k+1)\mathrm{log}|U|+\mathrm{log}|1\frac{t}{U^{k+1}}|.$$
(67)
In both cases, the second logarithm leads to an integral of the general form ($`w=\rho e^{i\theta }`$)
$`{\displaystyle _{|w|<C}}d^2w\mathrm{log}\left|1aw^b\right|,|aw^b|<1,`$ (68)
$`=`$ $`\frac{1}{2}{\displaystyle _0^{2\pi }}๐\theta {\displaystyle _0^C}\rho ๐\rho \left[\mathrm{log}(1a\rho ^be^{ib\theta })+\mathrm{log}(1\overline{a}\rho ^be^{ib\theta })\right],`$
$`=`$ $`\frac{1}{2}{\displaystyle _0^C}\rho ๐\rho {\displaystyle _0^{2\pi }}๐\theta \left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}a^n\rho ^{nb}e^{inb\theta }+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}\overline{a}^n\rho ^{nb}e^{inb\theta }\right]=0.`$
Now, if $`R^{}<|t|^{\frac{1}{2}}`$, then also $`|U|^{\frac{k+1}{2}}<R^{}<|t|^{\frac{1}{2}}`$, and we obtain:
$`{\displaystyle _{|U|^{\frac{k+1}{2}}<R^{}}}d^2U\mathrm{log}\left({\displaystyle \frac{R^2}{|U^{k+1}t|}}\right)`$ (69)
$`=`$ $`{\displaystyle _{|U|^{\frac{k+1}{2}}<R^{}}}d^2U\left(\mathrm{log}(R^2)\mathrm{log}|t|\mathrm{log}\left|1{\displaystyle \frac{U^{k+1}}{t}}\right|\right),`$
$`=`$ $`\pi (R^{})^{\frac{4}{k+1}}\mathrm{log}\left({\displaystyle \frac{R^2}{|t|}}\right),(R^{})^2<|t|,`$
where we used Eqs. (66) and (68) for the second and the third equality.
On the other hand, if $`|t|<(R^{})^2`$, then for part of the integral $`|U|^{\frac{k+1}{2}}<|t|^{\frac{1}{2}}<R^{}`$, and in the remaining part $`|t|^{\frac{1}{2}}|U|^{\frac{k+1}{2}}<R^{}`$. Splitting the integration accordingly, we obtain:
$`{\displaystyle _{|U|^{\frac{k+1}{2}}<R^{}}}d^2U\mathrm{log}\left({\displaystyle \frac{R^2}{|U^{k+1}t|}}\right)`$ (70)
$`=`$ $`{\displaystyle _{|U|^{\frac{k+1}{2}}<R^{}}}d^2U\mathrm{log}(R^2){\displaystyle _{|U|^{\frac{k+1}{2}}<R^{}}}d^2U\mathrm{log}\left|U^{k+1}t\right|,`$
$`=`$ $`\pi (R^{})^{\frac{4}{k+1}}\mathrm{log}(R^2){\displaystyle _{|U|^{\frac{k+1}{2}}<|t|^{\frac{1}{2}}<R^{}}}d^2U\left(\mathrm{log}|t|+\mathrm{log}\left|1{\displaystyle \frac{U^{k+1}}{t}}\right|\right)`$
$`{\displaystyle _{|t|^{\frac{1}{2}}|U|^{\frac{k+1}{2}}<R^{}}}d^2U\left(\mathrm{log}|U|^{k+1}+\mathrm{log}\left|1{\displaystyle \frac{t}{U^{k+1}}}\right|\right),`$
$`=`$ $`2\pi (R^{})^{\frac{4}{k+1}}\mathrm{log}(R)\pi |t|^{\frac{2}{k+1}}\mathrm{log}|t|(k+1)2\pi \left[\frac{1}{2}|U|^2\mathrm{log}|U|\frac{1}{4}|U|^2\right]_{|t|^{\frac{1}{k+1}}}^{(R^{})^{\frac{2}{k+1}}},`$
$`=`$ $`\frac{k+1}{2}\pi \left((R^{})^{\frac{4}{k+1}}|t|^{\frac{2}{k+1}}\right)+2\pi (R^{})^{\frac{4}{k+1}}\mathrm{log}\left({\displaystyle \frac{R}{R^{}}}\right),|t|<(R^{})^2,`$
where we used Eqs. (66) and (67) in the second, and (68) in the third equality.
To summarize,
$$K(t;2)=V(t)+4\pi ^2N|f_0|^2(R^{})^{\frac{4}{k+1}}I(t;R,R^{}),$$
(71)
where $`t=t(z,\overline{z})`$ and
$$I(t;R,R^{})=\{\begin{array}{cc}\mathrm{log}\left(\frac{R^2}{|t|}\right),\hfill & |t|>(R^{})^2,\hfill \\ \frac{k+1}{2}\left[1\left(\frac{|t|}{(R^{})^2}\right)^{\frac{2}{k+1}}\right]+2\mathrm{log}\left(\frac{R}{R^{}}\right),\hfill & |t|<(R^{})^2,\hfill \end{array}$$
(72)
Notice that this function is continuous at $`|t|=(R^{})^{k+1}`$, albeit not smooth.
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# Untitled Document
hep-th/0005091 IASSNS-HEP-00/38
A Derivation of $`K`$-Theory from $`M`$-Theory
Duiliu-Emanuel Diaconescu
School of Natural Sciences, Institute for Advanced Study
Olden Lane, Princeton, NJ 08540, USA
Gregory Moore
Department of Physics, Rutgers University
Piscataway, New Jersey, 08855-0849
Edward Witten
Dept. of Physics, Caltech, Pasadena, CA 91125
and CIT-USC Center for Theoretical Physics, USC, Los Angeles CA
We show how some aspects of the $`K`$-theory classification of RR fluxes follow from a careful analysis of the phase of the $`M`$-theory action. This is a shortened and simplified companion paper to โ$`E_8`$ Gauge Theory, and a Derivation of $`K`$-Theory from $`M`$-Theory.โ
On leave from Institute for Advanced Study, Princeton, NJ 08540.
May 3, 2000
1. Introduction
In the past few years, we have learned that $`D`$-brane charges should be thought of in the framework of $`K`$-theory \[1,,2,,3\]. More recently, it has been realized that the topological classification of RR fluxes in Type II string theory is also $`K`$-theoretic \[4,,5,,6\]. Other developments of the past few years, such as Matrix Theory, and the AdS/CFT correspondence, have shown that $`D`$-branes play an important role in the search for a more fundamental formulation of $`M`$-theory. It is natural, therefore, to ask how the $`K`$-theoretic formulation of RR charges and fluxes can be formulated in terms of $`M`$-theory.
In hindsight, the $`K`$-theoretic interpretation of RR fluxes and charges is almost inevitable, given the existence of Chan-Paton vector bundles on $`D`$-branes. But $`M`$-branes do not carry Chan-Paton bundles or vector fields, so the $`M`$-theoretic origin of $`K`$-theory is not manifest. In this letter, we will outline how some aspects of the $`K`$-theoretic formulation of RR charges and fluxes can in fact be derived from $`M`$-theory. In brief, when $`M`$-theory is formulated on an eleven-manifold of the form $`Y=S^1\times X`$, we can derive what might be called โintegral equations of motionโ for the topological class of the four-form flux of $`M`$-theory. These equations state that (on-shell) the four-form flux of $`M`$-theory is in fact $`K`$-theoretic, in a sense we will make more precise below. In order to keep technical complications to a minimum, we will make several simplifying assumptions on the topology of $`X`$. Further details, without the simplifying assumptions, can be found in .
Let $`X`$ be a compact spin 10-manifold with a fixed Riemannian metric $`g_{\mu \nu }`$. Consider Type IIA superstring theory on $`X`$ with metric $`tg_{\mu \nu }`$. We will study the partition function of the theory in the limit $`g_s0`$, where $`g_s`$ is the string coupling, and then $`t\mathrm{}`$. We will then consider $`M`$-theory on $`Y=X\times S^1`$, with metric
$$ds_M^2=g_s^{4/3}(dx^{11})^2+tg_s^{2/3}g_{\mu \nu }dx^\mu dx^\nu .$$
In $`M`$-theory, we will study the partition function in the limit $`t\mathrm{}`$ and then $`g_s0`$. Finally, we will show that the leading terms in the partition function of $`M`$ theory and Type IIA theory are in agreement.
One might at first think the agreement between the two expressions is trivial, since eleven-dimensional supergravity reduces on a circle to ten-dimensional supergravity. However, things are not so simple because the $`K`$-theoretic nature of RR fluxes implies that the sum over RR field-strengths is not simply a sum over all harmonic forms that obey conventional Dirac quantization. The quantization law is more subtle, and, as we will see, there are subtle phases in the action which are not present in the standard treatments of supergravity Lagrangians. The goal will be to derive these subtleties of the Type IIA theory starting from $`M`$-theory. In section 6, we will describe three results, of independent interest, which are corollaries of our analysis.
2. The Type IIA Partition Function
In the weak coupling limit described above, we consider the NS fields fixed at some classical (not necessarily on-shell) values. In this background, we will study the one-loop quantum mechanics of fermions and the free-field quantum mechanics of the RR fields. The partition function is accordingly
$$Z_{IIA}\mathrm{exp}\left(\frac{1}{g_s^2}S_{NS}\right)\frac{\mathrm{\Theta }_{IIA}}{\mathrm{\Delta }}$$
where $`S_{NS}`$ is the action for the NS-sector fields, $`\mathrm{\Delta }`$ is a product of bose and fermi determinants, and $`\mathrm{\Theta }_{IIA}`$, which is the factor that we will really focus on, is a theta function which arises from summing over the fluxes of the RR $`p`$-form fields.
A more complete treatment of the problem would involve integrating over the moduli of the NS fields โ notably the metric and the $`B`$-field โ but in the present discussion we hold these fixed and in particular set $`B=0`$.
A complete description of how to construct $`\mathrm{\Theta }_{IIA}`$ can be found in \[8,,4,,5\]. We consider the full RR field-strength $`G=G_0+G_2+\mathrm{}+G_{10}`$. This is a sum of real differential forms on $`X`$ of even degree. In IIA supergravity, there is a Bianchi identity $`dG=0`$, so $`G`$ has a characteristic cohomology class that one can regard as an element of the even degree cohomology $`H^{\mathrm{even}}(X;๐)`$. Because $`D`$-branes exist, $`G`$ obeys a Dirac-like quantization condition, but this condition does not merely state, as one might guess, that the periods of $`G/2\pi `$ are integral or in other words that $`G/2\pi `$ is associated with an element of the integral cohomology $`H^{\mathrm{even}}(X;๐)`$. Instead, from the existence of $`D`$-branes and their couplings to RR fields, one can deduce \[4,,5,,6\] that the topological sectors of RR fields are classified by an element $`xK(X)`$ of the $`K`$-theory of $`X`$. (In Type IIB one would use $`xK^1(X)`$ in a similar way.) For given $`x`$, the cohomology class of $`G`$ is
$$G(x)=\mathrm{ch}(x)\sqrt{\widehat{A}(X)}.$$
This is the $`K`$-theory version of Dirac quantization.
For our purposes, to construct the partition function $`\mathrm{\Theta }_{IIA}`$, we will be summing over on shell RR fields. For each $`xK(X)`$, there is a unique harmonic form in the cohomology class of $`G(x)`$; we use this in the sum over RR fields.
The second subtlety is that the total field-strength $`G`$ should be considered to be self-dual, $`G=G`$, an equation which at the classical level is hardly compatible with (2.1) for general metrics. The resolution of this paradox is that \[8,,4\] one interprets $`G=G`$ as a statement in the quantum theory. One sums over โhalfโ the fluxes, and these are indeed quantized by (2.1). More precisely $`\mathrm{\Gamma }=K(X)/K(X)_{tors}`$ (where $`K(X)_{tors}`$ is the torsion subgroup) is a lattice with an integral symplectic form given by the index of the Dirac operator coupled to $`x\overline{y}`$. If we let $`I(z)`$ denote the index of the Dirac operator coupled to a bundle $`z`$, and we use the Atiyah-Singer formula for the index, then the symplectic form $`\omega (x,y)`$ is defined by
$$\omega (x,y)=I(x\overline{y})=_X\mathrm{ch}(x\overline{y})\widehat{A}.$$
Poincarรฉ duality in $`K`$-theory implies that this form is unimodular. One sums over โhalf the fluxesโ by summing over fluxes associated with $`K`$-theory classes in a maximal Lagrangian sublattice $`\mathrm{\Gamma }_1\mathrm{\Gamma }`$. This sum is the theta function for the quantum self-dual RR field.
To sketch in somewhat more detail the definition of the theta function, we use the symplectic structure (2.1) to give the torus $`๐ซ_K(X)=(K(X)๐)/\mathrm{\Gamma }`$ the structure of a compact phase space. Moreover, this phase space has a metric
$$x^2=G(x)G(x).$$
There is therefore a unique translation invariant complex structure $`J`$ on $`๐ซ_K(X)`$ such that the metric (2.1) is of type $`(1,1)`$. Explicitly, $`G_{2p}(Jx)=(1)^{p+1}(G_{102p}(x))`$. Coherent state quantization with respect to this complex structure leads to a unique quantum state, since the symplectic volume of $`๐ซ_K(X)`$ is one. $`\mathrm{\Theta }_{IIA}`$ is the wavefunction of this quantum state. To write it more explicitly, we now choose a complementary Lagrangian sublattice $`\mathrm{\Gamma }_2`$ so that $`\mathrm{\Gamma }=\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$. The lattice vectors in $`\mathrm{\Gamma }_1,\mathrm{\Gamma }_2`$ define โ$`a`$-cyclesโ and โ$`b`$-cyclesโ in $`๐ซ_K(X)`$, and with respect to this decomposition we have a period matrix $`\tau `$, which is a quadratic form on $`\mathrm{\Gamma }_1๐`$ with positive imaginary part. Finally, we must define the characteristics of the theta function. This is the subtlest part of the quantization procedure. We introduce a function $`\mathrm{\Omega }:\mathrm{\Gamma }๐_2`$ such that
$$\mathrm{\Omega }(x+y)=\mathrm{\Omega }(x)\mathrm{\Omega }(y)e^{i\pi \omega (x,y)}.$$
Then, we may define the characteristics $`\theta \mathrm{\Gamma }_1,\varphi \mathrm{\Gamma }_2`$ by
$$\begin{array}{cc}\hfill \mathrm{\Omega }(x)=(1)^{\omega (x,\varphi )}\mathrm{for}x\mathrm{\Gamma }_1& \mathrm{\Omega }(x)=(1)^{\omega (x,\theta )}\mathrm{for}x\mathrm{\Gamma }_2.\hfill \end{array}$$
For an explanation of why the $`\mathrm{\Omega }`$ function is needed and the rationale for the definition of the characteristics, see \[8,,4\]. Finally, we may write the explicit formula for the theta function:
$$\mathrm{\Theta }_{IIA}=e^{i\pi \mathrm{Re}(\tau (\theta /2))}\underset{x\mathrm{\Gamma }_1}{}e^{i\pi \tau (x+\frac{1}{2}\theta )}\mathrm{\Omega }(x).$$
It remains to identify $`\mathrm{\Omega }`$. There is a (presumably unique) $`T`$-duality invariant choice \[8,,4\] given in terms of the mod two index of Atiyah and Singer . If $`V`$ is a real vector bundle on $`X`$, then we define $`q(V)`$ to be the number, modulo two, of chiral zero modes of the Dirac operator $`\text{/}D_V`$ coupled to $`V`$. For $`X`$ of dimension $`8k+2`$, $`q(V)`$ is a topological invariant. The definition of is:
$$\mathrm{\Omega }(x)=(1)^{q(x\overline{x})}.$$
It is often useful to regard $`x`$ as the charge of a $`D`$-brane in Type IIB theory. Then, by a Born-Oppenheimer argument , $`q(x\overline{x})`$ counts the number modulo two of fermion zero modes in the Ramond sector for open strings with boundary conditions defined by $`x`$ at each end. This makes the $`T`$-duality invariance of $`\mathrm{\Omega }`$ manifest.
A few facts about the mod two index will prove useful below. In general, the mod two index is not just the reduction modulo two of an ordinary index (which is, after all, simply zero in 10 dimensions for $`V`$ real). It is true that if the real bundle $`V`$ can (after complexification) be written as $`V=x\overline{x}`$ where $`xK(X)`$ and $`\overline{x}`$ is the complex conjugate of $`x`$, then $`q(V)`$ equals the mod two reduction of $`I(x)`$. This fact is used in showing that $`\mathrm{\Omega }`$ satisfies the cocycle relation (2.1).
There are many different choices of sublattices $`\mathrm{\Gamma }_1`$. Up to an overall normalization, different choices lead to different descriptions of the same partition function. In the problem discussed in this letter, there is a very natural choice of polarization. To motivate it, consider the behavior of the kinetic energy of a non-self-dual field $`G`$ as we scale the metric $`g_{\mu \nu }tg_{\mu \nu }`$. The kinetic energies scale as:
$$t^5G_0^2+t^3G_2^2+tG_4^2+t^1G_6^2+t^3G_8^2+t^5G_{10}^2.$$
We would like to choose a polarization so that only positive powers of $`t`$ show up in the exponential. Otherwise the sum over fluxes becomes less and less convergent as $`t\mathrm{}`$, and the terms in the sum do not accurately reflect the long-distance physics. We will have only positive powers of $`t`$ if we take $`\mathrm{\Gamma }_2`$ to be the set of $`K`$-theory classes $`x`$ with $`\mathrm{ch}_0(x)=\mathrm{ch}_1(x)=\mathrm{ch}_2(x)=0`$, and then take $`\mathrm{\Gamma }_1`$ to be a complementary Lagrangian sublattice. Working through the definitions of the quantization procedure one finds:
$$\begin{array}{cc}\hfill \mathrm{\Theta }_{IIA}=e^{i\pi \mathrm{Re}(\tau (\theta /2))}\underset{x\mathrm{\Gamma }_1}{}& e^{\pi (t^5G_0^2+t^3G_2^2+tG_4^2)}\hfill \\ & e^{i\pi {\scriptscriptstyle (G_0G_{10}G_2G_8+G_4G_6)}}\mathrm{\Omega }(x)\hfill \end{array}$$
where $`G`$ is understood to be given by (2.1) evaluated for $`x+\frac{1}{2}\theta `$. It might look at first sight like this is just the standard recipe for computing the RR partition function as a sum over fluxes of $`G_0`$, $`G_2`$, and $`G_4`$, with the higher RR fields eliminated using self-duality. However, the allowed values of the $`G_0,G_2`$, and $`G_4`$ fluxes differ from what one would conventionally guess. Moreover, in addition to a factor from the standard kinetic energy of $`G_0`$, $`G_2`$, and $`G_4`$, the action contains nonstandard phase factors. The factor $`e^{i\pi {\scriptscriptstyle (G_0G_{10}G_2G_8+G_4G_6)}}`$ (which arises by computing the real part of the $`\tau `$ function of the lattice) is, after imposing (2.1) (which constrains the $`G_{2p}`$ of $`p>2`$ in terms of those of $`p2`$) a $`120^{th}`$ root of unity that is given by a complicated cohomological formula and is not part of the standard supergravity formalism. The sign factor $`\mathrm{\Omega }(x)`$ is not given by any cohomological formula.
As we have stressed in the introduction, we want to focus on the behavior for $`t\mathrm{}`$. The dominant contributions come from $`K`$-theory classes $`x\mathrm{\Gamma }_1`$ such that $`G_0(x)=G_2(x)=0`$. A glance at (2.1) shows that these are classes of virtual dimension zero such that $`c_1(x)=0`$. Denoting by $`\mathrm{\Gamma }_1^{}`$ the sublattice of such classes, the leading term in the partition function may be simplified to
$$\mathrm{\Theta }_{IIA}=e^{i\pi \mathrm{Re}(\tau (\theta /2))}\underset{x\mathrm{\Gamma }_1^{}}{}e^{\pi tG_4^2}e^{i\pi {\scriptscriptstyle G_4G_6}}\mathrm{\Omega }(x).$$
It is also important to include $`G_2`$ for a more complete comparison to $`M`$-theory ($`G_0`$ has no known origin in $`M`$-theory, at least for general backgrounds), but for simplicity, in the present letter we consider only $`G_4`$.
3. The $`M`$-Theory Partition Function
The partition function for $`M`$-theory in the large volume limit is given by
$$Z_M\mathrm{exp}\left(_Y\sqrt{g}\right)\frac{1}{\mathrm{\Delta }_M}\mathrm{\Theta }_M$$
where the leading term is the Einstein action of a fixed Riemannian metric (1.1) on a spin 11-manifold $`Y`$, $`\mathrm{\Delta }_M`$ are one-loop determinants (which we take to be positive, absorbing the sign in $`\mathrm{\Theta }_M`$), and $`\mathrm{\Theta }_M`$ is the sum over the classical on-shell configurations of the $`C`$-field.
As in Type IIA, there is a subtle quantization law on $`G=dC`$ as well as a subtle phase-factor in the path integral . The topological quantization of 4-form field-strengths is given by choosing any element $`aH^4(Y;๐)`$ and taking $`G(a)`$ to be a certain de Rham representative satisfying
$$\frac{G(a)}{2\pi }=a\frac{1}{2}\lambda .$$
Here $`\lambda `$ is the degree four class represented at the level of differential forms by $`\frac{\mathrm{tr}RR}{16\pi ^2}`$. It is an integral class on a spin manifold and satisfies $`2\lambda =p_1`$. The contribution of a field $`C`$ in the topological sector $`a`$ to the $`M`$-theory partition function is
$$e^{G(a)^2}\mathrm{\Omega }_M(C),$$
where $`G(a)`$ is the on-shell field configuration.
The phase factor $`\mathrm{\Omega }_M(C)`$ is a globally well-defined version of the familiar supergravity interaction $`_YCGG+\mathrm{}`$. Since $`G(a)`$ is a nontrivial cohomology class, $`C`$ is not globally well-defined as a three-form, and the proper formulation of the phase is tricky . We first find a 12-manifold $`Z`$ such that $`Z=Y`$ and $`a`$ extends to $`\stackrel{~}{a}H^4(Z;๐)`$. The existence of the pair $`(Z,\stackrel{~}{a})`$ is highly nontrivial, but guaranteed by a result of Stong . We then define the phase by
$$\mathrm{\Omega }_M(C)=ฯต\mathrm{exp}\left[2\pi i_Z\left(\frac{1}{6}(\stackrel{~}{a}\frac{1}{2}\lambda )^3+(\stackrel{~}{a}\frac{1}{2}\lambda )\frac{\lambda ^2p_2}{48}\right)\right],$$
where $`ฯต`$ is the sign of the Pfaffian of the gravitino operator. In a topologically trivial situation we may identify $`\stackrel{~}{a}\frac{1}{2}\lambda =\stackrel{~}{G}=d\stackrel{~}{C}`$ and apply Stokesโ theorem to make contact with the more standard supergravity expressions.
The expression (3.1) is not manifestly well-defined since the choice of $`Z`$ is not unique. It was shown in that this difficulty is most elegantly eliminated by using $`E_8`$ index theory. We will also find the connection to $`E_8`$ useful below. Therefore, let us recall that on manifolds $``$ of dimension less than 16, cohomology classes $`bH^4(;๐)`$ are in 1-1 correspondence with topological classes of $`E_8`$ vector bundles $`V(b)`$ on $``$ . We hence consider the $`E_8`$ bundle $`V(\stackrel{~}{a})`$ on $`Z`$ in the adjoint representation and choose a connection $`A`$ on $`V(\stackrel{~}{a})`$ such that on $`Y`$
$$C=\frac{1}{16\pi ^2}\frac{1}{30}\mathrm{Tr}_{\mathrm{๐๐๐}}(AdA+\frac{2}{3}A^3)+\frac{1}{32\pi ^2}\mathrm{Tr}(\omega d\omega +\frac{2}{3}\omega ^3).$$
In other words, we interpret $`C`$ as a Chern-Simons three form of $`E_8`$ gauge theory plus gravity. (It is not obvious that given a $`C`$-field a corresponding gauge field $`A`$ always exists. A slightly longer argument can be made if a connection $`A`$ making (3.1) hold on the nose does not exist.) We then can evaluate the phase in terms of the $`\eta `$ invariants of the Dirac operator $`\text{/}D_{V(a)}`$ and the gauge-fixed Rarita-Schwinger operator $`D_{RS}=\text{/}D_{TX3๐ช}`$ on $`Y`$. Here $`๐ช`$ is the trivial real line bundle. Using the APS index theorem , and the fact that the index of $`\text{/}D_{V(a)}`$ is even in 12 dimensions, we can rewrite the phase as
$$\mathrm{\Omega }_M(C)=\mathrm{exp}\left[2\pi i\left(\frac{\eta (\text{/}D_{V(a)})+h(\text{/}D_{V(a)})}{4}+\frac{\eta (D_{RS})+h(D_{RS})}{8}\right)\right].$$
Here $`h`$ denotes the number of zero modes of the operator in question on $`Y`$. (The sign $`ฯต`$ in $`\mathrm{\Omega }_M`$ is absent as it cancels against a term that comes from the APS theorem.)
In order to compare to IIA superstrings, we will now restrict attention to $`Y=X\times ๐^1`$. We restrict to fields invariant under rotations of $`๐^1`$, and, since we have taken $`B=0`$ on the Type IIA side, we assume that the $`C`$-field is a pullback from $`X`$. Under these conditions, $`\mathrm{\Omega }_M(C)`$ is a topological invariant that depends only on the characteristic class $`a`$, and not on $`C`$, so we will denote it as $`\mathrm{\Omega }_M(a)`$. Moreover, the $`\eta `$ invariants vanish for $`X\times ๐^1`$ because of the reflection symmetry of $`๐^1`$, and only the contributions from the number $`h`$ of zero modes survive. The $`a`$-dependent factor is then simply
$$\mathrm{\Omega }_M(a)=\mathrm{exp}[i\pi h(\text{/}D_{V(a)})/2].$$
Using the standard relation between the radius $`R`$ of $`S^1`$ and the IIA string coupling we find:
$$\mathrm{\Theta }_M=\underset{aH^4(X;๐)}{}e^{G(a)^2}\mathrm{\Omega }_M(a)+๐ช(e^{1/g_s^2}).$$
where $`G(a)`$ is the harmonic form in the cohomology class $`a\frac{1}{2}\lambda `$ and the corrections correspond to field-strengths which have an index tangent to the $`M`$-theory circle or are not invariant under rotations of the circle. <sup>1</sup> In fact, the supergravity equation is $`\kappa dG=\frac{1}{2}GG+(\lambda ^2p_2)/48`$ for a suitable constant $`\kappa `$. The components of this equation which are pulled back from $`X`$ are enforced by subsequent integration over the $`B`$-field, which is held fixed (and zero) in this letter. The difference from the harmonic form is cohomologically trivial.
4. $`K`$-Theory vs. Cohomology
We would now like to compare (2.1) to (3.1). These are, a priori, rather different expressions. One involves a sum over a certain part of $`K(X)`$ that is related to $`H^4(X;๐)`$ and the other involves a sum over $`H^4(X;๐)`$. If we included the other RR fields, then on the Type IIA side, we would be summing over the lattice $`\mathrm{\Gamma }_1K(X)`$, and on the $`M`$-theory side we would be summing over degree two and degree four cohomology classes of $`X`$.
In general, $`K(X)`$ and $`H^{\mathrm{even}}(X;๐)`$ (the sum of the even degree integral cohomology groups of $`X`$) are closely related groups. If one tensors with the real numbers, they become isomorphic (by the map that takes an element of $`K(X)`$ to its Chern character). However the integral structures (which determine the Dirac quantization conditions) are different, and the torsion subgroups can be very different. For example, $`H^{\mathrm{even}}(RP^{2n+1};๐)๐๐_2^n`$ while $`K(RP^{2n+1};๐)๐๐_{2^n}`$. We now describe with some more precision the relation of $`K`$ and $`H^{\mathrm{even}}`$ at the integral level. The reader will find more detail in .
Let us first describe how integral cohomology arises from a $`K`$-theory class $`x`$. For every $`K`$-theory class $`x`$, there is a smallest integer $`i`$ such that $`x`$ can be represented as the class of a $`(2i1)`$-brane, wrapped on a $`2i`$-dimensional submanifold $`Q`$ of $`X`$. The Poincarรฉ dual of $`Q`$ is a $`(102i)`$-dimensional cohomology class $`\alpha `$ associated with $`x`$.
The map from a $`K`$-theory class $`x`$ to an associated cohomology class $`\alpha (x)`$ is the first step in a systematic procedure, known as the Atiyah-Hirzebruch spectral sequence, for comparing $`K`$-theory to cohomology. This map is relevant to us, because the Type IIA formula (2.1) involves a sum over $`K`$-theory classes $`x`$ for which $`\alpha (x)`$ is an element of $`H^4(X;๐)`$ (modulo those for which $`\alpha (x)`$ is of degree six or higher), while the $`M`$-theory formula (3.1) is a sum over the characteristic class $`aH^4(X;๐)`$. We will compare the $`M`$-theory sum over $`a`$ to the Type IIA sum over $`\alpha (x)`$.
Is it the case that every $`aH^4(X;๐)`$ is $`\alpha (x)`$ for some $`xK(X)`$? The answer to this question is โno.โ In ten dimensions, a necessary and sufficient condition for $`x`$ to exist is that
$$Sq^3(a)=0,$$
where $`Sq^3`$ is a certain cohomology operation, known as a Steenrod square. If $`x`$ exists, we call it a โ$`K`$-theory liftโ of $`a`$. Such an $`x`$ has virtual dimension zero, $`c_1(x)=0`$, and $`c_2(x)=a`$.
An introduction to the Steenrod squares $`Sq^i`$ is given in . In brief, if $`aH^k(X;๐)`$ then $`Sq^3(a)H^{k+3}(X;๐)`$ may be defined as follows. Let $`Q(a)`$ be a submanifold that is Poincarรฉ dual to $`a`$ in $`X`$. Then the normal bundle $`N(Q)`$ of $`Q`$ has integral characteristic classes $`W_i(N(Q))H^i(Q;๐)`$ for $`i`$ odd. These can be pushed into a tubular neighborhood of $`Q`$, allowing us to define
$$Sq^3(a)=W_3(N(Q))a.$$
Similarly, one defines the mod two Steenrod squares for all $`i`$ (not necessarily odd) as follows. For $`\overline{a}H^k(X;๐/2)`$, we set $`Sq^i(\overline{a})=w_i(N(Q))\overline{a}H^{k+i}(X;๐/2)`$. The Steenrod squares obey many identities; the ones we need are as follows (for sketches of proofs see ). First, the integral and mod two squares are related. In fact, $`Sq^3(a)=0`$ if and only if $`Sq^2(a)`$ has an integral lift, that is, if and only if there is an integral class $`b`$ whose mod two reduction is $`Sq^2(a)`$. Second, it is possible to โintegrate by partsโ with Steenrod squares. That is, for any $`a,b`$, $`_XaSq^2b=_XSq^2ab`$. Closely related to the criterion (4.1) for a $`K`$-theory lift to exist is the fact that if $`x`$ is a $`K`$-theory class with $`c_1(x)=0`$, then the higher Chern classes of $`x`$ obey
$$c_3(x)=Sq^2c_2(x)\mathrm{mod}2.$$
Finally, $`Sq^3Sq^3=0`$, so we may take its cohomology, namely, the kernel of $`Sq^3`$ acting on $`H^{\mathrm{even}}(X,๐)`$ modulo the image of $`Sq^3`$ acting on $`H^{\mathrm{odd}}(X,๐)`$. In our situation, the cohomology of $`Sq^3`$ is, essentially, $`K(X)`$.
While the equation $`Sq^3(a)=0`$ might seem somewhat exotic, it is a close cousin of a condition that has already appeared in the physics literature on $`D`$-branes. In particular, if we think of $`x`$ as determining the $`D`$-brane charge (in IIB string theory) of a brane wrapped on $`Q`$, then cancellation of worldsheet global anomalies implies that $`W_3(N(Q))=0`$ . Thus, by (4.1) it follows that $`Sq^3(a)=0`$ if $`a`$ is Poincarรฉ dual to $`Q`$.
Let us now apply these remarks to study (2.1). We would like to convert the sum over the sublattice $`\mathrm{\Gamma }_1^{}`$ in $`K`$-theory to a sum over cohomology elements, namely to a sum over classes $`aH^4(X;๐)`$ such that $`a`$ has a $`K`$-theory lift. At this point we run into an apparent difficulty. A $`K`$-theory lift of $`a`$, if it exists, is not unique because given one lift $`x`$ it is always possible to add a class $`y`$ to $`x`$ where $`y`$ is any element of the lattice $`\mathrm{\Gamma }_2`$ introduced earlier (thus $`G_{2p}(y)=0`$ for $`p2`$). The quantities $`G_6(x)`$ and $`\mathrm{\Omega }(x)`$ in (2.1) definitely depend on the choice of $`K`$-theory lift, but the product
$$e^{i\pi {\scriptscriptstyle G_4(x+{\scriptscriptstyle \frac{1}{2}}\theta )G_6(x+{\scriptscriptstyle \frac{1}{2}}\theta )}}\mathrm{\Omega }(x)$$
does not. This can be demonstrated using the facts noted above. Using (2.1), (2.1), and integration by parts, we first rewrite (4.1) as
$$\mathrm{exp}[i\frac{\pi }{4}c_2(\theta )c_3(\theta )]\mathrm{exp}\left[i\frac{\pi }{2}(c_2(x)+c_2(\theta ))c_3(x)\right]\mathrm{\Omega }(x).$$
If we change the $`K`$-theory lift by $`xx+y`$ with $`y`$ as above then using the cocycle condition (2.1), the definition (2.1) of $`\theta `$, and the index theorem one can verify that (4.1) is unchanged. At this point, we have shown that the Type IIA partition function can be written as a sum over $`G_4`$ fluxes (along with $`G_2`$ and $`G_0`$ if one chooses to include them), as one would naively expect, but with subtle shifts in the Dirac quantization condition and an exotic sign factor in the sum over fluxes.
Finally, it remains to translate the characteristic $`\theta `$ into cohomology. Again, it is useful to regard $`x`$ as the $`D`$-brane charge of a brane in IIB theory wrapping some worldvolume $`Q`$ of smallest possible dimension. As we noted above, the mod two index $`q(x\overline{x})`$ is given by the number mod two of the fermion zero modes of a singly-wrapped brane on $`Q`$. The classes $`x\mathrm{\Gamma }_2`$ correspond to $`D(1)`$, $`D1`$ and $`D3`$ instantons in $`X`$. In the first two cases, the number of zero modes is easily seen to be even. On the other hand, for a $`D3`$ instanton, the number of fermion zero modes is given by the index theorem to be $`_Q\lambda \mathrm{mod}2`$. We conclude that $`\mathrm{ch}(\theta )=\lambda +\mathrm{}`$. Thus, we can simplify (2.1) to
$$\mathrm{\Theta }_{IIA}=\underset{aH^4(X;๐):Sq^3(a)=0}{}e^{\pi ta\frac{1}{2}\lambda ^2}e^{i\frac{\pi }{2}{\scriptscriptstyle (c_2(x(a))+\lambda )c_3(x(a))}}\mathrm{\Omega }(x(a))$$
where $`x(a)`$ is any $`K`$-theory lift of $`a`$. We have now expressed the $`K`$-theory sum in terms of cohomology. It is now time to re-examine the $`M`$-theory sum (3.1).
5. The Integral Equation of Motion in $`M`$-Theory
In the previous section we reduced the $`K`$-theory partition function to a sum over a subgroup of $`H^4(X;๐)`$. This subgroup is of finite index, since for any $`a`$, $`Sq^3(a)`$ is of order 2, and hence $`Sq^3(2a)=0`$. By contrast, the $`M`$-theory partition function is a sum over the full group $`H^4(X;๐)`$. To show that the two expressions for the partition functions agree, we will argue that the $`M`$-theory phase $`\mathrm{\Omega }_M(a)`$ leads to an โintegral equation of motionโ $`Sq^3(a)=0`$ on the topological sectors in $`M`$-theory.
In this letter we will, for simplicity, show agreement of (4.1) and (3.1) under the assumption that $`\mathrm{\Omega }_M(c)=1`$ for torsion $`c`$, and that $`Sq^3(c)=0`$ for all torsion elements $`c`$. Suppose $`aH^4(X;๐)`$ is any class and $`cH^4(X;๐)_{tors}`$. The kinetic energy of $`G`$ in the topological sector $`a+c`$ is identical to that in the sector $`a`$ because the field-strength defined by (3.1) is a real differential form and hence $`G(a+c)=G(a)`$. Since the torsion subgroup is finite we may equally well write (3.1) as
$$\mathrm{\Theta }_M=\underset{aH^4(X;๐)}{}e^{G(a)^2}\mathrm{\Omega }_M^{\mathrm{av}}(a)$$
with
$$\mathrm{\Omega }_M^{\mathrm{av}}(a)=\frac{1}{|H^4(X;๐)_{tors}|}\underset{cH^4(X;๐)_{tors}}{}\mathrm{\Omega }_M(a+c).$$
This is useful because the $`M`$-theory phase $`\mathrm{\Omega }_M(a)`$ is not independent of the shift $`aa+c`$. Indeed, the bundles $`V(a+c)`$ and $`V(a)`$ are definitely not isomorphic, and, as we will demonstrate below, $`\mathrm{\Omega }_M^{\mathrm{av}}`$ is in fact a projection operator. Under a simplifying topological assumption (described below) this operator is:
$$\begin{array}{cc}\hfill \mathrm{\Omega }_M^{\mathrm{av}}(a)& =0\mathrm{if}Sq^3(a)0\hfill \\ & =\mathrm{\Omega }_M(a)\mathrm{if}Sq^3(a)=0.\hfill \end{array}$$
Moreover, when $`Sq^3(a)=0`$, so that $`a`$ has a $`K`$-theory lift $`xK(X)`$, we can compare $`M`$-theory and $`K`$-theory phases. We will show that they agree
$$\mathrm{\Omega }_M(a)=\mathrm{\Omega }(x)e^{\frac{i\pi }{2}_X(c_2(x)+\lambda )c_3(x)}.$$
The agreement of (4.1) with (3.1) immediately follows from the above pair of results. We will now sketch how they are derived, beginning with the proof of (5.1). It is here that the interpretation of the $`M`$-theory phase in terms of $`E_8`$ gauge theory is particularly effective. We are interested in $`Y=X\times S^1`$ with supersymmetric spin structure on the $`M`$-theory circle. In evaluating (3.1), we use the fact that a zero mode of $`\text{/}D_{V(a)}`$ is constant along the $`M`$-circle so that the phase just depends on the number of zero modes on $`X`$. These may be expressed in terms of the number of chiral zero modes in 10 dimensions via $`h(\text{/}D_{\pi ^{}V(a)})=h^+(\text{/}D_{V(a)})+h^{}(\text{/}D_{V(a)})=2h^+(\text{/}D_{V(a)})`$. We conclude that the phase is expressed in terms of a mod two index, $`f(a)=q(V(a))`$:
$$\mathrm{\Omega }_M(a)=(1)^{f(a)}.$$
The next step is to relate the $`E_8`$ bundle $`V(a)`$ to a $`K`$-theory class $`x`$. In general, $`K`$-theory classes are differences of vector bundles $`x=E_1E_2`$ where the structure group of $`E_i`$ must be taken to be $`U(N)`$ for some large $`N`$. However, we are working in 10-dimensions, and this dimension is sufficiently small that all $`K`$-theory classes on $`X`$ with $`c_1(x)=0`$ can be realized using $`SU(5)`$ bundles. The reason for this is that the classification of bundles on a ten-manifold depends only on the homotopy groups $`\pi _i(SU(N))`$ for $`i<10`$. (See for a description of this approach.) In 10 dimensions, $`SU(5)`$ is in the stable range: $`\pi _i(SU(5))=\pi _i(SU(\mathrm{}))`$, $`i<10`$. We can therefore take our $`K`$-theory lift to be $`x=EF`$ where $`F`$ is a trivial rank 5 bundle and $`E`$ is an $`SU(5)`$ bundle with $`\mathrm{ch}(E)=5+a+\mathrm{}`$. We can construct an $`E_8`$ bundle $`V(a)`$ with characteristic class $`a`$ using the โembeddingโ of $`SU(5)\times SU(5)`$ in $`E_8`$, taking the two $`SU(5)`$ bundles to be respectively $`E`$ and $`F`$. Using the decomposition of the adjoint representation of $`E_8`$ under $`SU(5)\times SU(5)`$ and the fact that $`F`$ is trivial, one finds that for mod 2 index theory (throwing away representations that appear an even number of times), the $`\mathrm{๐๐๐}`$ is equivalent to $`E\overline{E}๐ช+^2E+^2\overline{E}`$, where $`๐ช`$ is a trivial line bundle and $`^2`$ denotes the second antisymmetric product. Using the properties of the mod 2 index described above we now learn that
$$q(V(a))=q(E\overline{E}๐ช)+I(\mathrm{\Lambda }^2(E))=q(x\overline{x})+I(\mathrm{\Lambda }^2(E))+I(E)\mathrm{mod}2.$$
The formula (5.1) leads directly to (5.1). Indeed, the first term on the RHS of (5.1) corresponds to $`\mathrm{\Omega }(x)`$, while by an easy application of the index theorem, the second term is $`\frac{1}{2}(c_2(x)+\lambda )c_3(x)\mathrm{mod}2`$.
It remains to show (5.1). This is based on an analog of (2.1) for the $`E_8`$ mod two index $`f(a)`$. Namely, $`f`$ satisfies the bilinear identity
$$f(a+a^{})=f(a)+f(a^{})+_XaSq^2a^{}.$$
Unfortunately, there does not appear to be an elementary proof of (5.1). A proof using โcobordism theoryโ can be found in section 3.2 of . Granted this, we are now ready to complete the proof of (5.1). The argument simplifies considerably if we assume that $`f(c)=f_0`$ is independent of $`c`$ for torsion classes $`c`$. In this case we can write
$$\mathrm{\Omega }_M^{\mathrm{av}}(a)=(1)^{f(a)+f_0}\frac{1}{|H^4(X;๐)_{tors}|}\underset{cH^4(X;๐)_{tors}}{}e^{i\pi {\scriptscriptstyle c}Sq^2a}.$$
Now, for any $`bH^6(X;๐)`$ and any torsion $`c`$ it is always true that $`bc=\frac{1}{n}b(nc)=0`$. Using Poincarรฉ duality, one can prove the converse: if $`\overline{b}H^6(X;๐_2)`$ satisfies $`c\overline{b}=0`$ for all $`cH^4(X;๐)_{tors}`$, then $`\overline{b}`$ is the reduction of some integral class. Therefore $`\mathrm{\Omega }_M^{\mathrm{av}}(a)`$ projects onto the set of classes $`a`$ such that $`Sq^2a`$ has an integral lift. This is equivalent to $`Sq^3(a)=0`$, i.e., to the statement that $`a`$ has a $`K`$-theory lift $`x`$. Indeed $`Sq^2a`$ is the reduction modulo two of $`c_3(x)`$. This completes the proof of (5.1), and therefore establishes the equivalence of (2.1) and (3.1).
6. Three applications
The $`K`$-theory/$`E_8`$-formalism described above leads to three interesting physical effects which we will sketch very briefly here.
First, an easy consequence of (5.1) leads to a new topological consistency condition on string backgrounds. By (5.1) we have $`f(a+2c)=f(a)+f(2c)`$, and moreover $`f(2c)=cSq^2c`$. By a result of Stong $`cSq^2c=cSq^2\lambda `$. Now, the reasoning below (5.1) shows that $`\mathrm{\Omega }_M^{\mathrm{av}}(a)=0`$, and hence $`\mathrm{\Theta }_M=0`$ if $`Sq^3\lambda 0`$. In algebraic topology one shows that a certain characteristic class $`W_7(X)`$ of $`X`$ is $`Sq^3(\lambda )`$. Thus $`W_7(X)=0`$ is a necessary condition for a consistent background. Unfortunately, we do not know an intuitive interpretation of this condition.
Second, it turns out that the parity symmetry of $`M`$ theory on $`X\times S^1`$ (coming from reflection of the $`S^1`$) is anomaly-free, but this depends on a surprising anomaly cancellation between bosons and fermions. In IIA theory, this symmetry is $`(1)^{F_L}`$, and in IIB it is related to strong/weak coupling duality. By counting fermion zero-modes one can show that the gravitino measure $`\mu `$ transforms under parity as $`\mu (1)^{q(TX)}\mu `$. The mod-two index $`q(TX)`$ is nonvanishing for certain 10-folds, such as $`X=T^2\times \mathrm{๐๐}^2`$. The fermion anomaly is cancelled by the nontrivial transformation law of the $`G`$-flux partition function (3.1). Indeed, parity acts as $`GG`$ and by (3.1) $`a\lambda a`$. Using the bilinear identity (5.1) and $`aSq^2a=aSq^2\lambda `$, one finds that $`\mathrm{\Theta }_M(1)^{f(\lambda )}\mathrm{\Theta }_M`$. It follows that the total parity anomaly is $`(1)^{q(TX)+f(\lambda )}`$. On the other hand, the fermion measure of the heterotic string on $`X`$ transforms under $`(1)^F`$ by the same factor $`(1)^{q(TX)+f(\lambda )}`$. It is shown in that the heterotic string measure is well-defined, so we conclude that $`q(TX)+f(\lambda )=0\mathrm{mod}2`$, and hence that parity is a good symmetry of $`M`$-theory.
Our third application concerns the instability of some Type IIB branes wrapping homologically nontrivial cycles. The $`K`$-theory interpretation of $`D`$-branes means that $`D`$-branes cannot be wrapped on certain cycles; it also means that $`D`$-branes wrapped on certain cycles are unstable even though the cycles are nontrivial in homology. Let $`Q`$ be a cycle Poincarรฉ dual to an integral class $`aH^{\mathrm{even}}(X,๐)`$. If $`Sq^3(a)0`$ then, as we have mentioned, we cannot wrap a $`D`$-brane on $`Q`$. However, as stressed near (4.1), $`K(X)`$ is, essentially, the cohomology of $`Sq^3`$. Thus, if $`a`$ is โclosed,โ that is $`Sq^3(a)=0`$, then $`a`$ can be lifted to a $`K`$-theory class $`x`$, but if $`a`$ is โexact,โ that is $`a=Sq^3(a_0)`$ for some $`a_0`$, then one can take $`x=0`$. A $`D`$-brane whose lowest RR charge is given by such an $`a`$ can in fact be unstable, even though the class $`a`$ is nonzero in cohomology. Annihilation of such a $`D`$-brane occurs via nucleation and subsequent annihilation of $`D9\overline{D9}`$ pairs. This follows from the $`K`$-theoretic interpretation of the work of Sen on brane-antibrane annihilation.
7. Conclusions, Further Results, and Open Problems
The matching described above and in between the $`M`$-theory formalism based on $`E_8`$ and the Type IIA formalism based on $`K`$-theory gives considerable added confidence in both. In particular, we gain added confidence that not only $`D`$-brane charges, but also RR fluxes, should be classified by $`K`$-theory. This is an important conceptual change from the $`K`$-theoretic classification of $`D`$-brane charge; among other things, it suggests that RR fields, and not just $`D`$-branes, should somehow be associated with vector bundles.
We have focused here on the simplest case of the computation of in order to illustrate some of the central ideas. The simplifying topological assumptions we have made are relaxed in . Also, in we extend the computation sketched above to include $`G_20`$ in Type IIA theory; in $`M`$-theory, this corresponds letting $`Y`$ be a circle bundle over $`X`$ with Euler class $`G_2/2\pi `$. After a lengthier analog of the above computation with some additional ingredients added, the phases turn out to agree.
One interesting general lesson that emerges from is that when one takes torsion into account there is no direct relation between the flux $`G_4`$ in IIA theory and the four-form $`G`$ in $`M`$-theory. They have different underlying integral quantizations, and there is no 1-1 correspondence of the terms in the $`M`$-theory and the IIA theta functions. It is only after applying the โintegral equation of motion,โ $`Sq^3(a)=0`$ that one can compare results.
As for future directions, it should be very interesting to compare the absolute normalization of the $`M`$-theory and Type IIA partition functions; this depends on the one-loop determinants as well as some other overall normalization constants which arise when $`Sq^3`$ does not annihilate the torsion. One would like to extend the computation to include $`D`$-brane and $`M`$-brane instanton effects. Another, more difficult, open problem concerns the proper interpretation of nonzero values of $`G_0`$. While it is straightforward to include the effects of $`G_0`$ in the IIA partition function, comparing the results to $`M`$-theory presents an interesting and unsolved challenge.
Our computation confirms the utility of relating the $`C`$-field of $`M`$-theory to $`E_8`$ gauge theory as in . Other clues of a possible role of $`E_8`$ in the formulation of $`M`$-theory include the possibility of writing eleven-dimensional supergravity in terms of gauge fields of a noncompact form of $`E_8`$ , evidence for propagating $`E_8`$ gauge fields in $`M`$-theory on a manifold with boundary , and further issues considered in .
Finally, we mention that these considerations lead to an unresolved question in the case of Type IIB superstring theory. The problem is to reconcile the $`SL(2,๐)`$ symmetry of this theory with the $`K`$-theoretic interpretation of RR charges and fluxes. Although we have found some nontrivial partial results relevant to this problem, the main puzzle remains unsolved. Nevertheless, we hope that the clarification of the relation of $`M`$-theory and $`K`$-theory will play some role in the resolution.
Acknowledgments
We would like to thank M. F. Atiyah, D. Freed, J. Harvey, M. J. Hopkins, C. Hull, P. Landweber, J. Morgan, G. Segal, and I. M. Singer for discussions and explanations. The work of GM is supported by DOE grant DE-FG02-96ER40949. The work of EW has been supported in part by NSF Grant PHY-9513835 and the Caltech Discovery Fund. The work of DED has been supported by DOE grant DE-FG02-90ER40542. DED would also like to thank D. Christensen, K. Dasgupta, J. Gomis, C. Rezk and especially L. Nicolaescu and S. Stolz for useful discussions.
References
relax R. Minasian and G. Moore, โ$`K`$-theory and Ramond-Ramond charge,โ hep-th/9710230; JHEP 9711 (1997) 002. relax E. Witten, โ$`D`$-Branes And $`K`$-Theory,โ hep-th/9810188; JHEP 9812 (1998) 019. relax K. Olsen and R.J. Szabo, โConstructing $`D`$-Branes from $`K`$-Theory,โ hep-th/9907140. relax E. Witten, โDuality Relations among Topological Effects in String Theory,โ hep-th/9912086. relax G. Moore and E. Witten, โSelf-duality, RR fields and $`K`$-Theory,โ hep-th/9912279. relax D.S. Freed and M. J. Hopkins, โOn Ramond-Ramond Fields and $`K`$-Theory,โ hep-th/0002027. relax E. Diaconescu, G. Moore, and E. Witten, โ$`E_8`$ Gauge Theory, and a Derivation of $`K`$-Theory from $`M`$-Theory,โ hep-th/0005090 relax E. Witten, โFive-Brane Effective Action In $`M`$-Theoryโ, J. Geom. Phys. 22 (1997) 103, hep-th/9610234. relax M. F. Atiyah, I. M. Singer, โThe index of elliptic operators: Vโ Ann. Math. 93 (1971) 139. relax E. Witten, โOn Flux Quantization in $`M`$-Theory and the Effective Action,โ hep-th/9609122; Journal of Geometry and Physics, 22 (1997) 1. relax R. Stong, โCalculation of $`\mathrm{\Omega }_{11}^{spin}(K(๐,4))`$โ in Unified String Theories, 1985 Santa Barbara Proceedings, M. Green and D. Gross, eds. World Scientific 1986. relax E. Witten, โTopological Tools In Ten-Dimensional Physicsโ, in Unified String Theories, 1985 Santa Barbara Proceedings, M. Green and D. Gross, eds. World Scientific 1986. relax M. F. Atiyah, V. Patodi, and I. M. Singer, โSpectral asymmetry and Riemannian geometry.โ Math. Proc. Camb. Phil. Soc. 77 (1975) 43; Math. Proc. Camb. Phil. Soc. 78 (1975) 405; Math. Proc. Camb. Phil. Soc. 79(1976) 71. relax D.S. Freed and E. Witten, โAnomalies in String Theory with $`D`$-branes,โ hep-th/9907189. relax H. Nicolai, โ$`d=11`$ supergravity with local $`SO(16)`$ invariance,โ Phys. Lett. 187B(1987)316; S. Melosch and H. Nicolai, โNew Canonical Variables for $`d=11`$ supergravity,โ hep-th/9709227 relax P. Horava and E. Witten, โHeterotic and Type I String Dynamics from Eleven Dimensionsโ, Nucl.Phys. B460 (1996) 506, hep-th/9510209; โEleven-Dimensional Supergravity on a Manifold with Boundaryโ, Nucl.Phys. B475 (1996) 94, hep-th/9603142. relax P. Horava,โ $`M`$-Theory as a Holographic Field Theory,โ hep-th/9712130, Phys.Rev. D59 (1999) 046004; P. Horava and M. Fabinger, โCasimir Effect Between World-Branes in Heterotic $`M`$-Theory,โ hep-th/0002073
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# Knotted Spheres and Graphs in Balls
## 1 Knotted Spheres
This paper explores embedding punctured spheres in balls. I would like to thank Mike Freedman for suggesting the question and for insightful comments. I would also like to thank Cameron Gordon, John Luecke, Martin Scharlemann, and Ying-Qing Wu for helpful comments along the way.
We first must set up some definitions that greatly simplify the statements of the theorems. Throughout the paper let $`B`$ be the unit ball in $`^3`$. Let $`S`$ be the boundary of $`B`$.
###### Definition 1.1.
A complete graph on $`n`$ vertices whose vertices $`\{v_1,\mathrm{}v_n\}`$are disjoint disks on $`S`$, and whose edges $`\{e_{\{i,j\}},1i<jn\}`$are properly embedded geodesics in $`B`$ with the property that if $`st`$, $`v_sv_t=\mathrm{}`$, is a standard unlinked n-graph.
###### Definition 1.2.
A complete graph on $`n`$ vertices whose vertices $`\{v_1,\mathrm{}v_n\}`$are disjoint disks on $`S`$ and whose edges $`\{e_{\{i,j\}},1i<jn\}`$are properly-embedded, disjoint unknotted arcs in $`B`$ with the property that no two edges are linked in the ball, is called a pairwise unlinked n-graph.
###### Definition 1.3.
Given a graph $`\mathrm{\Gamma }`$ in $`B`$, and $`\mathrm{\Sigma }`$ a properly embedded n-holed sphere in $`B`$, such that $`\mathrm{\Sigma }=_{i=1}^nv_i`$ and $`\mathrm{\Sigma }e_{\{i,j\}}=\mathrm{}`$ for all $`i`$, $`j`$. $`\mathrm{\Sigma }`$ is called an enveloping n-holed sphere for $`\mathrm{\Gamma }`$. $`B\mathrm{\Sigma }`$ consists of two components. The part containing $`\{e_{\{i,j\}},1i<jn\}`$is called the inside of $`\mathrm{\Sigma }`$, the other component is called the outside.
###### Definition 1.4.
An enveloping n-holed sphere $`\mathrm{\Sigma }`$ is standard, if it is boundary parallel, meaning there is a product structure on the outside of $`\mathrm{\Sigma }`$ taking the interior of $`\mathrm{\Sigma }`$ to $`S_{i=1}^nv_i`$, but leaving $`\mathrm{\Sigma }`$ fixed.
###### Definition 1.5.
Let a core of an enveloping n-holed sphere $`\mathrm{\Sigma }`$ be a 1-complex, such that the boundary of a regular neighborhood of the complex co-bounds a product region with $`\mathrm{\Sigma }`$. Let the star core be the core with exactly $`n`$ edges and one vertex of valence of $`n`$. We shall call the vertices of the star core $`\{\nu _1,\mathrm{}\nu _n,\nu _{n+1}\}`$where $`\nu _iv_i`$ and $`\nu _{n+1}`$ is the vertex of valence $`n`$. We call the edges of the star core $`\{ฯต_1,\mathrm{}ฯต_n\}`$where $`ฯต_i`$ is the edge containing $`\nu _i`$.
## 2 Pairwise unlinked n-graphs
### 2.1 The Core Lemmas
In this section we introduce a couple of basic lemmas that will increase our insight to the theorem and will be useful in some of the cases.
###### Lemma 2.1.
Given $`ฯต_i`$ and $`ฯต_j,ij`$, two edges of the star core of an enveloping n-holed sphere for a pairwise unlinked n-graph, $`ฯต_iฯต_j`$ may never be a knotted arc in $`B`$.
###### Proof.
$`ฯต_iฯต_j`$ makes up the core of the cylinder that is left over if the n-holed sphere is compressed along disks parallel to all of the uninvolved vertices. If the core of the cylinder were knotted (see Figure 1), then by standard satellite knot theory so is any arc running through the cylinder, but the edge between these two vertices in our graph is unknotted and can be assumed to be inside the cylinder, so this cannot be the case. โ
###### Lemma 2.2.
Let $`ฯต_i`$, $`ฯต_j`$, and $`ฯต_k`$ be three edges of a star core of an enveloping n-holed sphere for a pairwise unlinked n-graph. Let $`ฯต_i^k`$ and $`ฯต_j^k`$ be edges obtained from $`ฯต_i`$ and $`ฯต_j`$ by contracting $`ฯต_k`$ and perturbing the arcs slightly so their end points are disjoint. Then $`ฯต_i^k`$ and $`ฯต_j^k`$ are not linked.
###### Proof.
If $`ฯต_i^k`$ and $`ฯต_j^k`$ were linked, then $`e_{\{i,k\}}`$ and $`e_{\{j,k\}}`$ (edges of the the original unlinked n-graph running from vertex $`i`$ to vertex $`k`$ and $`j`$ to $`k`$ respectively) would have to be linked.
\[C\] and \[GF\] are good places to look for an introduction to rational tangles. As explained in \[GF\] a rational tangle $`T`$ is assigned a rational number $`F(T)`$ corresponding to a simple continued fraction. Two tangles $`T_1`$, $`T_2`$ are isotopic if and only if $`F(T_1)=F(T_2)`$. Let $`N(T)`$ be the knot obtained by from $`T`$ by connecting the ends of the tangle in the manner called the numerator of $`T`$ and $`D(T)`$ the denominator (See Figure 2).
###### Lemma 2.3.
If $`T`$ is a rational tangle with $`D(T)`$ an unknot, then $`F(T)=p,p`$. Therefore $`T`$ may be undone by merely twisting two of the vertices around each other $`p`$ times leaving the other two vertices fixed.
###### Proof.
It is well known that if $`F(T)=p/q`$, then $`D(T)`$ produces a $`p/q`$ 2-bridge link, which is trivial if and only if q=1. Similarly, of course, $`N(T)`$ is an unknot if and only if $`p=1`$ so in that case $`F(T)=1/q,q`$ (the picture is just rotated ninety degrees). \[M\] section 9.3 and \[BZ\] sections 12.A and 12.B have expositions on the these facts. The classification was first done in \[S\].
### 2.2 The case, $`n=3`$.
###### Theorem 2.4.
Every enveloping 3-holed sphere for a pairwise unlinked 3-graph is standard.
###### Proof.
By Lemma 2.1 the star core is not knotted, so we may assume it consists of one straight edge $`ฯต_1ฯต_2`$ running from $`\nu _1`$ the north pole to $`\nu _2`$ at the south pole and another edge $`ฯต_3`$ meeting this edge at $`\nu _4`$ at the origin and winding around in some manner before ending up at a $`\nu _3`$ somewhere on the Southern hemi-sphere (as in Figure 3).
Act upon $`ฯต_3`$ by an ambient isotopy leaving $`ฯต_1ฯต_2`$ fixed until $`ฯต_3`$ lies entirely below the equatorial disk except for at $`\nu _4`$ (Figure 4). Now if we examine the two edges of the core in $`B^{}`$ the Southern hemi-ball (the ball we get by cutting the original ball in half along the equatorial disk) and pull $`ฯต_2`$ and $`ฯต_3`$ apart slightly so that they no longer intersect to get $`ฯต_2^{}`$ and $`ฯต_3^{}`$, the coreโs edges cannot be linked by Lemma 2.2 (this is the same as contracting $`ฯต_1`$ and looking at $`ฯต_2^1`$ and $`ฯต_3^1`$. This together with the fact that neither $`ฯต_2`$ nor $`ฯต_3`$ can be knotted in $`B^{}`$ by Lemma 2.1, means that $`ฯต_2^{}ฯต_3^{}`$ is just a rational tangle in $`B^{}`$.
Now by Lemma 2.1 we know that the rational tangle to which they correspond, must be an unknot if vertices on the equatorial disk are connected as are the vertices on the sphere. By Lemma 2.3 we may assume that the rational tangle is obtained merely by twisting the two vertices on the southern hemisphere around each other, and therefore that it may be untangled without affecting the rest of the core (Figure 5). Therefore the core is standard, as must be the enveloping n-holed sphere. โ
### 2.3 The Case $`n=4`$
The argument for the case $`n=4`$ is particularly interesting, because the methods for the previous case fall slightly shy of working, but the counterexample from the case $`n5`$ also just fails to disprove it. I would like to thank Ying-Qing Wu for ideas that were particularly helpful in this section.
###### Theorem 2.5.
Given $`n4`$ every enveloping n-holed sphere for a pairwise unlinked n-graph $`\mathrm{\Gamma }`$ is standard.
###### Proof.
Assume there is a counterexample for $`n=4`$ and examine what it must look like. $`B`$ minus the interior of the non-standard ball, bounded by the non-standard enveloping 4-holed sphere is homeomorphic to $`S^2\times I`$ minus an open neighborhood of four arcs $`\{ฯต_1,ฯต_2,ฯต_3,ฯต_4\}`$ which run from the inner sphere to the outer sphere. Since the enveloping 4-holed sphere is non-standard, the arcs cannot each simultaneously be isotoped in $`S^2\times I`$ to be $`pt\times I`$.
Since the original edges of $`\mathrm{\Gamma }`$ were unknotted and pairwise unlinked, they could be thought of as rational tangles. Examine the three pairs of edges of $`\mathrm{\Gamma }`$ corresponding to the three possible pairings of the vertices, $`(e_{\{1,2\}},e_{\{3,4\}})`$, $`(e_{\{1,3\}},e_{\{2,4\}})`$, $`(e_{\{1,4\}},e_{\{2,3\}})`$. Pick one pair, say $`(e_{\{1,2\}},e_{\{3,4\}})`$. Because the pair is a rational tangle, there is a disk $`D_1`$ in $`B`$ that separates the edges. Assume $`D_1`$ intersects $`\mathrm{\Sigma }`$ minimally. An innermost curve of intersection on $`D_1`$ yields a subdisk $`D_1^{}`$ which compresses $`\mathrm{\Sigma }`$ into two annuli $`A_{\{1,2\}}`$ and $`A_{\{3,4\}}`$ which are just the boundary of a regular neighborhood of $`e_{\{1,2\}}`$ and $`e_{\{3,4\}}`$ respectively. The same arguments can be made for the other two pairs. Thus, there are three disks that can be added to $`S^2\times 0`$ that separate the end points of the $`ฯต_i`$ into pairs, each yielding a rational tangle.
If we take the branched double cover of $`S^2\times I`$ over the four arcs $`ฯต_1,ฯต_2,ฯต_3,ฯต_4`$, we get a manifold $`M`$ with boundary two tori $`T_0`$, $`T_1`$ because the branched double cover of a sphere over four points is a torus. Since adding $`D_1^{}`$ to $`B`$ minus the โinsideโ of $`\mathrm{\Sigma }`$, thought of as $`S^2\times I`$ minus a neighborhood of the four branching arcs yields the exterior of a rational tangle, and the branched double cover of a ball over a rational tangle is, of course, just a solid torus. Thus, $`D_1^{}`$ lifts to a disk that gives a filling of $`T_0`$ that turns $`M`$ into a solid torus, so $`M`$ is just a solid torus minus a knot. Note that the three pairings of the vertices gives three different fillings of $`T_0`$ each of which yields a solid torus.
###### Theorem 2.6.
\[B Corollary 2.9\] If $`k`$ is a nontrivial knot in $`D^2\times S^1`$ such that $`k`$ is not parallel into $`D^2\times S^1`$ and there exists more than one nontrivial surgery on $`k`$ yielding $`D^2\times S^1`$, then $`k`$ is equivalent either to $`W_3^1W_7^3`$ or its mirror image $`W_3^1W_7^3`$.
$`W_3^1W_7^3`$ is the $`(2,3,7)`$ pretzel knot embedded as shown in Figure 8. We refer to the this as the Berge knot. Since the arcs were not standard, we must lift to the Berge knot or a knot $`k`$ parallel into $`T_1`$, the boundary of $`D^2\times S^1`$.
###### Proposition 2.7.
If $`k`$ is parallel into the boundary of $`D^2\times S^1`$ then $`k`$ fails to produce a counter example to Theorem 2.5.
The proof of the proposition requires two steps. First we prove that $`k`$ can be assumed to have the standard embedding in $`D^2\times S^1`$ by showing there is a unique strong inversion on $`k`$. Second we prove that the punctured sphere produced by quotienting out by the $`_2`$ symmetry is either boundary parallel or else violates Lemma 2.1 and therefore is not a counterexample.
###### Proof.
To prove that $`k`$ must be embedded in the standard manner we examine an annulus $`A`$ that runs from $`T_0`$, the torus boundary component corresponding to $`k`$ in the exterior of $`k`$, to $`T_1`$, the boundary of the solid torus. $`A`$ can be assumed to be embedded (see \[CF\]). Let $`F`$ be the involution of $`D^2\times S^1`$. Let $`F(A)=A^{}`$.
###### Lemma 2.8.
There is a unique strong inversion of $`k`$, a torus knot parallel to the boundary of a solid torus.
###### Proof.
Our first goal is to show that $`A`$ can be chosen such that $`A=A^{}`$. We choose $`A`$ to have a minimal number of intersections with $`A^{}`$. It is clear that $`A`$ can be assumed to be fixed by $`F`$, so perturbing $`A`$ slightly we can assume that $`AA^{}=\mathrm{}`$. Now $`AA^{}`$ must consist of simple closed curves. An innermost disk argument suffices to show that each of the simple closed curves is essential on $`A`$ and $`A^{}`$, so the intersection consists of disjoint simple closed curves $`c_1,c_2,\mathrm{}c_n`$ parallel to $`A`$ on $`A`$. Let $`c_1`$ be the curve of intersection closest to the boundary component of $`A`$ on $`T_1`$, $`c_2`$ be the second, and so on increasing the index as the curves move towards the boundary component on $`T_0`$. Likewise on $`A^{}`$ the intersection consists of parallel essential circles $`c_1^{},c_2^{},\mathrm{}c_n^{}`$ labeled in the same manner. Let $`A_i`$ be the sub-annulus of $`A`$ running from the boundary component on $`T_1`$ to $`c_i`$, and $`A_i^{}`$ be the sub-annulus of $`A^{}`$ running from the boundary component on $`T_1`$ to $`c_i^{}`$. Let $`c_j`$ be the curve of intersection on $`A`$ that corresponds to the intersection with $`A^{}`$ at $`c_1^{}`$. If we cut and paste $`A`$ replacing $`A_j`$ by a push off of $`A_1^{}`$ we reduce the number of intersections of $`A`$ and $`A^{}`$ by at least one. In order to preserve the property that $`F(A)=A^{}`$ we must also replace $`A_j^{}`$ by a push off of $`A_1`$. This, however, cannot increase the number of intersections, so we have a new annulus running from $`T_1`$ to $`T_0`$ that has fewer intersections with its image under $`F`$, contradicting minimality. Thus, we can assume that $`AA^{}=\mathrm{}`$.
This, however, implies that restricting to the solid torus between $`A`$ and $`A^{}`$, $`F`$ takes $`A\times I`$ to itself, exchanging $`A\times 0`$ with $`A\times 1`$. This in turn implies that there must be an annulus in $`A\times I`$ that is fixed by $`F`$.
Now that we know that $`A`$ is fixed, we use it to show that we have a standard inversion of $`k`$. Let $`D`$ be a meridional disk for $`T_1`$ that is fixed by $`F`$. Examine $`DA`$. In general the intersection pattern on $`D`$ will look something like Figure 6 with a collection of arcs running from $`T_0`$ (which punctures $`D`$ several times) to $`T_1`$ and a collection that run from one of the punctures from $`T_0`$ to another. We can remove the arcs running from $`T_0`$ to $`T_0`$ by picking an outermost arc on $`A`$ that runs from one component of $`A`$ to itself. This small disk gives an isotopy of $`A`$ together with $`k`$ that reduces the number of intersections of $`A`$ with $`D`$. because $`F(A)=A`$ we can simultaneously do a second isotopy of $`A`$ that also reduces the number of intersections of $`A`$ with $`D`$ and preserves the symmetry of $`k`$. Thus, we can assume that $`DA`$ consists solely of arcs running from $`k`$ to $`T_1`$. This, however, shows that we have the standard symmetry for $`k`$ because cutting the solid torus along $`D`$ turns it into a cylinder and $`A`$ becomes bands running from the top of the cylinder to the bottom in the unique way possible. โ
Now we need only argue that torus knots with standard embeddings fail to give a counterexample. Let $`k`$ be a torus knot embedded in $`D^2\times S^1`$, fixed by an involution $`F`$ of $`D^2\times S^1`$. Let $`A`$ be the fixed annulus above. Let $`M`$ be $`S^2\times I`$ with four branching arcs $`ฯต_1,\mathrm{}ฯต_4`$, the quotient of the exterior of $`k`$ in $`D^2\times S^1`$ by $`F`$. In the quotient, $`T_0`$ maps down to $`S^2\times 0`$ which we will designate $`S_0`$ and $`T_1`$ maps to $`S^2\times 1`$ designated $`S_1`$.
###### Lemma 2.9.
Either $`S_0ฯต_1ฯต_2ฯต_3ฯต_4`$ violates Lemma 2.1 or it is standard and therefore in either case is not a counterexample to Theorem 2.5.
###### Proof.
Because $`ฯต_1,\mathrm{}ฯต_4`$, each run from $`S_0`$ to $`S_1`$ the two arcs $`a_1`$, $`a_2`$ that are fixed in $`D^2\times S^1`$ by $`F`$ must each intersect $`k`$ in exactly one point. In the knot exterior $`a_1`$ and $`a_2`$ are then broken each into two arcs with one end point of each arc on $`T_1`$ and one end point on $`T_0`$. In $`M`$ the $`a_i`$ then become the four branching arcs $`ฯต_1\mathrm{}ฯต_4`$.
Under the quotient, the annulus $`A`$ becomes a rectangular disk $`D`$ which without loss of generality runs down $`ฯต_3`$ along $`S_0`$ up $`ฯต_4`$ and back along $`S_1`$ and is disjoint from $`ฯต_1`$ and $`ฯต_2`$. This is clear because $`A(a_1a_2)`$ consisted of the preimage of $`ฯต_3`$ and $`ฯต_4`$ but was disjoint from the preimage of $`ฯต_1`$ and $`ฯต_2`$ (recall that $`A`$, like the $`a_i`$ is fixed by $`F`$ and runs from $`T_0`$ to $`T_1`$). This means that $`M`$ looks exactly like Figure 7, where $`ฯต_1`$ and $`ฯต_2`$ form a rational tangle inside the ball designated $`T`$.
By Lemma 2.3 we see that $`ฯต_1ฯต_2`$ must be the tangle $`T`$ that results from two horizontal arcs whose eastern vertices are twisted $`n`$ times around each other or $`ฯต_1S_0ฯต_2`$ will be knotted violating Lemma 2.1. On the other hand, if $`ฯต_1ฯต_2`$ is the tangle $`T`$ above, then twisting the eastern portion of $`M`$, $`n`$ times shows $`M`$ is homeomorphic to the standard picture and therefore again fails to be a counterexample.
Note: One can in fact prove that (if $`k`$ is a torus knot other than the unknot) $`k`$ never produces a standard enveloping sphere, but instead always creates one that violates Lemma 2.1.
Lemma 2.9 completes the proof that $`k`$ cannot be parallel to $`T_1`$ and therefore must be a Berge knot if it is to produce a counterexample. โ
We may transform the traditional picture of the Berge Knot to a symmetric one as in Figure 8. Snappea \[W\] tells us that this knot (entered as a link) has exactly one $`_2`$ symmetry, which we can now see.
We quotient out by the symmetry to either get a counterexample, or proof that there is none. We get $`S^2\times I`$ minus four arcs $`ฯต_1,ฯต_2,ฯต_3,ฯต_4`$. We will show that three of the arcs violate Lemma 2.2 (See Figure 9).
Label the edge omitted from the picture $`ฯต_4`$. Label the horizontal edge that has one of its vertices on the eastern side of the sphere $`ฯต_3`$. Contraction of $`ฯต_3`$ takes edges $`ฯต_1`$ and $`ฯต_2`$ to $`ฯต_1^3`$ and $`ฯต_2^3`$ as pictured in Figure 10.
To see that $`ฯต_1^3`$ and $`ฯต_2^3`$ are in fact linked we call on the following well known facts about rational tangles. (See, for example, \[M\] Theorem 9.3.1.)
###### Theorem 2.10.
1. A 2-bridge knot (or link) is the denominator of some rational tangle
2. Conversely, the denominator of a rational tangle is a 2-bridge knot (or link).
###### Corollary 2.11.
If $`k`$ is the denominator of a rational tangle, then $`k`$ is prime.
The denominator of the tangle in Figure 10 is the connect sum of a trefoil and a figure eight knot and therefore it is not a rational tangle by Corollary 2.11. Since neither of the arcs is knotted, they must be linked violating Lemma 2.2. Since this was the only possible counterexample, Theorem 2.5 must be true.
### 2.4 The Case $`n5`$
###### Theorem 2.12.
There exist pairwise unlinked n-graphs with non-standard enveloping n-holed spheres for all $`n`$, $`n5`$.
###### Proof.
Figure 11 shows a $`\theta _n`$ curve in a ball. Such a graph is well known not to be standard, but every subgraph is standard. Let a $`\theta _n`$ curve be the star core of an enveloping n-holed sphere. Although it is not standard, it supports a pairwise unlinked n-graph $`\mathrm{\Gamma }`$. If we think of the enveloping sphere as bounding a central (round) ball with n tentacles running to the boundary, we can picture the pairwise unlinked graph as being a standard unlinked graph in the central ball which is extended by a product down each of the tentacles. Now since $`n5`$ any two edges of $`\mathrm{\Gamma }`$ miss at least one vertex, but there is an isotopy of any $`n1`$ arcs of a $`\theta _n`$ curve that makes those arcs appear standard. Likewise since the edges of $`\mathrm{\Gamma }`$ completely miss one of the tentacles, they may be pictured as being embedded in a ball bounded by a standard enveloping $`(n1)`$-holed sphere. The arcs remain standard within the central ball and are extended by a product down the tentacles throughout the entire process, so clearly the edges are not linked pairwise. (See Figure 11).
## 3 Standard unlinked n-graphs
###### Theorem 3.1.
Given a standard unlinked n-graph $`\mathrm{\Gamma }`$, every enveloping n-punctured sphere $`\mathrm{\Sigma }`$ is standard.
Note: Since the edges of our graph are geodesics in this case, Morse theory assures us that we can assume that there is an embedded disk $`D`$ in $`B`$ whose interior is disjoint from all of the edges of $`\mathrm{\Gamma }`$ and whose boundary consists of two arcs $`\alpha `$ and $`\beta `$, where $`\alpha `$ is one of the edges of the graph, $`\beta `$ is strictly contained in $`B`$, $`\alpha =\beta `$ and the interior of $`\beta `$ is disjoint from all of the edges of $`\mathrm{\Gamma }`$.
Letโs examine how $`D`$ meets $`\mathrm{\Sigma }`$.
###### Claim 3.2.
We may assume that $`D\mathrm{\Sigma }`$ contains no simple closed curves.
###### Proof.
Assume $`D`$ is chosen with a minimal number of intersections with $`\mathrm{\Sigma }`$. Examine an innermost curve $`\delta `$ on $`D`$. If $`\delta `$ is not essential on $`\mathrm{\Sigma }`$, then there is an obvious isotopy through which this intersection could have been eliminated, so we may assume it is essential on $`\mathrm{\Sigma }`$.
Therefore the innermost loop gives us a compressing disk for $`\mathrm{\Sigma }`$. Homology is enough to assure us that the disk must be on the inside of $`\mathrm{\Sigma }`$ (the component of $`B\mathrm{\Sigma }`$ containing $`\{e_{\{i,j\}},1i<jn\}`$). Since $`\delta `$ is assumed to be essential in $`\mathrm{\Sigma }`$, it must separate the vertices of $`\mathrm{\Sigma }`$ into two non-empty sets. With no loss of generality, let $`v_1`$ be in one set and $`v_2`$ be in the other. Since $`\delta `$ separates the vertices (as in Figure 12), the disk it bounds does too, and $`e_{12}`$ must intersect it. This, however, is a contradiction since the interior of $`D`$ is disjoint from the edges of $`\mathrm{\Gamma }`$. โ
Now we examine an outermost arc $`\gamma `$ on $`D`$. If $`\gamma `$ runs from one vertex of $`\mathrm{\Sigma }`$ to a different one, then it is obvious that the corresponding subdisk of $`D`$ is on the outside of $`\mathrm{\Sigma }`$. This gives us a compression disk that allows us to complete the proof by induction. (See Figure 13).
We may therefore assume that $`\gamma `$ connects a vertex, say $`v_n`$ to itself. Because $`D\mathrm{\Sigma }`$ is assumed to be minimal $`\gamma `$ must be essential on $`\mathrm{\Sigma }`$. If the disk $`\gamma `$ cuts off on $`D`$ is on the inside of $`\mathrm{\Sigma }`$, then once again it must separate $`\{v_1,\mathrm{}v_{n1}\}`$ into two sets and the argument proceeds as in the simple closed curve case. (See Figure 12).
Our final case therefore is that although $`\gamma `$ connects $`v_n`$ to itself, the disk is on the outside of $`\mathrm{\Sigma }`$. Compressing along this disk splits the n-punctured sphere $`\mathrm{\Sigma }`$ into two pieces $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_2`$. $`\mathrm{\Sigma }_1`$ is an r-punctured sphere and $`\mathrm{\Sigma }_2`$ is a $`n+1r`$-punctured sphere, where $`2rn1`$. By induction we may therefore assume that $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_2`$ are standard. $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_2`$ and their product structures are in different โhalvesโ of $`B`$. They are separated by the annulus formed by the boundary compressing $`v_n`$. The inverse of the boundary compression is a tunnel connecting the two boundary components of the annulus. Up to isotopy there is a unique arc running across the annulus, so there is a unique tunnel we can add to attain $`\mathrm{\Sigma }`$ from $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_2`$. Since the product structures on $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_2`$ are to the outside of the annulus, and the tunnel can be added extremely close to the boundary, it is clear that the product structure can be extended across the tunnel to give a product structure on $`\mathrm{\Sigma }`$, proving that it is standard.
Note: we needed the full strength of the complete graph here. Figure 14 shows a counter-example for $`n=3`$, if one edge is missing from $`\mathrm{\Gamma }`$. This counter-example may be generalized for any $`n`$ to give a counterexample for the complete graph on $`n`$ vertices, minus one edge.
## 4 The Infinite Case
###### Definition 4.1.
$`f:F\mathrm{\Sigma }^3`$ is a proper embedding of $`F`$, a sphere with a Cantor Set worth of punctures, in $`^3`$ if the pre-image of every compact set in $`^3`$ is a compact set in $`F`$ ($`H^3`$ could replace $`^3`$ throughout this section).
###### Definition 4.2.
Let $`l_\alpha `$ be a collection of geodesics in $`^3`$. Let $`l_\alpha `$ be contained in one connected component of $`^3\mathrm{\Sigma }`$. Again let the inside of $`\mathrm{\Sigma }`$ be the component of $`^3\mathrm{\Sigma }`$ containing $`\mathrm{\Gamma }`$ and the outside be the remaining component.
###### Definition 4.3.
$`\mathrm{\Sigma }`$ is said to be standard if there exists a product structure on the outside of $`\mathrm{\Sigma }`$ such that it is a product of the punctured sphere and a half open interval.
We may now ask, how many lines contained inside of $`\mathrm{\Sigma }`$ it takes to ensure that the embedding of $`\mathrm{\Sigma }`$ is standard.
###### Theorem 4.4.
Given $`\mathrm{\Sigma }`$, a proper embedding of $`F`$, a sphere with a Cantor Set worth of punctures in $`^3`$ and a set of geodesics which are dense in the Cantor Set, we may conclude that the punctured sphere is standard.
In this context saying that the geodesics are dense in the Cantor Set means that given any two punctures $`p_1`$ and $`p_2`$ of $`F`$ and any two neighborhoods of those punctures $`\mu _1`$ and $`\mu _2`$ on $`F`$, there exists a geodesic that runs from the image of some puncture in $`\mu _1`$ to the image of some puncture in $`\mu _2`$.
Note that though there are an uncountable number of points in the cantor set, we are only requiring a countable set of geodesics. Even if we wanted to satisfy the property of being dense for every point on the sphere, we would still only need a countable set of geodesics, since geodesics connecting all the points on the sphere with rational coordinates would suffice and the set of possible pairings of a countable set is itself a countable set.
Perhaps the easiest way to picture the scenario is to imagine the universal cover of a genus-two handlebody in hyperbolic three-space. The punctured sphere would be the boundary of this cover and the lines would be geodesics connecting points at infinity.
It is worthy of note that in the finite case we needed the complete graph, so we needed all possible geodesics, but in the infinite case with an uncountable number of punctures, we only need a countable number of edges.
We now begin the proof of the theorem.
###### Proof.
Choose a point on $`\mathrm{\Sigma }`$ to be the origin of $`^3`$.
Let $`S_n`$ be the sphere of radius $`n`$ centered at the origin in $`^3`$, and let $`B_n`$ be the ball that it bounds. We may alter $`\mathrm{\Sigma }`$ slightly if necessary so that we may assume that it intersects each $`S_i`$ transversally.
Fix $`i`$ and examine $`B_i`$. The pieces of $`\mathrm{\Sigma }`$ in $`B_i`$ may be broken into two sets. The first set consists of the connected piece containing the origin $`\mathrm{\Sigma }_{i1}`$, and the second set, all the other pieces, $`\{\mathrm{\Sigma }_{i2},\mathrm{}\mathrm{\Sigma }_{in}\}`$. We shall isotope $`\mathrm{\Sigma }`$ until $`B_i`$ contains one piece of the first type and none of the second on the โoutsideโ of $`\mathrm{\Sigma }`$ in $`B_i`$. The latter does not prevent us from claiming that the former is boundary parallel in $`B_i\mathrm{\Sigma }`$, so we do not worry about them. See Figure 15
We examine the pre-images of the $`\mathrm{\Sigma }_{ij}`$ in $`F`$. We now make an argument to show that we may assume that none of them have a boundary curve which is trivial in the fundamental group of $`F`$.
If there were a trivial boundary curve, we could choose an innermost one (on $`F`$). This would bound an embedded disk $`D`$ on $`\mathrm{\Sigma }`$ that meets $`S_i`$ in a simple closed curve. The boundary curve splits $`S_i`$ into two disks, each of which bounds a ball with $`D`$. If $`D`$ is on the interior of $`B_i`$, then we choose the ball that does not contain the surface $`\mathrm{\Sigma }_{i1}`$ ($`\mathrm{\Sigma }_{i1}`$, is connected and disjoint from $`D`$, so it can only be in one of the two balls). If $`D`$ is on the exterior of $`B_i`$ then we choose the smaller of the two balls (it is contained in the other ball). Either way we push $`D`$ across the ball and through $`S_i`$, eliminating at least its intersection with $`S_i`$ and possibly more extraneous intersections that were contained in the ball through which our isotopy was done.
Note that $`\mathrm{\Sigma }_{i1}`$ is unchanged away from its boundary and its boundary can only be changed by capping off trivial components. We continue this process until there are no more trivial components in the entire collection of $`\mathrm{\Sigma }_{ij}`$.
At the risk of sloppy notation we shall continue throughout to call the new surfaces $`\mathrm{\Sigma }_{ij}`$ carefully noting at each step that we still have done only a finite number of isotopies to a finite number of pieces.
We now notice that $`\mathrm{\Sigma }_{i1}`$ fits all of the criterion of the standard finite case. Since the geodesics are dense within the Cantor Set, at any finite stage there will be a complete graph in $`B_i`$ on the vertices that are given by the intersection of $`\mathrm{\Sigma }_{i1}`$ and $`S_i`$. Thus, by the previously proven finite case, $`\mathrm{\Sigma }_{i1}`$ is standard in $`B_i`$. We can use its inherited product structure to isotope $`\mathrm{\Sigma }`$ to make sure that it does not intersect $`B_i`$ on the outside of $`\mathrm{\Sigma }_{i1}`$.
Now we repeat the process for some $`k>i`$. We might worry that this process results in our pushing some piece of $`\mathrm{\Sigma }`$ an infinite number of times, but this is not the case, as every point in $`\mathrm{\Sigma }`$ is in some $`\mathrm{\Sigma }_{k1}`$ for large enough $`k`$ and our isotopies never affect points of $`\mathrm{\Sigma }_{k1}`$ that are not near the boundary of $`B_k`$. Thus, for each point there is some $`k`$ such that the point is left alone for good after $`k`$ steps. Since each step involved only a finite number of isotopies, each point is moved only a finite number of times.
The only thing left for us to check is that the product structure for $`\mathrm{\Sigma }_{k1}`$ can be chosen to correspond exactly with the product structure we already chose for $`\mathrm{\Sigma }_{i1}`$. $`B_i`$ may be left fixed as we do our operations for $`\mathrm{\Sigma }_{k1}`$, so naturally $`\mathrm{\Sigma }_{i1}`$ remains fixed, too.
Since $`\mathrm{\Sigma }_{i1}`$ is boundary parallel in $`B_i`$, we may substitute part of $`S_i`$ for it in $`\mathrm{\Sigma }_{k1}`$ and the resulting surface still contains a complete graph on one side and must be boundary parallel. If we concatenate its product lines in $`B_k`$ with the product lines of $`\mathrm{\Sigma }_{i1}`$ in $`B_i`$ we see a product structure that suits our desires as in Figure 16.
## 5 References
\[B\] J. Berge, The knots in $`D^2\times S^1`$ which have nontrivial Dehn surgeries that yield $`D^2\times S^1`$. Topology Appl. 38 (1991), no. 1, 1โ19.
\[BZ\] G. Burde and H. Zieschang, Knots. Walter de Gruyter, Inc. New York, NY, 1985.
\[C\] J. H. Conway, Enumeration of knots and links. Computational Problems in Abstract Algebra, Pergamon Press, 1970, 329โ359.
\[CF\] J. Cannon and C. Feustel Essential embeddings of annuli and mobius bands in 3-manifolds. Transactions of the AMS 215, (1976), 219โ239.
\[GF\] J. Goldman and L. Kauffman, Rational Tangles. Advances in Applied Mathematics 18 (1997), 300โ332.
\[HR\] J. Hempel and L. Roeling, Free factors of handlebody groups. preprint.
\[M\] K. Murasugi, Knot Theory and Its Applications, Birkhauser, Boston, 1996.
\[ST\] H. Schubert, Knoten mit zwei Brucken. Math. Zeit. 1956, no. 65, 133โ170
\[ST\] M. Scharlemann and A. Thompson, Detecting unknotted graphs in $`3`$-space. J. Differential Geom. 1991, no. 2, 539โ560
\[W\] Jeff Weeks, SnapPea: a computer program for creating and studying hyperbolic 3-manifolds, available by anonymous ftp from geom.umn.edu.
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# Using differential equations to compute two-loop box integralsPresented at Loops and Legs in Quantum Field Theory, April 2000, Bastei, Germany
## 1 Introduction
Perturbative corrections to many inclusive quantities have been computed to the two- and three-loop level in past years. From the technical point of view, these inclusive calculations correspond to the computation of multi-loop two-point functions, for which many elaborate calculational tools have been developed. In contrast, corrections to exclusive observables, such as jet production rates, could up to now only be computed at the one-loop level. These calculations demand the computation of multi-leg amplitudes to the required number of loops, which beyond the one-loop level turn out to be a calculational challenge obstructing further progress. Despite considerable progress made in recent times, many of the two-loop integrals relevant for the calculation of jet observables beyond next-to-leading order are still unknown. One particular class of yet unknown integrals appearing in the two-loop corrections to three jet production in electron-positron collisions, to two-plus-one jet production in electron-proton collisions and to vector boson plus jet production in proton-proton collisions are two-loop four-point functions with massless internal propagators and one external leg off-shell.
We elaborate on in this talk several techniques to compute multi-leg amplitudes beyond one loop. We demonstrate how integration-by-parts identities (already known to be a very valuable tool in inclusive calculations) and identities following from Lorentz-invariance (which are non-trivial only for integrals depending on at least two independent external momenta) can be used to reduce the large number of different integrals appearing in an actual calculation to a small number of master integrals. This reduction can be carried out mechanically (by means of a small chain of computer programs), without explicit reference to the actual structure of the integrals under consideration and can also be used for the reduction of tensor integrals beyond one loop.
The master integrals themselves, however, can not be computed from these identities. We derive differential equations in the external momenta for them. Solving these differential equations, it is possible to compute the master integrals without explicitly carrying out any loop integration, so that this technique appears to be a valuable alternative to conventional approaches for the computation of multi-loop integrals.
We demonstrate the application of these tools on several examples.
## 2 Reduction to master integrals
Any scalar massless two-loop integral can be brought into the form
$`I(p_1,\mathrm{},p_n)=`$ (1)
$`{\displaystyle \frac{\mathrm{d}^dk}{(2\pi )^d}\frac{\mathrm{d}^dl}{(2\pi )^d}\frac{1}{D_1^{m_1}\mathrm{}D_t^{m_t}}S_1^{n_1}\mathrm{}S_q^{n_q}},`$
where the $`D_i`$ are massless scalar propagators, depending on $`k`$, $`l`$ and the external momenta $`p_1,\mathrm{},p_n`$ while $`S_i`$ are scalar products of a loop momentum with an external momentum or of the two loop momenta. The topology (interconnection of propagators and external momenta) of the integral is uniquely determined by specifying the set $`(D_1,\mathrm{},D_t)`$ of $`t`$ different propagators in the graph. The integral itself is then specified by the powers $`m_i`$ of all propagators and by the selection $`(S_1,\mathrm{},S_q)`$ of scalar products and their powers $`(n_1,\mathrm{},n_q)`$ (all the $`m_i`$ are positive integers greater or equal to 1, while the $`n_i`$ are greater or equal to 0). Integrals of the same topology with the same dimension $`r=_im_i`$ of the denominator and same total number $`s=_in_i`$ of scalar products are denoted as a class of integrals $`I_{t,r,s}`$. The integration measure and scalar products appearing the above expression are in Minkowskian space, with the usual causal prescription for all propagators. The loop integrations are carried out for arbitrary space-time dimension $`d`$, which acts as a regulator for divergencies appearing due to the ultraviolet or infrared behaviour of the integrand (dimensional regularisation, ).
The number $`N(I_{t,r,s})`$ of the integrals grows quickly as $`r,s`$ increase, but the integrals are related among themselves by various identities. One class of identities follows from the fact that the integral over the total derivative with respect to any loop momentum vanishes in dimensional regularisation
$$\frac{\mathrm{d}^dk}{(2\pi )^d}\frac{}{k^\mu }J(k,\mathrm{})=0,$$
(2)
where $`J`$ is any combination of propagators, scalar products and loop momentum vectors. $`J`$ can be a vector or tensor of any rank. The resulting identities are called integration-by-parts (IBP) identities.
In addition to the IBP identities, one can also exploit the fact that all integrals under consideration are Lorentz scalars (or, perhaps more precisely, โ$`d`$-rotationalโ scalars) , which are invariant under a Lorentz (or $`d`$-rotational) transformation of the external momenta . These Lorentz invariance (LI) identities are obtained from:
$`(p_1^\nu {\displaystyle \frac{}{p_{1\mu }}}p_1^\mu {\displaystyle \frac{}{p_{1\nu }}}+\mathrm{}`$
$`\mathrm{}+p_n^\nu {\displaystyle \frac{}{p_{n\mu }}}p_n^\mu {\displaystyle \frac{}{p_{n\nu }}})I(p_1,\mathrm{},p_n)=0.`$
In the case of two-loop four-point functions, one has a total of 13 equations (10 IBP + 3 LI) for each integrand corresponding to an integral of class $`I_{t,r,s}`$, relating integrals of the same topology with up to $`s+1`$ scalar products and $`r+1`$ denominators, plus integrals of simpler topologies (i.e. with a smaller number of different denominators). The 13 identities obtained starting from an integral $`I_{t,r,s}`$ do contain integrals of the following types:
* $`I_{t,r,s}`$: the integral itself.
* $`I_{t1,r,s}`$: simpler topology.
* $`I_{t,r+1,s},I_{t,r+1,s+1}`$ : same topology, more complicated than $`I_{t,r,s}`$.
* $`I_{t,r1,s},I_{t,r1,s1}`$: same topology, simpler than $`I_{t,r,s}`$.
Quite in general, single identities of the above kind can be used to obtain the reduction of $`I_{t,r+1,s+1}`$ or $`I_{t,r+1,s}`$ integrals in terms of $`I_{t,r,s}`$ and simpler integrals - rather than to get information on the $`I_{t,r,s}`$ themselves.
If one considers the set of all the identities obtained starting from the integrand of all the $`N(I_{t,r,s})`$ integrals of class $`I_{t,r,s}`$, one obtains $`(N_{\mathrm{IBP}}+N_{\mathrm{LI}})N(I_{t,r,s})`$ identities which contain $`N(I_{t,r+1,s+1})+N(I_{t,r+1,s})`$ integrals of more complicated structure. It was first noticed by S. Laporta that with increasing $`r`$ and $`s`$ the number of identities grows faster than the number of new unknown integrals. As a consequence, if for a given $`t`$-topology one considers the set of all the possible equations obtained by considering all the integrands up to certain values $`r^{},s^{}`$ of $`r,s`$, for large enough $`r^{},s^{}`$ the resulting system of equations, apparently overconstrained, can be used for expressing the more complicated integrals, with greater values of $`r,s`$ in terms of simpler ones, with smaller values of $`r,s`$. An automatic procedure to preform this reduction by means of computer algebra is discussed in more detail in .
For any given four-point two-loop topology, this procedure can result either in a reduction towards a small number (typically one or two) of integrals of the topology under consideration and integrals of simpler topology (less different denominators), or even in a complete reduction of all integrals of the topology under consideration towards integrals with simpler topology. Left-over integrals of the topology under consideration are called irreducible master integrals or just master integrals.
## 3 Differential equations for master integrals
The IBP and LI identities allow to express integrals of the form (1) as a linear combination of a few master integrals, i.e. integrals which are not further reducible, but have to be computed by some different method.
At present, the complete set of master integrals for massless on-shell two-loop box integrals is known analytically up to finite terms in $`ฯต=(4d)/2`$. For massless two-loop box integrals with one off-shell leg, several topologies are yet to be computed analytically. A purely numerical approach for computing these integrals order by order in $`ฯต`$ has recently been proposed by Binoth and Heinrich in .
A method for the analytic computation of master integrals avoiding the explicit integration over the loop momenta is to derive differential equations in internal propagator masses or in external momenta for the master integral, and to solve these with appropriate boundary conditions. This method has first been suggested by Kotikov to relate loop integrals with internal masses to massless loop integrals. It has been elaborated in detail and generalised to differential equations in external momenta in ; first applications were presented in . In the case of four-point functions with one external off-shell leg and no internal masses, one has three independent invariants, resulting in three differential equations.
The derivatives in the invariants $`s_{ij}=(p_i+p_j)^2`$ can be expressed by derivatives in the external momenta:
$`s_{12}{\displaystyle \frac{}{s_{12}}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(+p_1^\mu {\displaystyle \frac{}{p_1^\mu }}+p_2^\mu {\displaystyle \frac{}{p_2^\mu }}p_3^\mu {\displaystyle \frac{}{p_3^\mu }}\right),`$
$`s_{13}{\displaystyle \frac{}{s_{13}}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(+p_1^\mu {\displaystyle \frac{}{p_1^\mu }}p_2^\mu {\displaystyle \frac{}{p_2^\mu }}+p_3^\mu {\displaystyle \frac{}{p_3^\mu }}\right),`$
$`s_{23}{\displaystyle \frac{}{s_{23}}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(p_1^\mu {\displaystyle \frac{}{p_1^\mu }}+p_2^\mu {\displaystyle \frac{}{p_2^\mu }}+p_3^\mu {\displaystyle \frac{}{p_3^\mu }}\right).`$
It is evident that acting with the right hand sides of (LABEL:eq:derivatives) on a master integral $`I_{t,t,0}`$ will, after interchange of derivative and integration, yield a combination of integrals of the same type as appearing in the IBP and LI identities for $`I_{t,t,0}`$, including integrals of type $`I_{t,t+1,1}`$ and $`I_{t,t+1,0}`$. Consequently, the scalar derivatives (on left hand side of (LABEL:eq:derivatives)) of $`I_{t,t,0}`$ can be expressed by a linear combination of integrals up to $`I_{t,t+1,1}`$ and $`I_{t,t+1,0}`$. These can all be reduced (for topologies containing only one master integral) to $`I_{t,t,0}`$ and to integrals of simpler topology by applying the IBP and LI identities. As a result, we obtain for the master integral $`I_{t,t,0}`$ an inhomogeneous linear first order differential equation in each invariant. For topologies with more than one master integral, one finds a coupled system of first order differential equations. The inhomogeneous term in these differential equations contains only topologies simpler than $`I_{t,t,0}`$, which are considered to be known if working in a bottom-up approach.
The master integral $`I_{t,t,0}`$ is obtained by matching the general solution of its differential equation to an appropriate boundary condition. Quite in general, finding a boundary condition is a simpler problem than evaluating the whole integral, since it depends on a smaller number of kinematical variables. In some cases, the boundary condition can even be determined from the differential equation itself.
In writing down the differential equations for all master integrals appearing in the reduction of two-loop four-point functions with up to one off-shell leg, one finds that all boundary conditions can be related to the following four one-scale two-loop integrals:
$`=`$ $`A_3\left(s\right)^{d3}`$
$`=`$ $`A_{2,LO}^2\left(s\right)^{d4}`$
$`=`$ $`A_4\left(s\right)^{d4}`$
$`=`$ $`A_6\left(s\right)^{d6}`$
These integrals fulfill homogeneous differential equations in $`s`$; their normalisation can therefore not be calculated in the differential equation method. They can however be evaluated explicitly using Feynman parameters .
## 4 Applications
Using the differential equation method, we computed in all master integrals with up to $`t=5`$ different denominators that can appear in the reduction of two-loop four-point functions with one off-shell leg.
Differential equations were also used in the computation of two-loop on-shell master integrals with planar and crossed topologies. These topologies contain two master integrals each. The differential equations were applied only to obtain a relation among both master integrals, while one of the master integrals was calculated by Smirnov for the planar case and Tausk for the non-planar case using a different method.
At this workshop, Nigel Glover reported first progress towards the computation of massless on-shell $`22`$ scattering amplitudes at two loops . In the reduction of planar amplitudes, it turned out that it is not sufficient to know the planar double box integrals up to their finite terms in $`ฯต`$. In the tensor reduction procedure, the following combination of master integrals arises:
$`{\displaystyle \frac{1}{ฯต}}\left[\text{}t\text{}\right]`$ (6)
$``$ $`{\displaystyle \frac{1}{ฯต}}\left[K_1(x,ฯต)+K_2(x,ฯต)\right],`$
where $`x=t/s`$. The functions $`K_1(x,ฯต)`$ and $`K_2(x,ฯต)`$ were computed only up to finite terms in $`ฯต`$ by Smirnov and Veretin in . The structure of the up to now unknown $`๐ช(ฯต)`$-term was subject of debate at this workshop. We did therefore decide after the workshop to attempt its computation using the differential equation method.
Starting from the differential equations for the massless double box integrals with one off-shell leg, which were obtained using the algorithms described in Section 3, we can derive the differential equations for the above on-shell integrals by constraining $`s_{12}+s_{13}+s_{23}=0`$ (see also ). Using $`s`$ and $`x`$ as independent variables, we obtain two homogeneous differential equations in $`s`$ (corresponding to a rescaling relation, ) and two coupled inhomogeneous equations in $`x`$, which can be employed to compute the two master integrals. Using the fact that the master integrals and their derivatives are regular in $`t`$ at $`t=s`$, the boundary condition at $`x=1`$ is inferred from the differential equations. This condition is however not yet sufficient to determine the boundary conditions for both master integrals, since it only fixes one combination of them. It does however turn out that it is sufficient to match the non-logarithmic terms in $`K_1(x,ฯต)`$, obtained by Smirnov in , up to finite order in $`ฯต`$ to find the boundary condition for $`K_1(x,ฯต)+K_2(x,ฯต)`$ up to $`๐ช(ฯต)`$. The differential equations are then solved order by order in $`ฯต`$ by expressing the unknown as sum of harmonic polylogarithms . Requiring the differential equation to be fulfilled and the boundary conditions to be matched, the coefficients of the harmonic polylogarithms in the ansatz can all be determined.
Using the same normalisation factor as , the coefficient of the $`๐ช(ฯต)`$-term of the above combination of double box integrals reads:
$`\left[K_1(x,ฯต)+K_2(x,ฯต)\right]|_{๐ช(ฯต)}`$ (7)
$`=`$ $`{\displaystyle \frac{4}{3}}\mathrm{ln}^4x+{\displaystyle \frac{4}{3}}\mathrm{ln}^3x4(4+3\pi ^2)\mathrm{ln}^2x`$
$`\left(46+21\pi ^2{\displaystyle \frac{16}{3}}\zeta (3)\right)\mathrm{ln}x`$
$`508\pi ^2{\displaystyle \frac{139}{60}}\pi ^4+{\displaystyle \frac{698}{3}}\zeta (3)`$
$`32S_{2,2}(x)+32\mathrm{ln}xS_{1,2}(x)`$
$`128\mathrm{Li}_4(x)+32(3\mathrm{ln}x\mathrm{ln}(1+x)`$
$`2)\mathrm{Li}_3(x)+32(\mathrm{ln}^2x+2\mathrm{ln}x`$
$`+\mathrm{ln}x\mathrm{ln}(1+x){\displaystyle \frac{5}{6}}\pi ^2)\mathrm{Li}_2(x)`$
$`+8\left(\mathrm{ln}^2x+\pi ^2\right)\mathrm{ln}^2(1+x)`$
$`+{\displaystyle \frac{16}{3}}(\mathrm{ln}^3x+6\mathrm{ln}^2x+6\pi ^22\pi ^2\mathrm{ln}x`$
$`+6\zeta (3))\mathrm{ln}(1+x)`$
$`+x[+{\displaystyle \frac{16}{3}}\mathrm{ln}^4x{\displaystyle \frac{52}{3}}\mathrm{ln}^3x+{\displaystyle \frac{46}{3}}\pi ^2\mathrm{ln}^2x`$
$`+\left(2{\displaystyle \frac{169}{3}}\pi ^2{\displaystyle \frac{496}{3}}\zeta (3)\right)\mathrm{ln}x`$
$`84{\displaystyle \frac{46}{3}}\pi ^2+{\displaystyle \frac{823}{360}}\pi ^4+536\zeta (3)`$
$`56S_{2,2}(x)+56\mathrm{ln}xS_{1,2}(x)`$
$`+40\mathrm{Li}_4(x)+(16\mathrm{ln}x56\mathrm{ln}(1+x)`$
$`104)\mathrm{Li}_3(x)(36\mathrm{ln}^2x+{\displaystyle \frac{8}{3}}\pi ^2`$
$`56\mathrm{ln}x\mathrm{ln}(1+x)104\mathrm{ln}x)\mathrm{Li}_2(x)`$
$`+14(\mathrm{ln}^2x+\pi ^2)\mathrm{ln}^2(1+x)`$
$`+{\displaystyle \frac{4}{3}}(16\mathrm{ln}^3x+39\mathrm{ln}^2x23\pi ^2\mathrm{ln}x`$
$`+39\pi ^2+42\zeta (3))\mathrm{ln}(1+x)].`$
## 5 Outlook
We have demonstrated how techniques developed for multi-loop calculation of two-point functions can be extended towards integrals with a larger number of external legs. In particular, we have shown that the use of differential equations in external invariants can be used as a powerful method to compute master integrals without carrying out explicit loop integrations. As a first example of the application of these tools in practice, we computed some up to now unknown two-loop four-point functions, relevant for jet calculus beyond the next-to-leading order. The most important potential application of these tools is the yet outstanding derivation of two-loop virtual corrections to exclusive quantities, such as jet observables.
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# 1 Introduction
## 1 Introduction
The Poisson cohomology of a Poisson manifold gives several informations on the geometry of the manifold. It was first introduced by Lichnerowicz in \[L\]. Unfortunately, the computation of these cohomology spaces is quite complicated and few explicit results have been found.
In the symplectic case, Poisson cohomology is naturally isomorphic to de Rham cohomology. The case of regular Poisson manifolds is discussed, for instance, in \[V\] and \[X\]. One can find some results on the Poisson cohomology of Poisson-Lie groups in \[GW\]. Some explicit computations are also done, for instance, in \[Co1\], \[Co2\] or \[G\].
In \[Cr\], Crainic links Poisson cohomology with the Morita equivalence. Finally, one can find some discussion on Poisson cohomology and Poisson homology in \[ELW\], \[B\] or \[FT\].
In the two-dimensional situation, some special cases on $`^2`$ have been studied. In \[V\], Vaisman began to compute the cohomology of $`(x^2+y^2)\frac{}{x}\frac{}{y}`$. His idea was to consider the homomorphism $`\iota ^{}:H^{}(^2)H^{}(^2\{(0,0)\})`$ induced by the inclusion $`\iota :^2\{(0,0)\}^2`$. A few years later, Nakanishi used this idea and computed the Poisson cohomology of quadratic Poisson structures on $`^2`$ (see \[N\]).
In the present paper, our approach is more direct and uses some tools arising from the theory of singularities. More precisely, we first study (in section 3) the normal forms of the โmost interestingโ germs at (0,0) of Poisson structures vanishing at (0,0), and we rediscover the list given by Arnold in \[A\]. These normal forms are of type
$$\mathrm{\Pi }=f(1+h)\frac{}{x}\frac{}{y},$$
where $`f`$ and $`h`$ are quasihomogeneous polynomials (there is a relation between their degrees).
Then, in section 4, we then compute locally the Poisson cohomology of Poisson structures of this type.
A generalisation of these results to the $`n`$-vectors on an $`n`$-dimensional manifold can be found in \[Mo2\].
## 2 Quasi-homogeneity
Throughout this text, $`๐`$ will indicate the field $``$ or $``$.
Let $`(\omega _1,\omega _2)^{}\times ^{}`$. We denote by $`W`$ the vector field $`\omega _1x\frac{}{x}+\omega _2y\frac{}{y}`$ on $`๐^2`$. Now, let $`T`$ be a non zero p-vector ($`\mathrm{p}\{0,1,2\}`$). We will say that $`T`$ is quasihomogeneous with weights $`\omega _1,\omega _2`$ and of (quasi)degree $`\mathrm{d}`$ if
$$_WT=\mathrm{d}T,$$
where $`_W`$ indicates the Lie derivative with respect to $`W`$. This condition can be written $`[W,T]=\mathrm{d}T`$ where \[.,.\] indicates the Schouten bracket. Note that $`T`$ is then polynomial.
If $`f`$ is a quasihomogeneous polynomial of degree d then $`\mathrm{d}=i\omega _1+j\omega _2`$ with $`(i,j)^2`$; therefore, an integer is not necessarily the quasidegree of a polynomial. If $`f๐\left[[x,y]\right]`$, we can write $`f=_{i=0}^{\mathrm{}}f_i`$ with $`f_i`$ quasihomogeneous of degree $`i`$ (we adopt the convention that $`f_i=0`$ if $`i`$ is not a quasidegree); $`f`$ is said to be of order d ($`\text{ord}(f)=\mathrm{d}`$) if all of its monomials have degree d or higher. For more details consult \[AGV\].
It is important to notice that $`\frac{}{x}`$ is a quasihomogeneous vector field of degree $`\omega _1`$ (in the same way $`\mathrm{deg}(\frac{}{y})=\omega _2`$); the minimal degree of a vector field is $`\mathrm{max}(\omega _1,\omega _2)`$. Note also that an integer can be the quasidegree of a vector field without being the quasidegree of a polynomial. Finally, note that $`\frac{}{x}\frac{}{y}`$ is quasihomogeneous of degree $`\omega _1\omega _2`$.
## 3 Local models of Poisson structures in dimension 2
In the reference \[A\], Arnold gives a list of normal forms for Poisson structures on a neighbourhood of $`(0,0)`$ in $`๐^2`$. In this section, we recall Arnoldโs theorem and we give the idea of a proof which is similar to Arnoldโs (the approach is a little bit different). For more details on this proof, consult \[Mo1\].
The particularity of the dimension two is that any 2-vector on a 2-dimensional manifold is a Poisson structure. For an introduction to Poisson structures, consult \[CW\] or \[V\].
The problem is the following: given $`\mathrm{\Pi }=F\frac{}{x}\frac{}{y}`$, a germ at 0 of Poisson structures on $`๐^2`$, we want to simplify its expression via a suitable local change of coordinates.
Notations : We denote by $`(๐^2)`$ (resp. $`๐ณ(๐^2),๐ฑ(๐^2)`$) the vector space of germs at (0,0) of (holomorphic if $`๐=`$ , analytic or $`๐^{\mathrm{}}`$ if $`๐=`$) functions (resp. vector fields, 2-vectors). We also denote by $`Diff_0(๐^2)`$ the group of local diffeomorphisms at (0,0) sending (0,0) to itself. Finally, $`_t(๐^2)`$ ($`๐ณ_t\left(๐^2\right),๐ฑ_t\left(๐^2\right)`$) indicates the space of germs depending differentiably on $`t`$.
Two germs $`\mathrm{\Pi }=f\frac{}{x}\frac{}{y}`$ and $`\mathrm{\Lambda }=g\frac{}{x}\frac{}{y}`$ are called equivalent if there exists $`\phi Diff_0(๐^2)`$ satisfying $`\phi _{}\mathrm{\Pi }=\mathrm{\Lambda }`$. This condition yields $`g\phi =(Jac\phi )f`$ where $`Jac\phi `$ indicates the Jacobian of $`\phi `$.
Two germs $`f`$ and $`g`$ are said to be R-equivalent if there exists $`\phi Diff_0(๐^2)`$ satisfying $`g\phi =f`$.
Actually, it is not possible to give normal forms for every Poisson structures. We only study Poisson structures determined by the germs of functions $`F`$ whose R-orbit is โinterestingโ enough. We will speak about it later.
The splitting theorem (\[W\]) allows us to assume that $`\mathrm{\Pi }_{(0,0)}=0`$. Moreover, it is quite easy to show that, if $`F`$ is regular at 0, then $`\mathrm{\Pi }`$ is, up to a change of coordinates, the germ $`x\frac{}{x}\frac{}{y}`$.
###### Remark 3.1
It is important to note that if two germs $`f`$ and $`g`$ are R-equivalent, then the germ $`\mathrm{\Pi }=f\frac{}{x}\frac{}{y}`$ will be equivalent to the germ of a Poisson structure of type $`ga\frac{}{x}\frac{}{y}`$ where $`a(0,0)0`$.
Now, we consider germs of Poisson structures of type
$$\mathrm{\Pi }=fa\frac{}{x}\frac{}{y},$$
where $`f`$ vanishes and is singular at $`(0,0)`$, and $`a(0,0)0`$.
Moreover, we suppose that $`f`$ is a quasihomogeneous polynomial of degree $`\mathrm{d}>0`$ with respect to $`W=\omega _1x\frac{}{x}+\omega _2y\frac{}{y}`$ ($`\omega _1`$ and $`\omega _2`$ are positive integers). This additional assumption will be justified later.
Here and throughout, the quasihomogeneity will be with respect to $`W`$.
Arnoldโs theorem is the following:
###### Theorem 3.2
\[A\] Up to a multiplicative constant, $`\mathrm{\Pi }`$ is equivalent to the germ of a Poisson structure of type $`f(1+h)\frac{}{x}\frac{}{y}`$ where $`h`$ is a quasihomogeneous polynomial of degree $`\mathrm{d}\omega _1\omega _2`$ (on condition that $`\mathrm{d}\omega _1\omega _2`$ is a quasidegree, otherwise this term disappears).
It is possible to show (see \[Mo1\]) that $`\mathrm{\Pi }`$ is equivalent to a germ of Poisson structures of type $`f(1+h+R)\frac{}{x}\frac{}{y}`$ where $`\mathrm{ord}\left(j_0^{\mathrm{}}(R)\right)>\mathrm{d}\omega _1\omega _2`$ ($`j_0^{\mathrm{}}(R)`$ indicates the $`\mathrm{}`$-jet of R at (0,0)) and $`h`$ is a quasihomogeneous polynomial of degree $`\mathrm{d}\omega _1\omega _2`$.
Thus, in order to prove the former theorem, we have to โremoveโ the term $`R`$. We are going to use Moserโs path method. For $`t`$, we put $`\mathrm{\Pi }_t=f(1+h+tR)\frac{}{x}\frac{}{y}`$ and we try to prove the existence of $`X_t๐ณ_t(๐^2)`$ satisfying $`[X_t,\mathrm{\Pi }_t]=\frac{\mathrm{d}\mathrm{\Pi }_t}{\mathrm{d}t}`$.
Actually we will look for an $`X_t`$ of type $`\alpha _tW`$ with $`\alpha _t_t(๐^2)`$.
Then, if we put $`R_t=\frac{R}{1+h+tR}`$ and $`\lambda _t=\mathrm{d}\omega _1\omega _2+\frac{W.(h+tR)}{1+h+tR}`$, it is sufficient to prove the existence of $`\alpha _t`$ in $`_t(๐^2)`$, such that
$$W.\alpha _t\lambda _t\alpha _t=R_t(E).$$
Let us note two things :
$``$ if $`\mathrm{\Pi }`$ is analytic ($`๐^{\mathrm{}}`$) then so are $`R_t`$ and $`\lambda _t`$
$``$ if $`\mathrm{d}\omega _1\omega _2`$ is a quasidegree, then $`\mathrm{ord}\left(j_0^{\mathrm{}}(R_t)\right)>\mathrm{d}\omega _1\omega _2`$.
Now, we just have to show that there exists $`\alpha _t`$ satisfying $`(E)`$.
Resolubility of equation $`(E)`$ : The results we give here will be useful in the computation of the Poisson cohomology. That is why we are going to give quite detailed proofs of them.
We can write $`\lambda _t=(\mathrm{d}\omega _1\omega _2)+\mu _t`$ where $`\mu _t_t(๐^2)`$ satisfies $`\mu _t(0,0)=0`$. In order to show that $`(E)`$ admits a solution :
1- we prove that there exists $`\beta _t_t(๐^2)`$ satisfying $`W.\beta _t\mu _t\beta _t=0`$ with $`\beta _t(0,0)0`$
2- we prove that there exists $`\gamma _t_t(๐^2)`$ satisfying $`W.\gamma _t(\mathrm{d}\omega _1\omega _2)\gamma _t=\frac{R_t}{\beta _t}`$
3- $`\alpha _t:=\beta _t\gamma _t`$ will then be a solution of $`(E)`$.
1- In order to show the first claim, we need the following result whose proof can be found in \[R\].
###### Theorem 3.3
Let $`X_t`$ be an element in $`๐ณ_t(๐^2)`$ having an isolated singularity at (0,0). Moreover, suppose that the eigenvalues of its linear component at (0,0) do not vanish. Take $`h_t`$ in $`๐_t^{\mathrm{}}(^2)`$ flat at (0,0). Then there exists $`g_t๐_t^{\mathrm{}}(^2)`$ flat at (0,0) satisfying $`X_t.g_t=h_t`$ for any $`t`$.
We also need the following lemma.
###### Lemma 3.4
If $`T_t_t(๐^2)`$ satisfies $`T_t(0,0)=0`$, then there exists $`\nu _t_t(๐^2)`$ such that $`W.\nu _t=T_t`$.
Proof of the lemma :
Formal case : Assume that $`T_t๐_t\left[[x,y]\right]`$; we have $`T_t=_{i>0}T_t^{(i)}`$ where $`T_t^{(i)}`$ is quasihomogeneous of degree $`i`$. If we put $`\nu _t=_{i>0}\frac{T_t^{(i)}}{i}`$ we get $`W.\nu _t=T_t`$.
Analytical case : Assume that $`T_t`$ is analytic at (0,0). Imitate the former proof noting that, so defined, $`\nu _t`$ is analytic at (0,0).
$`๐^{\mathrm{}}`$ case : Let $`\stackrel{~}{T_t}=j_0^{\mathrm{}}(T_t)`$ and $`\stackrel{~}{\epsilon _t}_t\left[[x,y]\right]`$ be such that $`W.\stackrel{~}{\epsilon _t}=\stackrel{~}{T_t}`$. Borelโs theorem ensures the existence of $`\epsilon _t๐^{\mathrm{}}(^2)`$ such that $`j_0^{\mathrm{}}(\epsilon _t)=\stackrel{~}{\epsilon _t}`$. Thus $`W.\epsilon _t=T_t+m_t`$ where $`m_t`$ is flat at (0,0). Let $`n_t`$ be flat at (0,0) and such that $`W.n_t=m_t`$ ($`n_t`$ exists by theorem 3.3); $`\nu _t=\epsilon _t+n_t`$ suits. $`\mathrm{}`$
Consequently, to prove 1-, we put $`\beta _t=exp\nu _t`$, where $`\nu _t`$ satisfies $`W.\nu _t=\mu _t`$.
2- Note first that if $`\mathrm{d}\omega _1\omega _2`$ is a quasidegree (for polynomials), then there exists $`(i,j)`$ in $`^2`$ such that $`\mathrm{d}\omega _1\omega _2=i\omega _1+j\omega _2`$ if not, $`\mathrm{d}\omega _1\omega _2=i\omega _1\omega _2`$ (or $`\omega _1+i\omega _2`$) with $`i`$. The following lemma will prove the second claim.
###### Lemma 3.5
Let $`k`$ and $`l`$ be in $``$ and $`T_t_t(๐^2)`$ with $`\mathrm{ord}((j_0^{\mathrm{}}(T_t))>k\omega _1+l\omega _2`$. Then there exists $`\gamma _t_t(๐^2)`$ satisfying $`W.\gamma _t(k\omega _1+l\omega _2)\gamma _t=T_t`$.
Let $`k`$ and $`T_t_t(๐^2)`$; then there exists $`\gamma _t_t(๐^2)`$ satisfying
$`W.\gamma _t(k\omega _1\omega _2)\gamma _t=T_t`$.
Proof : i) We use an induction :
For $`k=l=0`$: see lemma 3.4.
Now, assume that i) is true for $`(k,l)^2`$. We are going to show that it is true for $`k+1`$ and $`l`$ (for $`k`$ and $`l+1`$ the proof is the same).
Let $`T_t_t(๐^2)`$ with $`\text{ord}\left(j_0^{\mathrm{}}(T_t)\right)>(k+1)\omega _1+l\omega _2`$ and $`\delta _t_t(๐^2)`$ verifying $`W.\delta _t(k\omega _1+l\omega _2)\delta _t=\frac{T_t}{x}`$. Then we define $`\gamma _t`$ by $`\gamma _t(x,y)=_0^x\delta _t(u,y)๐u`$ for $`(x,y)`$ in a neighbourhood of (0,0). An easy computation shows that $`W.\gamma _t((k+1)\omega _1+l\omega _2)\gamma _t=T_t`$.
ii) We use again an induction :
For $`k=0`$: we know that there exists $`\delta _t_t(๐^2)`$ such that $`W.\delta _t=_0^yT_t(x,u)๐u`$. If we put $`\gamma _t=\frac{\delta _t}{y}`$ then we get $`W.\gamma _t+\omega _2\gamma _t=T_t`$.
The end of the proof can be achieved as in i). $`\mathrm{}`$
A list of normal forms : We recall that a germ of Poisson structures on $`๐^n`$ is determined by the choice of a germ of functions.
We consider a germ $`\mathrm{\Pi }=f\frac{}{x}\frac{}{y}`$, where $`f`$ vanishes and is singular at $`(0,0)`$. We suppose, in addition, that the germ $`f`$ is of finite codimension. It means that the vector space $`Q_f=(๐^2)/I_f`$ ($`I_f`$ is the ideal spanned by $`\frac{f}{x}`$ and $`\frac{f}{y}`$) is of finite dimension.
Why do we suppose that $`f`$ is of finite codimension? In fact, one can see $`I_f`$ as the tangent space of the orbit of $`f`$ (with respect to the R-equivalence). Thus, the finite-codimensional germs are those whose orbit is โbigโ enough.
###### Remark 3.6
It is important to note the following fact:
According to Tougeronโs theorem (see for instance \[AGV\]), if $`f`$ is of finite codimension, then $`f`$ is R-equivalent to its $`k`$-jets for $`k`$ sufficiently large. The set $`f^1(\{0\})`$ is then, from the topological point of view, the same as the set of zeroes of a polynomial. Therefore, if $`g`$ is a germ at 0 of functions which satisfies $`fg=0`$, then $`g=0`$.
Moreover, we suppose that the germ $`f`$ is simple. It means that a sufficiently small neighbourhood (with respect to Whitneyโs topology; see \[AGV\]) of $`f`$ intersects only a finite number of R-orbits. Simple germs are those who present a certain kind of stability under deformation.
Note that simple germs are necessarily of finite codimension (see for instance \[AGV\]).
We have a classification of such germs in the following theorem.
###### Theorem 3.7
\[AGV\] Simple germs at (0,0) of functions are given, up to R-equivalence, in the following list:
| $`A_kk1`$ | $`D_kk4`$ | $`E_6`$ | $`E_7`$ | $`E_8`$ |
| --- | --- | --- | --- | --- |
| $`x^2\pm y^{k+1}`$ | $`x^2y\pm y^{k1}`$ | $`x^3\pm y^4`$ | $`x^3+xy^3`$ | $`x^3+y^5`$ |
If $`๐=`$ (or if $`k`$ is even in the real $`A_k`$ case), then the symbol $`\pm `$ disappears.
It is important to note that these models are quasihomogeneous polynomials.
Now, applying theorem 3.2 to these models, we can state the following theorem.
###### Theorem 3.8
\[A\] Let $`f`$ be a simple germ at (0,0) of finite codimension. Suppose that $`f`$ has at (0,0) a critical point with critical value 0. Then, the germ $`\mathrm{\Pi }=f\frac{}{x}\frac{}{y}`$ is equivalent, up to a multiplicative constant, to a germ of type $`g\frac{}{x}\frac{}{y}`$, where $`g`$ is in the following list:
$`A_{2p}`$ $`:`$ $`x^2+y^{2p+1}p1`$
$`A_{2p1}^\pm `$ $`:`$ $`(x^2\pm y^{2p})(1+\lambda y^{p1})p1`$
$`D_{2p}^\pm `$ $`:`$ $`(x^2\pm y^{2p})(1+\lambda y^{p1})p2`$
$`D_{2p+1}`$ $`:`$ $`(x^2y+y^{2p})(1+\lambda x)p2`$
$`E_6`$ $`:`$ $`x^3+y^4`$
$`E_7`$ $`:`$ $`(x^3+xy^3)(1+\lambda y^2)`$
$`E_8`$ $`:`$ $`x^3+y^5`$
If $`๐=`$, the symbol $`\pm `$ disappears.
## 4 Poisson cohomology
In this section we compute the Poisson cohomology of some Poisson structures. In fact, we work locally and we study the โgermifiedโ Poisson cohomology. This means that we work with germs of Poisson structures, functions, vector fields, 2-vectorsโฆ
We recall that $`(๐^2)`$ ($`๐ณ(๐^2)`$, $`๐ฑ(๐^2)`$) indicates the space of germs at 0 of functions (vector fields, 2-vectors).
Let $`\mathrm{\Pi }`$ be a germ of Poisson structure on $`๐^2`$. We have then the complex
$$0\stackrel{\delta _0}{}(๐^2)\stackrel{\delta _1}{}๐ณ(๐^2)\stackrel{\delta _2}{}๐ฑ(๐^2)\stackrel{\delta _3}{}0$$
where $`\delta _0=0`$, $`\delta _3=0`$, $`\delta _1(g)=[g,\mathrm{\Pi }]`$ and $`\delta _2(X)=[X,\mathrm{\Pi }]`$ (\[.,.\] indicates the Schouten bracket).
We will denote by $`Z^i(\mathrm{\Pi })=\text{Ker }\delta _{i+1}`$, $`B^i(\mathrm{\Pi })=\text{Im }\delta _i`$ and $`H^i(\mathrm{\Pi })=Z^i(\mathrm{\Pi })/B^i(\mathrm{\Pi })`$.
If we assume that $`\mathrm{\Pi }=F\frac{}{x}\frac{}{y}`$ where $`F(๐^2)`$ then, for $`g(๐^2)`$ , we have
$$\delta _1(g)=F\frac{g}{y}\frac{}{x}F\frac{g}{x}\frac{}{y}.$$
We will denote by $`X_g`$ this vector field (it is the Hamiltonian of $`g`$ with respect to $`\mathrm{\Pi }`$) and $`H_g`$ the vector field $`\frac{g}{y}\frac{}{x}\frac{g}{x}\frac{}{y}`$.
On the other hand, for $`X๐ณ(๐^2)`$, we have
$$\delta _2(X)=(X.F(divX)F)\frac{}{x}\frac{}{y}.$$
We will denote $`H^2(F)`$ the space $`(๐^2)/\{X.F(divX)F|X๐ณ(๐^2)\}`$. This space is clearly isomorphic to $`H^2(\mathrm{\Pi })`$.
Actually, we will compute the cohomology of Poisson structures of a particular type.
Let $`(\omega _1,\omega _2)^{}\times ^{}`$. Here and throughout, the quasihomogeneity will be understood as being in the sense of $`(\omega _1,\omega _2)`$ ($`W`$ will again indicate the vector field $`\omega _1x\frac{}{x}+\omega _2y\frac{}{y}`$).
We take a quasihomogeneous polynomial $`f`$ of degree d and we assume that $`f`$ is a germ at 0 of finite codimension $`c`$ (recall that it means that the vector space $`Q_f=(๐^2)/I_f`$, where $`I_f`$ is the ideal spanned by $`\frac{f}{x}`$ and $`\frac{f}{y}`$, is of dimension $`c`$). We also give us a quasihomogeneous polynomial $`h`$ of degree $`\mathrm{d}\omega _1\omega _2`$ (if $`\mathrm{d}\omega _1\omega _2>0`$).
Now we consider two germs of Poisson structures
$$\mathrm{\Pi }_0=f\frac{}{x}\frac{}{y}\text{ and }\mathrm{\Pi }=f(1+h)\frac{}{x}\frac{}{y}$$
and we are going to compute the cohomology of these Poisson structures.
In the former section, we saw that the โmost interestingโ Poisson structures are of this type.
Note that in the sequel,we do not suppose that the germ $`f`$ is simple.
In our computation, it is very important to assume that $`f`$ is of finite codimension (see the role played by the second claim of lemma 4.2, and the remark 4.8).
It is easy to see that, since $`f`$ is of finite codimesion, the spaces $`H^0(\mathrm{\Pi }_0)`$ and $`H^0(\mathrm{\Pi })`$ are isomorphic to $`๐`$ (see remark 3.6).
###### Remark 4.1
It is important to note that, since $`f`$ is quasihomogeneous, the computation of $`H^{}(\mathrm{\Pi }_0)`$ can be done โdegree by degreeโ. For instance, if $`_iX^{(i)}`$ is the $`\mathrm{}`$-jet of $`X`$ and if $`X`$ is in $`Z^1(\mathrm{\Pi }_0)`$ (resp. $`B^1(\mathrm{\Pi }_0)`$) then $`X^{(i)}`$ is also in $`Z^1(\mathrm{\Pi }_0)`$ (resp. $`B^1(\mathrm{\Pi }_0)`$) for each $`i`$. Moreover, if $`X`$ is polynomial, then $`X`$ is in $`Z^1(\mathrm{\Pi }_0)`$ (resp. $`B^1(\mathrm{\Pi }_0)`$) if and only if each of its quasihomogeneous components is in $`Z^1(\mathrm{\Pi }_0)`$ (resp. $`B^1(\mathrm{\Pi }_0)`$). We have the same properties for $`B^2(\mathrm{\Pi }_0)`$.
The computation of the cohomology of $`\mathrm{\Pi }`$ does not present this property.
The following result will be useful in the sequel.
###### Lemma 4.2
Let $`X`$ be in $`๐ณ(๐^2)`$.
1- If $`divX=0`$, then there exists $`g(๐^2)`$ such that $`X=H_g`$.
2- If $`X.f=0`$, then $`X=\alpha H_f`$ with $`\alpha (๐^2)`$.
Proof : We can write $`X=A\frac{}{x}+B\frac{}{y}`$. We consider the 1-form $`\omega =Bdx+Ady`$.
1- If $`divX=0`$ then $`d\omega =0`$, which implies that $`\omega =dg`$ with $`g(๐^2)`$, and so $`X=H_g`$.
2- If $`X.f=0`$ then $`df\omega =0`$. Since $`f`$ has finite codimension, de Rhamโs division theorem (see \[dR\] or \[M\]) enables us to conclude. $`\mathrm{}`$
### 4.1 Computation of $`H^1`$
Computation of $`H^1(\mathrm{\Pi }_0)`$ :
###### Lemma 4.3
Let $`XZ^1(\mathrm{\Pi }_0)`$. Then there exists $`\alpha (๐^2)`$ such that
$`X=\alpha H_f+\frac{divX}{\mathrm{d}}W`$.
Proof : Direct application of lemma 4.2.
The main idea in the computation of the space $`H^1(\mathrm{\Pi }_0)`$, is to show that every 1-cocycle whose $`\mathrm{}`$-jet has a sufficiently large order is a cobord.
###### Lemma 4.4
Let $`X๐ณ(๐^2)`$ be such that $`\mathrm{ord}\left(j_0^{\mathrm{}}(X)\right)>\mathrm{d}\omega _1\omega _2`$.
If $`XZ^1(\mathrm{\Pi }_0)`$ then $`XB^1(\mathrm{\Pi }_0)`$.
Proof : first case : $`divX=0`$. We then show that $`f`$ divides X.
Since $`X.f(divX)f=0`$, we have $`X.f=0`$ and then $`X=\gamma H_f`$ (lemma 4.2) with $`\gamma (๐^2)`$. Note that, since $`\text{ord}\left(j_0^{\mathrm{}}(X)\right)>\mathrm{d}\omega _1\omega _2`$ and $`H_f`$ is quasihomogeneous of degree $`\mathrm{d}\omega _1\omega _2`$, we have $`\text{ord}\left(j_0^{\mathrm{}}(\gamma )\right)>0`$.
We prove that $`f`$ divides $`\gamma `$. Let $`\mu (๐^2)`$ be such that $`W.\mu =\gamma `$ ($`\mu `$ exists according to lemma 3.4). Note that, if we write the $`\mathrm{}`$-jet in the relation $`W.\mu =\gamma `$, we see that the order of the $`\mathrm{}`$-jet of $`\mu `$ is also strictly positive (because $`W`$ is quasihomogeneous of degree 0). Therefore, the order of the $`\mathrm{}`$-jet of $`H_f.\mu `$ is strictly larger than $`\mathrm{d}\omega _1\omega _2`$, because $`H_f`$ is quasihomogeneous of degree $`\mathrm{d}\omega _1\omega _2`$.
We have
$$H_f.(W.\mu )=W.(H_f.\mu )+[H_f,W].\mu =W.(H_f.\mu )+((\mathrm{d}\omega _1\omega _2)H_f).\mu $$
(because $`H_f`$ is of degree $`\mathrm{d}\omega _1\omega _2`$).
Since $`H_f.(W.\mu )=H_f.\gamma =divX=0`$, we have
$$W.(H_f.\mu )=(\mathrm{d}\omega _1\omega _2)H_f.\mu $$
and so $`H_f.\mu `$ is either 0 or quasihomogeneous of degree $`\mathrm{d}\omega _1\omega _2`$.
Thus, as $`\text{ord}\left(j_0^{\mathrm{}}(H_f.\mu )\right)>\mathrm{d}\omega _1\omega _2`$, we have $`H_f.\mu =0`$.
Now, $`H_\mu .f=H_f.\mu =0`$ so there exists $`\nu (๐^2)`$ such that $`\frac{\mu }{x}=\nu \frac{f}{x}`$ and $`\frac{\mu }{y}=\nu \frac{f}{y}`$ (lemma 4.2). Therefore, $`W.\mu =\nu W.f`$, that is $`\gamma =\nu (\mathrm{d}\times f)`$.
We deduce that $`X=fZ`$ with $`Z๐ณ(๐^2)`$.
Finally, since $`XZ^1(\mathrm{\Pi }_0)`$, $`divZ=0`$ and then $`Z=H_g`$ for some $`g(๐^2)`$ (lemma 4.2). Hence $`X=fH_g=X_g`$.
Second case : $`divX0`$. If we find $`\beta (๐^2)`$ such that $`divX=divX_\beta `$, then the 1-cocycle $`XX_\beta `$ satisfies $`div(XX_\beta )=0`$, which implies (see the first case) that $`X=X_\beta +X_\epsilon `$ where $`\epsilon (๐^2)`$. Since $`divX_\beta =H_\beta .f=H_f.\beta `$, we are looking for $`\beta `$ such that $`H_f.\beta =divX`$.
We have $`X=\alpha H_f+\frac{divX}{\mathrm{d}}W`$ with $`\alpha (๐^2)`$ (lemma 4.3) so that $`divX`$ satisfies the equation
$$W.(divX)(\mathrm{d}\omega _1\omega _2)divX=\mathrm{d}\times H_f.\alpha .$$
Note that, if we write the $`\mathrm{}`$-jets at 0 in the relation $`X=\alpha H_f+\frac{divX}{\mathrm{d}}W`$, we see that the order of the $`\mathrm{}`$-jet of $`\alpha `$ is strictly positive so, according to lemma 3.4 we can take $`\beta (๐^2)`$ such that $`W.\beta =\mathrm{d}\times \alpha `$ (the order of the $`\mathrm{}`$-jet of $`\beta `$ is also strictly positive). Since
$$W.(H_f.\beta )=H_f.(W.\beta )+[W,H_f].\beta =\mathrm{d}(H_f.\alpha )+\left((\mathrm{d}\omega _1\omega _2)H_f\right).\beta ,$$
we have $`W.(divX+H_f.\beta )=(\mathrm{d}\omega _1\omega _2)(divX+H_f.\beta )`$.
We deduce that $`divX+H_f.\beta `$ is either 0 or quasihomogeneous of degree $`\mathrm{d}\omega _1\omega _2`$. Therefore, $`divX=H_f.\beta `$ (because $`\text{ord}\left(j_0^{\mathrm{}}(divX+H_f.\beta )\right)>\mathrm{d}\omega _1\omega _2`$). $`\mathrm{}`$
For $`XZ^1(\mathrm{\Pi }_0)`$, we denote by $`[X]_{_{\mathrm{\Pi }_0}}`$ its class modulo $`B^1(\mathrm{\Pi }_0)`$. We also denote by $`\{e_1,\mathrm{},e_r\}`$ a basis of the vector space of quasihomogeneous polynomials of degree $`\mathrm{d}\omega _1\omega _2`$ (in fact, in order to obtain a vector space, we have to add 0 to the set of quasihomogeneous polynomials of degree $`\mathrm{d}\omega _1\omega _2`$).
###### Theorem 4.5
The family $`\{[H_f]_{_{\mathrm{\Pi }_0}},[e_1W]_{_{\mathrm{\Pi }_0}},\mathrm{},[e_rW]_{_{\mathrm{\Pi }_0}}\}`$ is a basis of $`H^1(\mathrm{\Pi }_0)`$. In particular, $`H^1(\mathrm{\Pi }_0)`$ is a finite-dimensional vector space of dimension $`r+1`$.
Proof : First we prove that $`H^1(\mathrm{\Pi }_0)`$ is spanned by this family. Lemma 4.4 says that every $`X`$ in $`Z^1(\mathrm{\Pi }_0)`$ is cohomologous to a polynomial vector field of maximum degree $`\mathrm{d}\omega _1\omega _2`$. Indeed, if $`XZ^1(\mathrm{\Pi }_0)`$ then $`j_0^{\mathrm{d}\omega _1\omega _2}(X)`$ is also in $`Z^1(\mathrm{\Pi }_0)`$ ($`j_0^{\mathrm{d}\omega _1\omega _2}(X)`$ indicates the jet of degree $`\mathrm{d}\omega _1\omega _2`$ of $`X`$ at 0). Thus, $`Xj_0^{\mathrm{d}\omega _1\omega _2}(X)`$ is in $`Z^1(\mathrm{\Pi }_0)`$ and the order of its $`\mathrm{}`$-jet at 0 is strictly higher than $`\mathrm{d}\omega _1\omega _2`$.
Therefore, using remark 4.1, we can assume that $`X`$ is quasihomogeneous of degree $`\mathrm{d}\omega _1\omega _2`$ or lower.
\- If $`XZ^1(\mathrm{\Pi }_0)`$ is quasihomogeneous with $`\mathrm{deg}X<\mathrm{d}\omega _1\omega _2`$ then $`X=0`$. Indeed, according to lemma 4.3, we have $`X=\frac{divX}{\mathrm{d}}W`$, and so
$$divX=\frac{divX}{\mathrm{d}}divW+W.\left(\frac{divX}{d}\right),$$
which implies that $`(\mathrm{d}\omega _1\omega _2\mathrm{deg}X)divX=0`$.
\- Let $`XZ^1(\mathrm{\Pi }_0)`$ be quasihomogeneous of degree $`\mathrm{d}\omega _1\omega _2`$. We have (lemma 4.3) $`X=\alpha H_f+\frac{divX}{\mathrm{d}}W`$ where $`\alpha ๐`$ and $`divX`$ is a quasihomogeneous polynomial of degree $`\mathrm{d}\omega _1\omega _2`$.
Therefore the family generates $`H^1(\mathrm{\Pi }_0)`$.
Now, we prove that this family is free. Suppose that $`_i\lambda _ie_iW+\alpha H_fB^1(\mathrm{\Pi })`$ where $`\alpha ,\lambda _1,...,\lambda _r`$ are scalars. Then $`_i\lambda _ie_iW+\alpha H_f=0`$. Indeed, if $`g`$ is a quasihomogeneous polynomial, then $`\mathrm{deg}X_g=\mathrm{deg}H_f+\mathrm{deg}g=\mathrm{d}\omega _1\omega _2+\mathrm{deg}g`$, which is strictly larger than $`\mathrm{d}\omega _1\omega _2`$ as soon as $`g0`$.
Consequently, $`div\left(_i\lambda _ie_iW+\alpha H_f\right)=0`$ i.e. $`_i\lambda _ie_i=0`$. We deduce that $`\lambda _1,\mathrm{},\lambda _r`$ are 0, and so $`\alpha =0`$. $`\mathrm{}`$
Computation of $`H^1(\mathrm{\Pi })`$ :
If $`XZ^1(\mathrm{\Pi })`$ we denote by $`[X]__\mathrm{\Pi }`$ its class modulo $`B^1(\mathrm{\Pi })`$.
###### Theorem 4.6
$`\{[(1+h)H_f]__\mathrm{\Pi },[(1+h)e_1W]__\mathrm{\Pi },\mathrm{},[(1+h)e_rW]__\mathrm{\Pi }\}`$ is a basis of $`H^1(\mathrm{\Pi })`$. In particular, $`H^1(\mathrm{\Pi })H^1(\mathrm{\Pi }_0)`$.
Proof : It is sufficient to notice that $`XZ^1(\mathrm{\Pi })`$ (resp. $`XB^1(\mathrm{\Pi })`$) if and only if $`\frac{X}{1+h}Z^1(\mathrm{\Pi }_0)`$ (resp. $`B^1(\mathrm{\Pi }_0)`$). $`\mathrm{}`$
### 4.2 Computation of $`H^2`$
Computation of $`H^2(\mathrm{\Pi }_0)`$ :
###### Lemma 4.7
Let $`g`$ be a germ at 0 of functions on $`๐^2`$.
1. If the $`\mathrm{}`$-jet at $`0`$ of $`g`$ does not contain a component of degree $`2\mathrm{d}w_1w_2`$ then
$$gB^2(f)gI_f.$$
2. If $`g`$ is quasihomogeneous of degree $`2\mathrm{d}w_1w_2`$ then
$$gB^2(f)gI_f.$$
Proof : If $`g=X.f(divX)fB^2(f)`$ where $`X๐ณ(๐^2)`$ then , if we put
$$Y=X\frac{divX}{\mathrm{d}}W,$$
we have $`g=Y.f`$
This proves the second claim and the first part of the first one.
Now, we prove the converse of the first claim: we assume that $`gI_f`$ (and the $`\mathrm{}`$-jet of $`g`$ does not contain a component of degree $`2\mathrm{d}w_1w_2`$) and we are going to show that $`gB^2(f)`$.
Formal case : Let $`g={\displaystyle \underset{i0}{}}g^{(i)}๐[[x,y]]`$ and $`X={\displaystyle \underset{i\mathrm{d}\mathrm{max}(\omega _1,\omega _2)}{}}X^{(i\mathrm{d})}`$ (with $`g^{(i)}`$ of degree $`i`$ and $`X^{(i\mathrm{d})}`$ of degree $`i\mathrm{d}`$) such that $`g=X.f`$. Note that $`X^{(\mathrm{d}\omega _1\omega _2)}=0`$.
If we put
$$Y=X+\underset{i2\mathrm{d}\omega _1\omega _2}{}\frac{divX^{(i\mathrm{d})}}{2d\omega _1\omega _2i}W,$$
a direct computation gives $`Y.f(divY)f=X.f=g`$.
Analytical case : If $`X`$ is analytic at (0,0), then $`divX`$ is analytic too and since
$`lim_{i+\mathrm{}}\frac{1}{2\mathrm{d}\omega _1\omega _2i}=0`$ the vector field defined above is also analytic in (0,0).
$`๐^{\mathrm{}}`$ case : Let us denote by $`\stackrel{~}{g}=j_0^{\mathrm{}}(g)`$ and $`\stackrel{~}{X}=j_0^{\mathrm{}}(X)`$.
If we write the $`\mathrm{}`$-jets in the relation $`g=X.f`$, we get $`\stackrel{~}{g}=\stackrel{~}{X}.f`$. Thus, there exists a formal vector field $`\stackrel{~}{Y}`$ such that
$$\stackrel{~}{g}=\stackrel{~}{Y}.f(div\stackrel{~}{Y})f.$$
Let $`Y`$ be a $`๐^{\mathrm{}}`$ vector field such that $`\stackrel{~}{Y}=j_0^{\mathrm{}}(Y)`$. Since $`Y.f(divY)f`$ and $`g`$ have the same $`\mathrm{}`$-jet, this vector field satisfies
$$Y.f(divY)f=g+\epsilon $$
where $`\epsilon `$ is flat at 0.
Now, since $`Y.f(divY)fI_f`$ (see the beginning of the proof), $`\epsilon `$ is in $`I_f`$ so that $`\epsilon =P.f`$ where $`P`$ is a flat vector field. According to lemma 3.5, there exists $`\alpha (๐^2)`$ such that $`W.\alpha (\mathrm{d}\omega _1\omega _2)\alpha =divP`$. Consequently, setting $`Z=P+\alpha W`$, we have $`Z.f(divZ)f=\epsilon `$. $`\mathrm{}`$
###### Remark 4.8
1- This lemma is true even if $`f`$ is not of finite codimension.
2- This lemma gives $`B^2(f)I_f`$. Thus, there is a surjection from $`H^2(f)`$ onto $`Q_f`$. Therefore, if $`f`$ is not of finite codimension, then $`H^2(\mathrm{\Pi }_0)`$ is an infinite-dimensional vector space.
3- Finally, according to this lemma, if $`\xi I_f`$, then there exists a quasihomogeneous polynomial $`\overline{\xi }`$ of degree $`2\mathrm{d}\omega _1\omega _2`$ such that $`\xi +\overline{\xi }B^2(f)`$.
If $`g(๐^2)`$, $`[g]_{_{\mathrm{\Pi }_0}}`$ indicates its class modulo $`B^2(f)`$. Recall that $`\{e_1,\mathrm{},e_r\}`$ is a basis of the space of quasihomogeneous polynomials of degree $`\mathrm{d}\omega _1\omega _2`$. Finally, we denote by $`\{u_1,\mathrm{},u_c\}`$ a monomial basis of $`Q_f=(๐^2)/I_f`$ (for the existence of such a basis, see \[AGV\]).
###### Theorem 4.9
The family $`\{[e_1f]_{_{\mathrm{\Pi }_0}},\mathrm{},[e_rf]_{_{\mathrm{\Pi }_0}},[u_1]_{_{\mathrm{\Pi }_0}},\mathrm{},[u_c]_{_{\mathrm{\Pi }_0}}\}`$ is a basis of $`H^2(f)`$.
In particular, $`H^2(\mathrm{\Pi }_0)`$ is a finite-dimensional vector space of dimension $`r+c`$.
Proof :- This family generates $`H^2(f)`$:
Let $`g(๐^2)`$. We can write $`g=_{i=1}^c\lambda _iu_i+\xi `$ where $`\lambda _i๐`$ and $`\xi I_f`$. According to lemma 4.7 (and remark 4.8), we can write
$$g=\underset{i=1}{\overset{c}{}}\lambda _iu_i+\overline{g}\text{mod}B^2(f)$$
where $`\overline{g}`$ is a quasihomogeneous polynomial of degree $`2\mathrm{d}\omega _1\omega _2`$.
We can again write $`\overline{g}=_{i=1}^c\overline{\lambda _i}u_i\text{ mod }I_f`$ where $`\overline{\lambda _i}๐`$. Now, we know (see \[AGV\] p.200) that $`\mathrm{max}\{\mathrm{deg}u_1,\mathrm{},\mathrm{deg}u_c\}=2\mathrm{d}2\omega _12\omega _2`$ which is strictly lower than $`\mathrm{deg}\overline{g}`$. So, $`\overline{g}`$ is in $`I_f`$, i.e. $`\overline{g}=X.f`$ with $`X`$ quasihomogeneous of degree $`\mathrm{d}\omega _1\omega _2`$.
Therefore, $`\overline{g}=(divX)f+(X.f(divX)f)`$ with $`divX`$ quasihomogeneous of degree $`\mathrm{d}\omega _1\omega _2`$.
\- This family is free: Let $`g_1=_{i=1}^r\lambda _ie_i`$ and $`g_2=_{j=1}^c\mu _ju_j`$ with $`\lambda _i`$ and $`\mu _j`$ in $`๐`$ for any $`i`$ and $`j`$. We assume that $`g_1f+g_2B^2(f)`$. Since $`\mathrm{max}\{\mathrm{deg}u_1,\mathrm{},\mathrm{deg}u_c\}<\mathrm{deg}(g_1f)`$, $`g_1f`$ and $`g_2`$ are both in $`B^2(f)`$ (see remark 4.1).
On the one hand, we then have $`g_2I_f`$ which is possible only if $`\mu _1=\mathrm{}=\mu _c=0`$.
On the other hand, since $`g_1f=X.f(divX)f`$ for some quasihomogeneous vector field $`X`$ of degree $`\mathrm{d}\omega _1\omega _2`$, if $`Y`$ denotes the vector field $`\frac{g_1+divX}{\mathrm{d}}W`$ then $`(XY).f=0`$. Therefore (lemma 4.2), $`X=Y+\alpha H_f`$ with $`\alpha ๐`$. Hence $`X.f(divX)f=0`$, which implies $`\lambda _1=\mathrm{}=\lambda _r=0`$. $`\mathrm{}`$
Computation of $`H^2(\mathrm{\Pi })`$ :
###### Lemma 4.10
Let $`g(๐^2)`$. If the order of the $`\mathrm{}`$-jet of $`g`$ is larger than or equal to $`2\mathrm{d}\omega _1\omega _2`$, then there exists a quasihomogeneous polynomial $`\epsilon `$ of degree $`\mathrm{d}\omega _1\omega _2`$ such that $`g=\epsilon f\mathrm{mod}B^2(f+fh)`$.
Proof : According to theorem 4.9, we can write
$$\frac{g}{1+h}=\underset{i=1}{\overset{c}{}}\lambda _iu_i+\epsilon f+X.f(divX)f$$
where the $`\lambda _i`$ are in $`๐`$ and $`\epsilon `$ is quasihomogeneous of degree $`\mathrm{d}\omega _1\omega _2`$.
But since $`\text{ord}\left(j_0^{\mathrm{}}\left(\frac{g}{1+h}\right)\right)2\mathrm{d}\omega _1\omega _2>\mathrm{max}\{\mathrm{deg}u_1,\mathrm{},\mathrm{deg}u_c\}`$, if we write the $`\mathrm{}`$-jets in the former relation, we see that the $`\lambda _i`$ are zero and that the order of the $`\mathrm{}`$-jet of $`X`$ is higher than $`\mathrm{d}\omega _1\omega _2`$.
Now,
$$X.(f+fh)(divX)(f+fh)=g+f(X.h)\epsilon f(1+h).$$
We put $`\lambda =(\mathrm{d}\omega _1\omega _2)\left(1+\frac{h}{1+h}\right)`$. If $`X.h\epsilon h=0`$, then the lemma is shown. Now, if we suppose that $`X.h\epsilon h`$ is not 0, the order of its $`\mathrm{}`$-jet is strictly larger than $`\mathrm{d}\omega _1\omega _2`$, and so we can take $`\alpha (๐^2)`$ such that $`W.\alpha \lambda \alpha =\frac{X.h\epsilon h}{1+h}`$ (see the โResolubility of equation $`(E)`$โ in the former section).
If we put $`Z=X+\alpha W`$, we have $`Z.(f+fh)(divZ)(f+fh)+f\epsilon =g`$. $`\mathrm{}`$
###### Theorem 4.11
The family $`\{[e_1f]__\mathrm{\Pi },\mathrm{},[e_rf]__\mathrm{\Pi },[u_1]__\mathrm{\Pi },\mathrm{},[u_c]__\mathrm{\Pi }\}`$ is a basis of $`H^2(f+fh)`$. In particular, $`H^2(f+fh)H^2(f)`$ (the space $`H^2(\mathrm{\Pi })`$ is then of dimension $`r+c`$).
Proof :
\- This family generates $`H^2(f+fh)`$.
Given $`g(๐^2)`$, we have $`g=_{i=1}^c\lambda _{i,0}u_i+P_0f+X_0.f(divX_0)f`$ (theorem 4.9) where $`P_0`$ is a quasihomogeneous polynomial of degree $`\mathrm{d}\omega _1\omega _2`$, $`\lambda _{i,0}๐`$ for any $`i`$ and $`\text{ord}\left(j_0^{\mathrm{}}(X_0)\right)\text{max}(\omega _1,\omega _2)`$ (remember that $`\text{max}(\omega _1,\omega _2)`$ is the smallest degree for quasihomogeneous vector fields). Since
$$X_0.(f+fh)(divX_0)(f+fh)=g\underset{i=1}{\overset{c}{}}\lambda _{i,0}u_iP_0f+X_0.(fh)(divX_0)(fh),$$
we can write
$$g=\underset{i=1}{\overset{c}{}}\lambda _{i,0}u_i+P_0f(X_0.(fh)(divX_0)(fh))\text{mod}B^2(f+fh).$$
Now, we have $`X_0.(fh)(divX_0)fh=\lambda _{i,1}u_iP_1f+X_1.f(divX_1)f`$ where $`\lambda _{i,1}`$ is in $`๐`$ for any $`i`$, $`P_1`$ is a quasihomogeneous polynomial of degree $`\mathrm{d}\omega _1\omega _2`$ and $`\text{ord}\left(j_0^{\mathrm{}}(X_1)\right)\mathrm{d}\omega _1\omega _2\text{max}(\omega _1,\omega _2)`$. So, in the same way,
$$X_0.(fh)(divX_0)fh=\underset{i=1}{\overset{c}{}}\lambda _{i,1}u_iP_1f(X_1.(fh)(divX_1)(fh))\text{mod}B^2(f+fh).$$
Hence
$$g=\underset{i=1}{\overset{c}{}}(\lambda _{i,0}+\lambda _{i,1})u_i+(P_0+P_1)f(X_1.(fh)(divX_1)(fh))\text{mod}B^2(f+fh).$$
In this way, we get
$$g=\underset{i=1}{\overset{c}{}}(\lambda _{i,0}+\mathrm{}+\lambda _{i,k})u_i+(P_0+\mathrm{}+P_k)f(X_k.(fh)(divX_k)(fh))\text{mod}B^2(f+fh)$$
where $`k`$ is the smallest integer such that $`k(\mathrm{d}\omega _1\omega _2)\text{max}(\omega _1,\omega _2)0`$, $`P_j`$ is a quasihomogeneous polynomial of degree $`\mathrm{d}\omega _1\omega _2`$ for any $`j`$, $`\lambda _{i,j}๐`$ for any $`i`$ and $`j`$ and $`\text{ord}\left(j_0^{\mathrm{}}(X_k)\right)k(\mathrm{d}\omega _1\omega _2)\text{max}(\omega _1,\omega _2)`$.
But since
$`\text{ord}\left(j_0^{\mathrm{}}(X_k.(fh)(divX_k)fh)\right)`$ $``$ $`2\mathrm{d}\omega _1\omega _2+k(\mathrm{d}\omega _1\omega _2)\text{max}(\omega _1,\omega _2)`$
$``$ $`2\mathrm{d}\omega _1\omega _2,`$
lemma 4.10 gives $`X_k(divX_k)f=Qf\text{mod}B^2(f+fh)`$ for some quasihomogeneous polynomial $`Q`$ of degree $`\mathrm{d}\omega _1\omega _2`$.
\- This family is free. Let $`\lambda _1,\mathrm{}\lambda _r`$ be scalars and $`P`$ a quasihomogeneous polynomial of degree $`\mathrm{d}\omega _1\omega _2`$. Suppose that
$$\underset{i=1}{\overset{c}{}}\lambda _iu_i+Pf=X.(f+fh)(divX)(f+fh)()$$
with $`X๐ณ(๐^2)`$. Since $`fh`$ is quasihomogeneous of degree $`2\mathrm{d}\omega _1\omega _2`$ and the order of the $`\mathrm{}`$-jet of $`X`$ is larger than $`\text{max}(\omega _1,\omega _2)`$, we have
$`\text{ord}\left(j_0^{\mathrm{}}(X.(fh)(divX)(fh))\right)`$ $``$ $`2\mathrm{d}\omega _1\omega _2\text{max}(\omega _1,\omega _2)`$
$`>`$ $`2(\mathrm{d}\omega _1\omega _2)=\text{max}\{\text{deg}u_1,\mathrm{},\text{deg}u_c\}.`$
Therefore, if we write the $`\mathrm{}`$-jets in the relation $`()`$, we have $`_{i=1}^c\lambda _iu_iB^2(f)`$ and so $`\lambda _1=,...,=\lambda _c=0`$ (theorem 4.9).
We obtain
$$Pf=X.(f+fh)(divX)(f+fh)(),$$
where $`Pf`$ is a quasihomogeneous polynomial of degree $`2\mathrm{d}\omega _1\omega _2`$.
Now, we can write $`j_0^{\mathrm{}}(X)=_{i\delta }X^{(i)}`$ ($`X^{(i)}`$ is quasihomogeneous of degree $`i`$). If $`\delta <\mathrm{d}\omega _1\omega _2`$ then $`X^{(\delta )}Z^1(\mathrm{\Pi }_0)`$ and so, $`X^{(\delta )}=0`$ (cf proof of theorem 4.5).
In the same way, we can prove that $`X^{(i)}=0`$ for any $`i<\mathrm{d}\omega _1\omega _2`$.
Consequently, if we write the $`\mathrm{}`$-jets in the relation $`()`$, we get
$$Pf=X^{(\mathrm{d}\omega _1\omega _2)}.fdivX^{(\mathrm{d}\omega _1\omega _2)}f$$
that is, $`PfB^2(f)`$, which is possible only if $`P=0`$ (cf proof of theorem 4.9). $`\mathrm{}`$
### 4.3 Examples
We are going to make explicit the cohomology of some Poisson structures given in theorem 3.8.
The regular case : We suppose that $`\mathrm{\Pi }=x\frac{}{x}\frac{}{y}`$.
In this case, $`\omega _1=\omega _2=\mathrm{d}=1`$, and so $`\mathrm{d}\omega _1\omega _2<0`$. Moreover, $`Q_x=\{0\}`$.
Therefore, $`H^1(\mathrm{\Pi })๐\frac{}{y}`$ and $`H^2(\mathrm{\Pi })=\{0\}`$.
Morseโs singularity $`(A_1)`$: We suppose that $`\mathrm{\Pi }=(x^2+y^2)\frac{}{x}\frac{}{y}`$.
Here, we have $`\omega _1=\omega _2=1`$ and $`\mathrm{d}=2`$. The only monomials of degree $`\mathrm{d}\omega _1\omega _2`$ are the scalars. Moreover, $`Q_{x^2+y^2}๐.1`$. Then, we have
$`H^1(\mathrm{\Pi })`$ $``$ $`๐.(y{\displaystyle \frac{}{x}}x{\displaystyle \frac{}{y}})๐.(x{\displaystyle \frac{}{x}}+y{\displaystyle \frac{}{y}})`$
$`\text{and }H^2(\mathrm{\Pi })`$ $``$ $`๐.{\displaystyle \frac{}{x}}{\displaystyle \frac{}{y}}๐.f{\displaystyle \frac{}{x}}{\displaystyle \frac{}{y}}.`$
The singularity $`D_{2p+1}(p2)`$: We suppose that $`\mathrm{\Pi }=(x^2y+y^{2p})(1+x)\frac{}{x}\frac{}{y}`$.
In this case, we can see that $`\omega _1=2p1`$, $`\omega _2=2`$ and $`\mathrm{d}=4p`$. The only monomials of degree $`\mathrm{d}\omega _1\omega _2`$ are of type $`\lambda x`$ (with $`\lambda ๐`$). Moreover, the family $`\{1,x,y,y^2,\mathrm{},y^{2p}\}`$ is a monomial basis of $`Q_{x^2y+y^{2p}}`$.
Therefore,
the family $`\{\left[(1+x)\left((x^2+2py^{2p1})\frac{}{x}2xy\frac{}{y}\right)\right],\left[(1+x)xW\right]\}`$ is a basis of $`H^1(\mathrm{\Pi })`$
and the family
$`\{[x(x^2y+y^{2p})\frac{}{x}\frac{}{y}],[\frac{}{x}\frac{}{y}],[x\frac{}{x}\frac{}{y}],[y\frac{}{x}\frac{}{y}],[y^2\frac{}{x}\frac{}{y}],...,`$
$`...,[y^{2p}\frac{}{x}\frac{}{y}]\}`$ is a basis of $`H^2(\mathrm{\Pi })`$.
In particular, $`H^1(\mathrm{\Pi })`$ is of dimension $`2`$ and $`H^2(\mathrm{\Pi })`$ of dimension $`2p+3`$.
###### Remark 4.12
Our first approach to these problems was to use the spectral sequence associated to our complex, filtred by the valuation (whith respect to the quasihomogeneous degree). But the method we present here gives better results.
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# 1 Introduction. Fireball model
## 1 Introduction. Fireball model
Observed $`\gamma `$-bursts (GRB โ Gamma Ray Burst) is extremely interesting and still unexplained phenomenon (see reviews \[1-4\] and refs. therein).
Optical identification of the $`\gamma `$-bursts with โhostโ galaxies has proved that, at least, a part of them occurs in the galaxies with red shift of $`Z1`$, i.e. has the cosmological origin. This agrees well with a fully isotropic distribution of GRB over the sky and with statistic distribution of the burst events over their intensity.
Optical identification of GRB has allowed one to determine the distance to them and establish that a huge energy of $`10^{52}รท10^{54}`$ erg. in the $`\gamma `$-range (30-500 keV) is emitted in this phenomenon. Many of observed GRB characteristics are explained in the framework of the โfireballโ model \[5-7\], i.e. in terms of the electron-positron cloud, expanding with ultrarelativistic velocities. Ultrarelativistic velocities of the expansion (naturally appearing in the electron-positron plasma) , allow one to solve the problem of the GRB source compactness \[6-8\] and match the nonthermal GRB spectrum with short characteristic time of the GRB variability $`(\delta t10`$ms.)<sup>1</sup><sup>1</sup>1It is very important here that the bound on the number of baryons contained, which has to be small enough , is a condition of the ultrarelativistic expansion of lepton-photon plasma. Basing on the fireball model one has succeeded to explain both the observed effects of the long-lasting optical GRB afterglow, which appears as a result of the interaction between relativistic expanding fireball and inter-stars medium , and effect of the early afterglow, which intersects in time with GRB of high duration \[12-14\]. Thus, the model of the relativistic fireball containing a small number of baryons allows to agree the observed GRB characteristics and explain accompanying phenomena. However, some of questions remain unresolved yet:
1. The mechanism of the fireball production.
2. Large energy in the fireball.
3. The presence in some GRBโs of large number $`(10^2รท10^3)`$ of pulsation of the $`\gamma `$ emittion intensity with characteristic time of $`\delta t`$10 ms. We believe that this effect could serve as basic key to resolve the $`\gamma `$-bursts puzzle. In this paper we appeal the attention to the fact that the oscillations of the $`\gamma `$-quantum flux can naturally appear during the hydrodynamic collapse of some compact, massive, and nonrotating stars at final stage of their evolution.
## 2 The possibility of the oscillator burning of thermonuclear fuel during the process of hydrodinamic collapse
It is well-known that the production of sufficiently large iron core in the process of the star evolution is a reason of the hydrodinamic collapse of massive stars with masses $`M10M_{}`$. In this case the star core, having exhausted the source of thermonuclear energy, tends to compress and heat. The resulting increase of preasure, however, is unable to stop the compression, since the thermal energy is spent to the endothermal reaction of the iron nucleus decay and further - to the core neutronization. As a result, the core compression is transferred to the catastrophic hydrodinamic collapse, which is followed by the production of the hot neutron star. According to the idea by Fowler and Hoyle , the accretion of the nuclear fuel, which is left in the star envelope, onto the hot neutron star leads to its explosion and pollution manifesting itself as the burst of supernova. However, selfconsistent hydrodinamic calculations did not prove this assumption. It turned out that consequent account for the neutrino emittion leads to the delay of the collapse, which stops only when the core matter becomes nontransparent to the neutrino radiation. As a result of this delay, the burning of the accreting nuclear fuel occurs in a deep gravitational potential and, thus, the emitted thermonuclear energy is insufficient for the envelope ejection (see, for instance, review and refs. therein). The shock wave, occurring during the collapse delay, dumps only a small fraction of the envelope with the energy about $`10^{49}`$erg , that is two orders of magnitude less than a characteristic energy of the supernova explosion, $`10^{51}`$erg. Thus, this phenomenon was named as โsoundlessโ or silent collapse. Further attempts to explain the supernova bursts at the spherically symmetric collapse of massive stars had not led to a desirable result and now the most of specialists tend to the idea that the observed supernovas bursts with massive early-supernovas are anyway connected with effects of the collapsing star rotation: magnetic preasure onto the envelope, the Relay-Taylor instability, or breakage of the neutron star into two components (see, for instance, \[21-22\]).
However, one should take into account that the burning of the thermonuclear fuel during its accretion onto a hot neutron star can have the oscillating behavior. This effect is well-known and manifest itself in the consideration of the final evolution stage of stars with small masses, $`3M_{}M10M_{}`$. In such stars the oxigene-carbon core with degenerated electronic gas is produced as a result of the evolution. In this case the thermal explosion in the degenerated star core is the reason of the instability, when the core mass achieves the value close to the Chandler limit . In this case the oscillating character of thermonuclear fuel burning could be easily understood, if one takes into account a relatively small calorific power of this fuel with carbon-oxigene etc. content. The energy emitted during the thermal explosion leads to elimination of the degeneracy and increase of the thermal preasure resulting in the star expansion. As a result of the expansion the star temperature is decreased. This leads to the consequent compression of the star and enhancement of the thermonuclear burning which, in its turn, leads to the expansion, etc. All said above can be illustrated in Fig. 1, where one can see the results of calculations . Observed oscillation not only retained, but have been enhanced due to the account for convection, thus leading to delayed detonation with explosion energy $`10^{51}`$erg (report by V.S.Imshennik, seminar devoted to the memory of S.I.Syrovatsky, March 2, 2000).
It is possible that such oscillations appears as well in the layers of the thermonuclear fuel accreting onto a hot neutron star. In contrast to the case discussed above, they can have only local character, evolving in layers, adjacent to the surface of the hot neutron star. The period of these oscillations can be estimated by the arguments of dimension:
$$\tau \frac{1}{\sqrt{G_N\overline{\rho }}},$$
(1)
where $`G_N`$ is the gravitational constant, $`\overline{\rho }`$ is the matter density in the vicinity of the neutron star. According to the calculations (see Fig. 2), $`\overline{\rho }10^{11}`$g/cm<sup>3</sup>. Thus, the period of oscillations turns out to be equal to
$$\tau 10^2s,$$
(2)
that intriguesly coincides with oscillation period of the $`\gamma `$-quantum flux in some GRB<sup>2</sup><sup>2</sup>2It should be noted that the conditions for oscillation excitation are also realized in the hot neutron stars, produced as a result of the iron core collapse. However, their frequency is, at least, two orders of magnitude more. Density and temperature oscillations close to the surface of hot neutron star have to generate diverging shock waves in the surrounding envelope (with decreasing density vs. radius increase), repeating with the oscillation frequency.
## 3 The possibility of the electro-positron plasma stratification
One of the most important effects, which should be taken into account in description of shock waves passing through the star envelope, is a possible stratification of the electron-positron plasma, which happens without violation of the electroneutrality of ordinary matter (nuclei and electrons, which compensate their electric charge). Such stratification is possible under the condition of light preasure in propagating shock wave, as the Eddington limit for electron-positron plasma is 3600 times lower than that for an ordinary matter with nuclei $`A2Z`$. So, the electron-positron plasma, when exiting to the star surface, will contain a relatively low concentration of baryons (that is just needed to agree the fireball model and observed data). It should also be noted that under the condition of the rarefacted star atmosphere the equilibrium electro-positron plasma can appear at relatively low temperatures, since under these conditions at $`kT<<m_ec^2`$:
$`n_{e^+}n_e^{}`$ $``$ $`{\displaystyle \frac{1}{(2\pi ^3)^{1/2}}}\left({\displaystyle \frac{m_ec}{\mathrm{}}}\right)^3e^{1/x}x^{3/2}`$
$`x`$ $`=`$ $`{\displaystyle \frac{kT}{mc^2}}<<1`$
while in the dense matter the positron concentration will be proportional to $`\mathrm{exp}\left(\frac{2mc^2}{kT}\right)`$ (see, for instance, ). The fact that sufficiently high temperatures $`(kTm_ec^2)`$ are achieved close to the center of the collapsing star is proved by the presence of the process of explosive nucleosynthesis of the $`{}_{}{}^{56}Ni`$ nuclei with subsequent production of $`{}_{}{}^{56}Co`$, as it comes from the observation data on SN$`1987`$ . (According to the calculations , in the region of neutrinosphere $`kT5,6`$ MeV).
When leaving the star atmosphere the expanded cloud of electron-positron plasma inevitably gains the ultrarelativistic character (see ). This, as it is well-known, leads to variations of the $`\gamma `$-radiation momentum, received by a remote observer.
$$\delta t\frac{R}{2c\mathrm{\Gamma }^2},$$
(4)
where $`R`$ โ the cloud radius, $`\mathrm{\Gamma }=(1v^2/c^2)^{1/2}`$ โ the Lorentz factor, corresponding to expansion velocity, $`v`$. It is evident that the oscillation of the $`\gamma `$-quantum flux will be observed if
$$\delta t\tau .$$
(5)
In opposite case $`(\delta t>>\tau )`$, the oscillations in the $`\gamma `$-quantum flux for a remote observer are smeared. This helps to explain the fact that there are no observed oscillations in some GRB. Stratification of the electron-positron plasma from ordinary matter should lead to situation, when each oscillation in the vicinity of the neutron star will produce shock waves in the form of two shells expanding with different velocities. Here, it could happen that the electron-positron shell, emitted in the following oscillation can overtake that one containing baryons and emitted in previous oscillation. Thus, there can be an interaction of shock waves inside the fireball itself .
Ultrarelativistic character of the fireball expansion $`(\mathrm{\Gamma }10^2)`$ allows to conclude that at observed values of $`\delta t10^2`$s. the fireball size can be large enough (and the plasma density is low enough, correspondingly), to consider the fireball as a โthinโ source. It allows to explain the nonthermal (power) spectrum of GRB. โIntrinsicโ shock waves, producing the fireball, can also generate high energy particles by means of well-known mechanism of acceleration. This can explains the observation in some GRB the high energy $`\gamma `$-quanta (up to 18 GeV).
## 4 Possible progenitors of GRB
The scenario of the oscillation origin in GRB developed above confines the class of object, which could be the progenitors of GRB.
First, these should be sufficiently massive stars with masses $`M>(1520)M_{}`$. Quite stage of such star evolution has to be ended within a time period about $`10^6`$ years or less. Star mass, large enough, is also required to explain the GRB energy.
Second, these should be nonrotating (or with low angular velocity) stars. It seems that stars with high angular velocity should explode due to the effects connected with their rotation, as ordinary supernovas of the II-type with ejection of relatively massive envelope .
And third, these should be compact stars devoided of extent hydrogen ad, probably, in part helium envelope, which is able to prevent the outer ejection of the electron-positron plasma due to the processes of positron annihilation.
The stars of the Wolf-Rayet type (WR) meet all these requirements โ they are the most massive compact start, which have lost almost all their hydrogen and, in part, helium envelope during their evolution. It is possible that, namely, due to the loss of the balk of their envelope these stars have lost their rotatory impulse. Anyway, the rotation is observed only for 15% of the WR-stars . In the studies by A.M. Cherepashchuk et al. (see and refs. therein) it was found that one can neglect the decrease of the WR-star masses in the process of their further evolution (which is caused by the stellar wind). It allows one to compare the masses of the WR-stars and their -nuclei with masses of the relativistic objects (neutron stars and โblack holesโ), for which the WR-stars are the progenitors. Basing on the masses measurement of the X-ray sources in double systems A.M. Cherepashchuk has drawn the extremely important conclusion that the distribution of the X-ray sources masses has clear bimodal character. There is a mass gap between the neutron stars โ pulsars, whose masses are ranged in the narrow band $`(12)M_{}`$ with average mass $`(1.35\pm 0.15)M_{}`$, and masses of candidates to black holes, which are distributed in the range of $`(5รท15)M_{}`$ and have average mass of $`(8รท10)M_{}`$. The bimodal mass distribution and the presence of the gap serves as the indication to different origin of these objects. As for the massive candidates to black holes, the correlation (discovered by A.M. Cherepashchuk ) between their masses and masses of the WR stars, which are in the range of $`(555)M_{}`$ and for which the average value of their CO-nuclei is $`(8รท12)M_{}`$, close to the average value of masses of the observed candidates to black holes, seems to be very important to understand their origin. Thus, there are arguments to assume that, at least, some of the WR stars collapse into massive objects through the โsoundlessโ collapse without any significant ejection of their envelope. In the first turn it concerns more evolved WC stars with reach content of C-nuclei in their envelope (produced due to the thermonuclear burning of helium) and with average mass of $`13,4M_{}`$. The data presented in are the strong argument in the favour of the assumption formulated by P.Conti in 1982 y. that the WR-stars more often disappears in the form of the โwhimperโ, rather than explosion <sup>3</sup><sup>3</sup>3It is possible that some WR-stars are early supernovas 1b. (I would like to acknowledge this remark by V.S. Imshennik).
The compact structure of the WR-stars allows one to assume that the electron-positron plasma shells, which appears due to the stratification in shock waves, can leave the star surface and even a small part of the large gravitational energy emitted in the massive star collapse can explain the GRB energy.
## 5 The GRB energy
To determine the GRB energy it is necessary to have self-consistent hydrodinamic calculations of the process of soundless massive star collapse with the account for the possibility of the $`e^+e^{}`$-plasma stratification from the rest of the matter. However, one can try to estimate a possible GRB energy from physical (though, not very reliable) arguments. If as the result of the collapse into the hot neutron star with mass $`M=1,5M_{}`$ the gravitational energy $`ฯต510^{53}`$erg. is emitted in the form of the neutrino radiation, then for masses $`M=(15รท20)M_{}`$ (in the case without the envelope ejection) one can expect the energy emittion about $`ฯต10^{56}`$erg.<sup>4</sup><sup>4</sup>4It should noted that a hot neutron star should be stable up to the mass of $`M_{NS}70M_{}`$ .. So, to provide the GRB energy $`10^{53}`$erg. it is sufficient to have $`(e^+e^{})`$-plasma ejection accumulated $``$0,1% of emitting gravitational energy. Analogous estimate one can obtain using (taking some risk) the calculation results of the hydrodinamic collapse of the iron-oxygene star core . Though the shock wave generated in this process is subjected to the attenuation due to neutrino radiation and its power is not sufficient to explain the supernova explosion, nevertheless, the energy of the emitted shell can be $`10^{49}`$erg. For the obtained velocity of the shell expansion of $`v1,510^3`$km/s the mass of the ejected shell is $`\mathrm{\Delta }M0,44M_{}`$. Were such mass being ejected in the form of the $`e^+e^{}`$-plasma, the energy of $`810^{53}`$erg. should be emitted during its subsequent annihilation. The comparison given here is not proved well, but it gives an idea on the possible effect value. Thus, the bulk of the GRB energy in the mechanism considered has the gravitational origin. The heating of the collapsing star leads to the production of the dense and hot $`e^+e^{}`$-plasma, and energy emitted in oscillatoric burning of the thermonuclear fuel is spent to generation of shock waves pushing the $`e^+e^{}`$-plasma beyond the star.
One should also take into account the possible additive energy to the expanding $`e^+e^{}`$-plasma due to neutrinos and antineutrinos radiated during the collapse process (since, during their scattering on the electrons and positrons they can pass to laters the bulk of their energy). The question of $`\gamma `$-radiation spectra from GRB requires special treatment. It is quite possible that the absence of the 511keV line from the $`e^+e_{}`$-annihilation at the rest is connected with the ultrarelativistic fireball expansion.
## Discussions
On the basis of the observation data there are forcible arguments to assume that massive, compact, and nonrotating stars of the Wolf-Rayet type are subjected to the relativistic collapse without any significant ejection of their envelope. This assumption agrees well with the fact that within the hydrodinamic calculations one fails to obtain the envelope ejection sufficient to explain the supernova bursts. The gravitational energy emitted during the relativistic collapse of such objects can be about $`10^{55}รท10^{56}`$erg.
The hypothesis developed in this paper is that the burning of the thermonuclear fuel accreting onto a hot neutron star can proceed in the oscillatoric regime, which generate the shock waves, which, in their turn, push out the $`e^+e^{}`$-plasma outside the star surface (in the presence of its stratification). This hypothesis qualitatively explains the origin of the relativistic fireball with a low baryon content and oscillations observed in GRB (moreover, the oscillation period is explained quantitatively in the order of magnitude). The duration of GRB ($`20`$s) agrees with the time of the outer envelope accretion onto the neutron star and time of its cooling. A number of observed data, provided by B. Paczynski, in particular the indication that GRB happen in the regions of intensive star production , tell in the favor of the fact that the WR-stars can be the progenitors of GRB.
According to the hypothesis suggested in this paper, the collapse of the WR-star and ejection of the $`e^+e^{}`$-plasma happen spherical symmetrically. Successful description of the optical afterglow at time-period $`200`$ days , obtained in the framework of this hypothesis, proves in favour of the spherical symmetry of a number of the $`\gamma `$-bursts. (However, there are some indications to the fact in that jets can appear in some of GRB.)
In the conclusion the author would like to thank G.V. Domogazky, A.M. Dyhne, A.A. Logunov, V.S. Imshennik, D.K. Nadezhin, K.A. Postnov for the stimulating interest to this work and valuable remarks. Special thanks to A.M. Cherepashchuk for letting me know his work and data on the WR stars.
This work is supported, in part, by the RFBR grants 99-02-16558 and 00-15-96645.
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# The morphology of the Magellanic Clouds revealed by stars of different age: results from the DENIS survey
## 1 Introduction
The Magellanic Clouds are our closest neighbors allowing direct observation of individual constituent objects. They are bound to the Galaxy and show signs of strong interaction with the Milky Way about $`0.2`$ Gyr ago (Westerlund west (1997)). The LMC is classified as an irregular dwarf galaxy, its most prominent feature is a central bar, much like those found in barred spiral galaxies. Its eastern side is closer than its western side (Caldwell & Coulson cc (1986)). Underlying the bar is a circular disk of older stars (Westerlund west (1997)). The appearance of the SMC is characterized by a much less pronounced bar, and an eastern extension called the Wing. Lines-of-sight through the SMC appear to cover extensive depths; the Wing and the northeastern part of the Bar are closer than the southern parts (Westerlund west (1997)).
Newly obtained large photometric data sets at different wavelengths and with improved sensitivity and spatial coverage allow us to investigate the large scale properties of the Magellanic Clouds. In particular, data in the near infrared allow us to access stages of stellar evolution that are marginally covered by optical data, such as the RGB and AGB phases.
Very recently Zaritsky et al. (zar (2000)) found that the asymmetric appearance of the SMC is primarily caused by the distribution of young stars, and that the older stars have a very regular distribution. It is not possible from their Figures to evaluate the behavior of the density towards the center of the Cloud. Weinberg & Nikolaev (wn (2000)) point out the presence of intervening tidal debris up to $`15`$ kpc from the LMC.
## 2 The Data
Our morphological study of the Magellanic Cloud is based on a sample of stars extracted from the DCMC catalogue (Cioni et al. cio (2000)). The sample includes all sources detected in both $`I`$ and $`J`$, irrespective of detection in $`K_S`$. The DCMC catalogue contains sources detected in at least two of the three DENIS bands ($`I`$: $`0.8\mu `$m, $`J`$: $`1.25\mu `$m and $`K_S`$: $`2.15\mu `$m) within a surface area of $`20\times 16`$ square degrees centered on $`(\alpha ,\delta )=(5^h27^m20^s`$,$`69\mathrm{ยฐ}00\mathrm{}00\mathrm{})`$ toward the LMC and $`15\times 10`$ square degrees centered on $`(\alpha ,\delta )=(1^h02^m40^s`$,$`73\mathrm{ยฐ}00\mathrm{}00\mathrm{})`$ toward the SMC; J2000 coordinates are used throughout this paper. The observations have been performed with the DENIS instrument (Epchtein et al. epal (1997)) on the 1mโESO telescope.
We have used the ($`IJ`$, $`I`$) colourโmagnitude diagrams (Figure 1 for the LMC only) to select three classes of objects in each Cloud. Sources labelled (A) with $`I<4.64\times (IJ)+19.78`$ represent the youngest population in the Magellanic Clouds: the brightest dwarf stars, blueโloop stars and supergiants (third vertical sequence from the left), together with an unrelated foreground component of dwarfs and giants (first two vertical sequences from the left). Sources labelled (B) with $`I>4.64\times (IJ)+19.78`$, located above the tip of the red giant branch (TRGB โ Cioni et al. 2000a ) are mainly asymptotic giant branch stars (AGB). Sources labelled (C) with $`I>4.64\times (IJ)+19.78`$ located below the TRGB are mostly red giant branch stars (RGB) and represent the oldest population in the Clouds. For the sake of clarity, we have plotted in Fig. 1 only those sources that were detected in all three wave bands and that occur in the very central part of the Cloud. Sources detected only in $`I`$ and $`J`$ predominantly populate the lower part of the diagram. The position of the TRGB is indicated by a horizontal line. The $`I`$, $`J`$ and $`K_S`$ sensitivity limits are $`18`$, $`16`$ and $`14`$ mag respectively. Photometric errors widen the sequences towards fainter magnitudes.
For each class of objects in each of the two Clouds, we show their distribution in the plane of the sky by counting the sources in bins of $`0.2\mathrm{ยฐ}\times 0.2\mathrm{ยฐ}`$, applying a light smoothing to the resulting structure (Figs. 27); the contour values increase logarithmically. Regions corresponding to missing data (strips at constant RA indicated by diamonds) were filled in by interpolation. Their effect is mostly negligible except in Fig. 4 where strips of possibly lower photometric quality may be causing discontinuities in the outermost contours.
## 3 Spatial distribution
The contribution due to Galactic foreground stars has not been subtracted from the maps. Its influence is most clearly seen in Fig. 2 in the direction of the Galactic Plane. In the other maps, the foreground contribution is rather constant and does not affect the morphology of the Clouds.
### 3.1 Structure of the LMC
The lower contours in Figs. 2โ4 show an almost circular outline (axial ratios consistent with an inclination $`i=30\mathrm{ยฐ}40\mathrm{ยฐ}`$) centered in all three cases near $`\alpha =5^h20^m,\delta =69\mathrm{ยฐ}`$ with major axis at about $`13^{}`$. Westerlund (1997) gives a similar diameter for the stars of the old disk. This stellar disk also coincides in shape and extent with the HI disk (Kim et al. 1998). The center of the disk is offset from the center of the Bar by about 30$``$ to the north (see Fig. 2). We confirm the conclusion by Westerlund (1997) that the LMC consists of two systems: a circular disk and an off center bar. Half of the total number of stars are in the bar and this factor (Fig. 2) increases for younger objects. Unless this is a transient configuration, it thus seems that the LMC must be embedded in a gravitational potential produced by an unseen mass component (see also Sofue 1999). This is in agreement with the conclusion by Stil (1999) that the class of dwarf galaxies to which the LMC belongs (โfast rotatorsโ) is dominated by dark matter.
The youngest component (younger than $`0.5`$ Gyr) is composed of very bright main-sequence dwarf stars, blueโloop stars and supergiants. Their distribution (Fig. 2) is clumpy and irregular. The Bar, extending over about $`4\mathrm{ยฐ}`$, is prominent and contains a well defined nuclear concentration at its center. The region of 30 Dor is represented by the small feature just above the northeastern side of the Bar, and the Shapley Constellation III is the large structure at $`\delta 67\mathrm{ยฐ}`$. Elongations at either end of the Bar indicate the presence of spiral arms most clearly seen in the northwest at the location of the giant HII region complex N 11. Similar structures are seen in the distribution of stellar complexes (Maragoudaki et al. mara (1998)), associations and HII regions (Bica et al. bical (1999)). Clusters (Bica et al. bisch (1995), Kontizas et al. kon (1990)) have a distribution more similar to the one of AGB/RGB stars.
The distribution of AGB stars (Fig. 3), also relatively young (around $`1`$ Gyr) likewise reveals a prominent Bar and nucleus. Shapley Constellation III is inconspicuous in AGB stars. A broad and faint spiral arm begins at the northwestern end of the Bar and bifurcates around $`\alpha =05^h10^m,\delta =66.5\mathrm{ยฐ}`$. The spiral arm feature originating at the southeastern end of the Bar is clearly delineated in the AGB star population and can easily be followed to $`\alpha =4^h45^m,\delta =73\mathrm{ยฐ}`$. It was noted before by Bothun & Thompson (bt (1988)) in their surface photometry study of the Magellanic Clouds โ see their $`1.1<BR<1.35`$ diagram. At least this spiral arm might be due to tidal action, as it appears to be connected to the Magellanic Cloud Bridge (cf. Staveley-Smith et al. ss (2000)). The outernmost contour well matches the carbon stars by Kunkel et al. (kunk (1997)).
The oldest population (from $`1`$ to $`5`$ Gyr), represented by RGB stars (Fig. 4), once again reveals a prominent Bar which is significantly broader than that defined by the younger populations. Galactic foreground stars may affect the outermost contours. The southern spiral arm is inconspicuous, but the two faint northern spiral arms seen in Fig. 3 (AGB) have weak counterparts in the form of extensions at $`\alpha =5^h6^h,\delta =64\mathrm{ยฐ}68\mathrm{ยฐ}`$.
Bothun & Thompson ( bt (1988)) conclude that the LMC has a relatively large scale length more appropriate for galaxies with obvious spiral structure than for other dwarf galaxies. It is interesting that the asymmetric spiral structure delineated by the different components in Figs. 24 is in fairly good agreement with the HI map shown by Gardiner et al. (gard (1998)) and is nicely reproduced by their dynamical model.
### 3.2 Structure of the SMC
The structure of the SMC is still not understood (Westerlund 1997). Our maps show that populations of different age have different distributions. The youngest component has an asymmetric distribution(Fig. 5) elongated along a NEโSW axis (PA $`45\mathrm{ยฐ}`$). In the south, the outermost contour defines four protuberances which might be associated with tidal features: at least the eastern (coincident with the SMC Wing) and western protuberances are aligned with that of the Magellanic Cloud Bridge (cf. Staveley-Smith et al. (ss (2000)). Higher contours show an extension in the northeast, aligned with the main body of the SMC Bar. The Bar structure itself is similar to that seen in the distribution of young clusters (Bica & Dutra bd (2000)) and in the upper main-sequence map by Zaritsky et al. (zar (2000)). Clusters, associations and HII regions (Bica & Schmitt bisch (1995)) are also found at the locations of the southern protuberances. The young stars are strongly concentrated in the southwestern part of the SMC Bar. Outside the main body of the SMC, the two Galactic globular clusters NGC 104 = 47 Tuc (west) and NGC 362 (north) can be discerned. The HI column density contours in the map presented by Stanimirovicฬ et al. (stan (1998)) outline the distribution of the young stars quite well.
The AGB stars have a more regular distribution (Fig. 6) with two prominent central concentrations matching the carbon stars by Hardy et al. (hardy (1989)). The easternmost also coincides with the peak of the youngโstar distribution. The AGB distribution axis is much less inclined (PA $`75\mathrm{ยฐ}`$) than that of the younger and very similar to that of the RGB star distribution. As in the case of the LMC, the stellar distributions become more regular and smoother with increasing age, also apparent in the $`B`$ and $`V`$ band images by Zaritsky et al. (zar (2000)) and for the outer contour in the carbon stars by Kunkel et al. (kdi (2000)); carbon stars by Rebeirot et al. (reb (1983)) fill the second level countour (Fig. 6).
The distribution of RGB stars (Fig. 7) is similar to that of the AGB stars and also exhibits two major concentrations. The western most is more pronounced in RGB than in AGB stars. The eastern concentration on average appears to be significantly younger than the western concentration dominated by the older stars. Remarkably, the strongest HI concentration in the SMC map by Stanimirovicฬ et al. (stan (1998)) appears to be just between the concentration of younger stars and that of older stars. It is also remarkable that the older star distribution extends over the full length of the suspected southwestern tidal feature, about $`2\mathrm{ยฐ}`$ from the main body of the Bar. With respect to the overall distribution of the older stars, that of the HI appears to be displaced towards the east. The SMC Wing, prominent in HI and also traceable in the younger stellar population, has no counterpart in the older stars.
## 4 Conclusions
Counts of sources towards the Magellanic Clouds extracted from the DCMC (Cioni et al. cio (2000)) allow differentiation in the ($`IJ`$, $`I`$) colourโmagnitude diagram into three groups of objects with different mean ages. The spatial distribution of the three age groups is quite different: in either Cloud, the youngest stars exhibit an irregular structure characterized by spiral arms and tidal features while the older stars are smoothly and regularly distributed. The distribution of younger stars is wellโcorrelated with those of clusters, associations, HII regions and HI. The significant offset of the LMC Bar with respect to the overall circular disk suggests that the LMC potential is dominated by dark matter. The wellโdefined southern spiral arm may be due to tidal interaction with the SMC. The nature of the two northern spiral arms is uncertain. In the SMC, the regular, but doubleโpeaked, structure of the AGB and RGB stars is remarkable, as is its offset from the HI distribution, and the mean age difference of the two maxima. Relatively faint eastโwest features in the younger star population (including the Wing) are probably also due to tidal interaction.
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# Perturbative method for generalized spectral decompositions.
## I Introduction.
Following previous ideas about the use of analytic continuation in the study of unstable quantum systems , , , , in papers , and we have developed our own version, obtaining new results. But, as there were not a perturbative method adapted to our formalism, we were forced to study only solvable simple models. Therefore, this paper is devoted to fill this gap introducing a perturbation method adapted to our formalism. Let us briefly review the problem and the history of the proposed solution.
The decay of unstable systems in quantum mechanics is usually described by a Hamiltonian of the form $`H=H_{free}+H_{int}`$, where $`H_{free}`$ has a discrete spectrum imbedded in the continuous spectrum, and $`H`$ has a continuous spectrum. In non equilibrium quantum statistical mechanics, the generator of the time evolution is the Liouville-Von Newmann operator which in some cases can be written as $`๐=L_{free}+๐_{int}`$. There is an overlapping between the continuous and the discrete part of the spectrum of $`๐_{free}`$, where in this case the discrete part of the spectrum contains the zero eigenvalue of the generator of the free evolution. It is expected that the infinite dimensional degeneration of the zero eigenvalue to be partially removed by the interaction, so that the zero eigenspace of $`๐`$ turns out to be the equilibrium subspace.
E.g., for the Friedrichs model, a prototype of a decay system in quantum mechanics, with Hamiltonian
$`H=\mathrm{\Omega }|11|+{\displaystyle _0^{\mathrm{}}}๐\omega \omega |\omega \omega |+{\displaystyle _0^{\mathrm{}}}๐\omega V_\omega \left\{|\omega 1|+|1\omega |\right\},\mathrm{\Omega }^+,V_\omega =\overline{V_\omega },`$
it is possible to find conditions on the interaction $`V_\omega `$ for which the spectrum of $`H`$ is $`^+`$. The discrete part of the spectrum of $`H_0=\mathrm{\Omega }|11|+_0^{\mathrm{}}๐\omega \omega |\omega \omega |`$ is eliminated by the interaction. A complete set of generalized right (left) eigenvectors $`|\omega ^+`$ ($`\omega ^+|`$) of the Hamiltonian $`H`$ with eigenvalue $`\omega ^+`$ can be explicitly obtained, thus $`H=_0^{\mathrm{}}๐\omega \omega |\omega ^+\omega ^+|`$. When $`V_\omega 0`$, $`|\omega ^+`$ tends to $`|\omega `$ and therefore the discrete eigenvector $`|1`$ is not recovered.
E.C.G.Sudarshan et al obtained a generalized spectral decomposition for the Friedrichs model by considering a Hamiltonian obtained from the previous expression by replacing the integrals over $`^+`$ by integrals over a curve $`\mathrm{\Gamma }`$ in the lower complex half plane. This extended Hamiltonian $`H_\mathrm{\Gamma }`$ is meaningful as the generator of the time evolution for state vectors $`\mathrm{\Phi }`$ for which $`\omega |\mathrm{\Phi }`$ has a well defined analytic extension to the lower complex half plane. The spectrum of $`H_\mathrm{\Gamma }`$ is the curve $`\mathrm{\Gamma }`$ and a point $`z_0`$ located in the lower half plane between the curve $`\mathrm{\Gamma }`$ and the real axis. When $`V_\omega 0`$, $`z_0`$ tends to $`\mathrm{\Omega }`$, and the discrete part of the spectrum of the โfreeโ Hamiltonian is recovered (see also reference ).
The same complex spectral decomposition for the Friedrichs model was also obtained by T.Petrosky, I.Prigogine and S.Tasaki , with a perturbative method using a โtime ordering ruleโ by which an small imaginary part is included to avoid small denominators, with a different sign according to the type of transition involved.
Latter, I.Antoniou et.al. , developed a formalism to deal with singular observables and generalized states, which they applied in the framework of the subdynamics theory. We gave our version of this formalism in papers and . We also considered singular observables and generalized states to describe the time evolution of a quantum oscillator interacting with a field, describing the evolution of the system to equilibrium in the thermodynamic limit .
The analytic dilation method is also a well established procedure to determine the contributions of the singularities to the time evolution of unstable quantum systems.
In this paper we consider states and observables with suitable analytic properties (as in papers and ), such that the integrals over the continuous part of the spectrum can be deformed into integrals over curves in the complex plane. The physical motivation of these analytic properties is explained in . In this way we can avoid the ill defined terms involving the differences between discrete and continuous eigenvalues of the unperturbed generator of the time evolution, and implement a well defined algorithm to obtain a complete biorthogonal set of generalized eigenvectors in one to one correspondence with the basis of the unperturbed problem. This biorthogonal system can be ued to give a description of the decay process.
The paper is organized as follows:
In section II we present the perturbative algorithm for the case of pure states, compute the complete set of right and left generalized eigenvectors and discuss the roll of the complex eigenvalues in the approximated expressions for the decay process. The method is applied to study the tunneling through a barrier in a one dimensional problem and to the Friedrichs model, were the results are compared with the exact solutions of reference . In section III the case of generalized states and singular observables for the Friedrichs model is studied with the perturbative method. In the conclusions of section IV we summarize the results and discuss the physical interpretation of the complex generalized spectral decomposition.
## II Pure States.
Let us consider a system where the free Hamiltonian $`\stackrel{0}{H}`$ has a continuous spectrum plus a single discrete eigenvalue, and there is a general interaction $`\stackrel{1}{H}`$:
$`H`$ $`=`$ $`\stackrel{0}{H}+\stackrel{1}{H}`$ (1)
$`\stackrel{0}{H}`$ $``$ $`\mathrm{\Omega }|11|+{\displaystyle _0^{\mathrm{}}}๐\omega \omega |\omega \omega |\mathrm{\Omega }^+,`$ (2)
$`\stackrel{1}{H}`$ $``$ $`{\displaystyle _0^{\mathrm{}}}๐\omega V_\omega |\omega 1|+{\displaystyle _0^{\mathrm{}}}๐\omega \overline{V_\omega }|1\omega |+{\displaystyle _0^{\mathrm{}}}๐\omega {\displaystyle _0^{\mathrm{}}}๐\omega ^{}V_{\omega \omega ^{}}|\omega \omega ^{}|,V_{\omega \omega ^{}}=\overline{V_{\omega ^{}\omega }},\mathrm{\Omega }^+,`$ (3)
We suppose that the functions $`V_\omega `$ and $`V_{\omega \omega ^{}}`$ are endowed with some analytic properties that we will specify below equation (8). The generalized right (left) eigenvectors $`|1`$ and $`|\omega `$ ($`1|`$ and $`\omega |`$) of $`\stackrel{0}{H}`$ form a complete biorthogonal set that we can use to expand any vector state $`|\mathrm{\Phi }`$, i.e.
$$1|1=1,\omega |\omega ^{}=\delta (\omega \omega ^{}),1|\omega =\omega |1=0,|\mathrm{\Phi }=|11|\mathrm{\Phi }+_0^{\mathrm{}}๐\omega |\omega \omega |\mathrm{\Phi }.$$
(4)
The amplitude of a state vector $`\mathrm{\Phi },`$ observed in the state vector $`\mathrm{\Psi }`$, is given by
$$\mathrm{\Psi }|\mathrm{\Phi }=\mathrm{\Psi }|11|\mathrm{\Phi }+_0^{\mathrm{}}๐\omega \mathrm{\Psi }|\omega \omega |\mathrm{\Phi },$$
(5)
and therefore
$`\mathrm{\Psi }|\mathrm{\Phi }=\mathrm{\Psi }|I|\mathrm{\Phi },I|11|+{\displaystyle _0^{\mathrm{}}}๐\omega |\omega \omega |.`$
If we try to construct a complete biorthogonal set of eigenvectors of the Hamiltonian $`H=\stackrel{0}{H}+\stackrel{1}{H}`$, depending analytically on the small interaction parameter which we assume included in $`\stackrel{1}{H}`$, the standard methods of perturbation theory are not applicable, because the superposition of the discrete and continuous part of the spectrum of $`\stackrel{0}{H}`$ will produce non defined terms in the perturbation expansion. However, this problem can be avoided if we are satisfied to obtain information about terms of the form $`\mathrm{\Psi }|\mathrm{\Phi }`$, $`\mathrm{\Psi }|H|\mathrm{\Phi }`$ or $`\mathrm{\Psi }|\mathrm{exp}(iHt)|\mathrm{\Phi }`$, where $`\mathrm{\Phi }`$ is a vector for which $`\phi (\omega )=\omega |\mathrm{\Phi }`$ has a well defined analytic extension $`\phi (z)`$ to the lower half plane, and $`\mathrm{\Psi }`$ is a vector for which $`\psi (\omega )=\omega |\mathrm{\Psi }`$ has a well defined analytic extension $`\psi (z)`$ to the upper half plane. In this case the functionals $`z|`$ and $`|z`$ are defined through the analytic expressions:
$`z|\mathrm{\Phi }\phi (z),\mathrm{\Psi }|z\overline{\psi (\overline{z})}.`$
Then, the Cauchy theorem can be used in equation (5) to write
$$\mathrm{\Psi }|\mathrm{\Phi }=\mathrm{\Psi }|11|\mathrm{\Phi }+_0^{\mathrm{}}๐\omega \mathrm{\Psi }|\omega \omega |\mathrm{\Phi }=\mathrm{\Psi }|11|\mathrm{\Phi }+_\mathrm{\Gamma }๐z\mathrm{\Psi }|zz|\mathrm{\Phi },$$
(6)
where $`\mathrm{\Gamma }`$ is the curve in the lower complex half plane indicated in Fig. 1. Therefore, for the states $`\mathrm{\Phi }`$ and observables $`\mathrm{\Psi }`$ endowed with the analyticity properties mentioned above, we can define the new โidentity operatorโ $`I_\mathrm{\Gamma }`$
$`I_\mathrm{\Gamma }|11|+{\displaystyle _\mathrm{\Gamma }}๐z|zz|,\mathrm{\Psi }|\mathrm{\Phi }=\mathrm{\Psi }|I_\mathrm{\Gamma }|\mathrm{\Phi }.`$
The functionals $`|1`$, $`1|`$, $`|z`$ and $`z|`$ satisfy
$$1|1=1,z|z^{}=\delta _\mathrm{\Gamma }(zz^{}),1|z=z|1=0,$$
(7)
where $`\delta _\mathrm{\Gamma }(zz^{})`$ is the $`\delta `$ distribution defined on the curve $`\mathrm{\Gamma }`$ ($`_\mathrm{\Gamma }๐z^{}\delta _\mathrm{\Gamma }(zz^{})f(z^{})=f(z)`$).
We can also define $`H_\mathrm{\Gamma }=\underset{\mathrm{\Gamma }}{\overset{0}{H}}+\underset{\mathrm{\Gamma }}{\overset{1}{H}}`$, where
$$\underset{\mathrm{\Gamma }}{\overset{0}{H}}\mathrm{\Omega }|11|+_\mathrm{\Gamma }dzz|zz|,\underset{\mathrm{\Gamma }}{\overset{1}{H}}_\mathrm{\Gamma }dzV_z|z1|+_\mathrm{\Gamma }dz\overline{V_{\overline{z}}}|1z|+_\mathrm{\Gamma }dz_\mathrm{\Gamma }dz^{}V_{zz^{}}|zz^{}|,$$
(8)
where $`V_z`$, $`\overline{V_{\overline{z}}}`$ and $`V_{zz^{}}`$ are the analytic extensions to the lower half plane of $`V_\omega `$, $`\overline{V_\omega }`$ and $`V_{\omega \omega ^{}}`$. Therefore we postulate that these functions can be analytically extended at least up to the curve $`\mathrm{\Gamma }`$. The operators $`\underset{\mathrm{\Gamma }}{\overset{0}{H}}`$ and $`\underset{\mathrm{\Gamma }}{\overset{1}{H}}`$ verify $`\mathrm{\Psi }|\stackrel{0}{H}|\mathrm{\Phi }=\mathrm{\Psi }|\underset{\mathrm{\Gamma }}{\overset{0}{H}}|\mathrm{\Phi }`$ and $`\mathrm{\Psi }|\stackrel{1}{H}|\mathrm{\Phi }=\mathrm{\Psi }|\underset{\mathrm{\Gamma }}{\overset{1}{H}}|\mathrm{\Phi }`$ for the restricted class of vectors defined above, and therefore $`H_\mathrm{\Gamma }`$can be considered as the generator of the time evolution. Being $`\mathrm{\Omega }`$ a real number, there is no superposition between the discrete and the continuous part of the spectrum of $`\underset{\mathrm{\Gamma }}{\overset{0}{H}}`$.
We can now proceed with a perturbation approach to the right eigenvalue problem
$`H_\mathrm{\Gamma }|\mathrm{\Phi }=\lambda |\mathrm{\Phi }.`$
By assuming the expansion
$`|\mathrm{\Phi }=|\stackrel{0}{\mathrm{\Phi }}+|\stackrel{1}{\mathrm{\Phi }}+|\stackrel{2}{\mathrm{\Phi }}+\mathrm{},\lambda =\stackrel{0}{\lambda }+\stackrel{1}{\lambda }+\stackrel{2}{\lambda }+\mathrm{},`$
with respect to the small interaction parameter that we suppose contained in $`\underset{\mathrm{\Gamma }}{\overset{1}{H}}`$, we obtain order by order the set of equations
$`(\stackrel{0}{\lambda }\underset{\mathrm{\Gamma }}{\overset{0}{H}})|\stackrel{0}{\mathrm{\Phi }}`$ $`=`$ $`0`$ (9)
$`(\stackrel{0}{\lambda }\underset{\mathrm{\Gamma }}{\overset{0}{H}})|\stackrel{1}{\mathrm{\Phi }}`$ $`=`$ $`(\underset{\mathrm{\Gamma }}{\overset{1}{H}}\stackrel{1}{\lambda })|\stackrel{0}{\mathrm{\Phi }}`$ (10)
$`(\stackrel{0}{\lambda }\underset{\mathrm{\Gamma }}{\overset{0}{H}})|\stackrel{2}{\mathrm{\Phi }}`$ $`=`$ $`(\underset{\mathrm{\Gamma }}{\overset{1}{H}}\stackrel{1}{\lambda })|\stackrel{1}{\mathrm{\Phi }}\stackrel{2}{\lambda }|\stackrel{0}{\mathrm{\Phi }}`$ (12)
$`\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
$`(\stackrel{0}{\lambda }\underset{\mathrm{\Gamma }}{\overset{0}{H}})|\stackrel{n}{\mathrm{\Phi }}`$ $`=`$ $`(\underset{\mathrm{\Gamma }}{\overset{1}{H}}\stackrel{1}{\lambda })|\stackrel{n1}{\mathrm{\Phi }}\stackrel{2}{\lambda }|\stackrel{n2}{\mathrm{\Phi }}\mathrm{}\stackrel{n}{\lambda }|\stackrel{0}{\mathrm{\Phi }}`$ (14)
$`\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
Analogously, for the left eigenvalue problem we have
$`\mathrm{\Psi }|H_\mathrm{\Gamma }=\lambda \mathrm{\Psi }|,\mathrm{\Psi }|=\stackrel{0}{\mathrm{\Psi }}|+\stackrel{1}{\mathrm{\Psi }}|+\stackrel{2}{\mathrm{\Psi }}|+\mathrm{},\lambda =\stackrel{0}{\lambda }+\stackrel{1}{\lambda }+\stackrel{2}{\lambda }+\mathrm{},`$
and we obtain
$`\stackrel{0}{\mathrm{\Psi }}|(\stackrel{0}{\lambda }\underset{\mathrm{\Gamma }}{\overset{0}{H}})`$ $`=`$ $`0`$ (15)
$`\stackrel{1}{\mathrm{\Psi }}|(\stackrel{0}{\lambda }\underset{\mathrm{\Gamma }}{\overset{0}{H}})`$ $`=`$ $`\stackrel{0}{\mathrm{\Psi }}|(\underset{\mathrm{\Gamma }}{\overset{1}{H}}\stackrel{1}{\lambda })`$ (16)
$`\stackrel{2}{\mathrm{\Psi }}|(\stackrel{0}{\lambda }\underset{\mathrm{\Gamma }}{\overset{0}{H}})`$ $`=`$ $`\stackrel{1}{\mathrm{\Psi }}|(\underset{\mathrm{\Gamma }}{\overset{1}{H}}\stackrel{1}{\lambda })\stackrel{0}{\mathrm{\Psi }}|\stackrel{2}{\lambda }`$ (18)
$`\mathrm{}\mathrm{}\mathrm{}`$
$`\stackrel{n}{\mathrm{\Psi }}|(\stackrel{0}{\lambda }\underset{\mathrm{\Gamma }}{\overset{0}{H}})`$ $`=`$ $`\stackrel{n1}{\mathrm{\Psi }}|(\underset{\mathrm{\Gamma }}{\overset{1}{H}}\stackrel{1}{\lambda })\stackrel{n2}{\mathrm{\Psi }}|\stackrel{2}{\lambda }\mathrm{}\stackrel{0}{\mathrm{\Psi }}|\stackrel{n}{\lambda }`$ (20)
$`\mathrm{}\mathrm{}\mathrm{}`$
It is convenient to define the projectors
$$P_d|11|,P_\mathrm{\Gamma }_\mathrm{\Gamma }๐z|zz|,P_z|zz|,$$
(21)
endowed with the following properties
$`P_d+P_\mathrm{\Gamma }`$ $`=`$ $`I_\mathrm{\Gamma },P_dP_d=P_d,P_\mathrm{\Gamma }P_\mathrm{\Gamma }=P_\mathrm{\Gamma },P_zP_z^{}=\delta _\mathrm{\Gamma }(zz^{})P_z,P_dP_\mathrm{\Gamma }=P_\mathrm{\Gamma }P_d=0`$ (22)
$`P_d\underset{\mathrm{\Gamma }}{\overset{0}{H}}`$ $`=`$ $`\underset{\mathrm{\Gamma }}{\overset{0}{H}}P_d=\mathrm{\Omega }P_d,P_\mathrm{\Gamma }\underset{\mathrm{\Gamma }}{\overset{0}{H}}=\underset{\mathrm{\Gamma }}{\overset{0}{H}}P_\mathrm{\Gamma },P_z\underset{\mathrm{\Gamma }}{\overset{0}{H}}=\underset{\mathrm{\Gamma }}{\overset{0}{H}}P_z=zP_z.`$ (23)
With all these equations in hand we can begin our computations.
### A The discrete spectrum.
Let us start computing the eigenvalues and eigenvectors of $`H_\mathrm{\Gamma }`$ corresponding to the unperturbed eigenvalue $`\mathrm{\Omega }`$, so we take $`|\stackrel{0}{\mathrm{\Phi }}=|1`$ and $`\stackrel{0}{\lambda }=\mathrm{\Omega }`$. With this choice equation (9) is verified. Then equation (10) gives
$`{\displaystyle _\mathrm{\Gamma }}๐z|z(\mathrm{\Omega }z)z|\stackrel{1}{\mathrm{\Phi }}={\displaystyle _\mathrm{\Gamma }}๐z|zz|\underset{\mathrm{\Gamma }}{\overset{1}{H}}|1\stackrel{1}{\lambda }|1,`$
and therefore
$`\stackrel{1}{\lambda }=0,P_\mathrm{\Gamma }|\stackrel{1}{\mathrm{\Phi }}={\displaystyle \frac{1}{(\mathrm{\Omega }\underset{\mathrm{\Gamma }}{\overset{0}{H}})}}P_\mathrm{\Gamma }\underset{\mathrm{\Gamma }}{\overset{1}{H}}|1.`$
Equation (12), gives
$`{\displaystyle _\mathrm{\Gamma }}dz|z(\mathrm{\Omega }z)z|\stackrel{2}{\mathrm{\Phi }}=\underset{\mathrm{\Gamma }}{\overset{1}{H}}|\stackrel{1}{\mathrm{\Phi }}\stackrel{2}{\lambda }|1,`$
and therefore
$`\stackrel{2}{\lambda }=1|\underset{\mathrm{\Gamma }}{\overset{1}{H}}|\stackrel{1}{\mathrm{\Phi }},P_\mathrm{\Gamma }|\stackrel{2}{\mathrm{\Phi }}={\displaystyle \frac{1}{(\mathrm{\Omega }\underset{\mathrm{\Gamma }}{\overset{0}{H}})}}P_\mathrm{\Gamma }\underset{\mathrm{\Gamma }}{\overset{1}{H}}|\stackrel{1}{\mathrm{\Phi }}.`$
For $`|\stackrel{0}{\mathrm{\Phi }}=|1`$ and $`\stackrel{0}{\lambda }=\mathrm{\Omega }`$, equations (9)-(14) give no information about $`1|\stackrel{n}{\mathrm{\Phi }}`$ ($`n=1,2,\mathrm{}`$). Thus, we impose the condition $`1|\stackrel{n}{\mathrm{\Phi }}=0`$ ($`n=1,2,\mathrm{}`$), namely the usual choice in quantum perturbation theory. This choice will fix the normalization of the eigenvector. In summary, up to second order, we obtain the right eigenvector and eigenvalue
$`|\mathrm{\Phi }_\mathrm{\Omega }`$ $`=`$ $`|1+\left({\displaystyle \frac{1}{(\mathrm{\Omega }\underset{\mathrm{\Gamma }}{\overset{0}{H}})}}P_\mathrm{\Gamma }\underset{\mathrm{\Gamma }}{\overset{1}{H}}\right)|1+\left({\displaystyle \frac{1}{(\mathrm{\Omega }\underset{\mathrm{\Gamma }}{\overset{0}{H}})}}P_\mathrm{\Gamma }\underset{\mathrm{\Gamma }}{\overset{1}{H}}\right)^2|1,`$ (24)
$`\lambda _\mathrm{\Omega }`$ $`=`$ $`\mathrm{\Omega }+1|\underset{\mathrm{\Gamma }}{\overset{1}{H}}{\displaystyle \frac{1}{(\mathrm{\Omega }\underset{\mathrm{\Gamma }}{\overset{0}{H}})}}P_\mathrm{\Gamma }\underset{\mathrm{\Gamma }}{\overset{1}{H}}|1.`$ (25)
Using the expressions for $`\underset{\mathrm{\Gamma }}{\overset{0}{H}}`$ and $`\underset{\mathrm{\Gamma }}{\overset{1}{H}}`$ given in equations (8), we obtain
$`\lambda _\mathrm{\Omega }=\mathrm{\Omega }+{\displaystyle _\mathrm{\Gamma }}๐z{\displaystyle \frac{V_z\overline{V_{\overline{z}}}}{\mathrm{\Omega }z}}.`$
If $`V_z`$ and $`\overline{V_{\overline{z}}}`$ are assumed to be analytic functions of $`z`$ in the lower half plane, the integral over $`\mathrm{\Gamma }`$ can be transformed into an integral over $`^+`$:
$`\lambda _\mathrm{\Omega }`$ $`=`$ $`\mathrm{\Omega }+{\displaystyle _0^{\mathrm{}}}๐\omega {\displaystyle \frac{V_\omega \overline{V_\omega }}{\mathrm{\Omega }\omega +i0}}=\mathrm{\Omega }+{\displaystyle _0^{\mathrm{}}}๐\omega \left\{i\pi \delta (\omega \mathrm{\Omega })๐ซ\left({\displaystyle \frac{1}{\omega \mathrm{\Omega }}}\right)\right\}\left|V_\omega \right|^2=`$ (26)
$`=`$ $`\mathrm{\Omega }{\displaystyle _0^{\mathrm{}}}๐\omega ๐ซ\left({\displaystyle \frac{1}{\omega \mathrm{\Omega }}}\right)\left|V_\omega \right|^2i\pi \left|V_\mathrm{\Omega }\right|^2,`$ (27)
where we see that the perturbation of the eigenvalue $`\mathrm{\Omega }>0`$ gives, up to second order, a complex eigenvalue $`\lambda _\mathrm{\Omega }`$ with negative imaginary part (see Fig.1). Notice that the imaginary part would not appear if $`\mathrm{\Omega }<0`$.
Starting from $`\stackrel{0}{\mathrm{\Psi }}|=1|`$ and $`\stackrel{0}{\lambda }=\mathrm{\Omega }`$, equations (15)-(20) and the imposed conditions $`\stackrel{n}{\mathrm{\Psi }}|1=0`$ ($`n=1,2,\mathrm{}`$) provide a well defined algorithm to obtain the left eigenvector. E.g. up to second order we obtain
$$\mathrm{\Psi }_\mathrm{\Omega }|=1|+1|\left(\underset{\mathrm{\Gamma }}{\overset{1}{H}}P_\mathrm{\Gamma }\frac{1}{(\mathrm{\Omega }\underset{\mathrm{\Gamma }}{\overset{0}{H}})}\right)+1|\left(\underset{\mathrm{\Gamma }}{\overset{1}{H}}P_\mathrm{\Gamma }\frac{1}{(\mathrm{\Omega }\underset{\mathrm{\Gamma }}{\overset{0}{H}})}\right)^2,$$
(28)
with the same eigenvalue $`\lambda _\mathrm{\Omega }`$. As we are searching for a complete biorthogonal system, it is necessary at every step to define the normalized eigenvectors as
$$|f_\mathrm{\Omega }|\mathrm{\Phi }_\mathrm{\Omega }/\sqrt{\mathrm{\Psi }_\mathrm{\Omega }|\mathrm{\Phi }_\mathrm{\Omega }},\stackrel{~}{f}_\mathrm{\Omega }|\mathrm{\Psi }_\mathrm{\Omega }|/\sqrt{\mathrm{\Psi }_\mathrm{\Omega }|\mathrm{\Phi }_\mathrm{\Omega }},$$
(29)
which satisfy the condition $`\stackrel{~}{f}_\mathrm{\Omega }|f_\mathrm{\Omega }=1`$.
### B The continuous spectrum.
Let us now compute the eigenvalues and eigenvectors of $`H_\mathrm{\Gamma }`$ corresponding to the continuous spectrum of $`\underset{\mathrm{\Gamma }}{\overset{0}{H}}`$. So we now consider equations (9)-(14) with $`|\stackrel{0}{\mathrm{\Phi }}=|u`$ and $`\stackrel{0}{\lambda }=u`$ ($`u\mathrm{\Gamma }`$). Equation (9) is satisfied, and equation (10) gives
$`{\displaystyle _\mathrm{\Gamma }}๐z|z(uz)z|\stackrel{1}{\mathrm{\Phi }}+(u\mathrm{\Omega })|11|\stackrel{1}{\mathrm{\Phi }}={\displaystyle _\mathrm{\Gamma }}๐z|zz|\underset{\mathrm{\Gamma }}{\overset{1}{H}}|u+|11|\underset{\mathrm{\Gamma }}{\overset{1}{H}}|u\stackrel{1}{\lambda }{\displaystyle _\mathrm{\Gamma }}๐z|z\delta _\mathrm{\Gamma }(zu),`$
or equivalently
$$(uz)z|\stackrel{1}{\mathrm{\Phi }}=z|\underset{\mathrm{\Gamma }}{\overset{1}{H}}|u\stackrel{1}{\lambda }\delta _\mathrm{\Gamma }(zu),(u\mathrm{\Omega })1|\stackrel{1}{\mathrm{\Phi }}=1|\underset{\mathrm{\Gamma }}{\overset{1}{H}}|u.$$
(30)
If we assume that $`V_{zu}`$ is an analytic function of the two variables $`z`$ and $`u`$, the term $`z|\underset{\mathrm{\Gamma }}{\overset{1}{H}}|u=V_{zu}`$ does not include a term proportional to $`\delta _\mathrm{\Gamma }(zu)`$, and therefore equations (30) give
$$\stackrel{1}{\lambda }=0P_\mathrm{\Gamma }|\stackrel{1}{\mathrm{\Phi }}=\frac{1}{(u+i0\underset{\mathrm{\Gamma }}{\overset{0}{H}})}P_\mathrm{\Gamma }\underset{\mathrm{\Gamma }}{\overset{1}{H}}|u,P_d|\stackrel{1}{\mathrm{\Phi }}=\frac{1}{(u\underset{\mathrm{\Gamma }}{\overset{0}{H}})}P_d\underset{\mathrm{\Gamma }}{\overset{1}{H}}|u$$
(31)
Equation (12) gives
$`{\displaystyle _\mathrm{\Gamma }}๐z|z(uz)z|\stackrel{2}{\mathrm{\Phi }}+(u\mathrm{\Omega })|11|\stackrel{2}{\mathrm{\Phi }}={\displaystyle _\mathrm{\Gamma }}๐z|zz|\underset{\mathrm{\Gamma }}{\overset{1}{H}}|\stackrel{1}{\mathrm{\Phi }}+|11|\underset{\mathrm{\Gamma }}{\overset{1}{H}}|\stackrel{1}{\mathrm{\Phi }}\stackrel{2}{\lambda }{\displaystyle _\mathrm{\Gamma }}๐z|z\delta _\mathrm{\Gamma }(zu),`$
and therefore
$$(uz)z|\stackrel{2}{\mathrm{\Phi }}=z|\underset{\mathrm{\Gamma }}{\overset{1}{H}}|\stackrel{1}{\mathrm{\Phi }}\stackrel{2}{\lambda }\delta _\mathrm{\Gamma }(zu),(u\mathrm{\Omega })1|\stackrel{2}{\mathrm{\Phi }}=1|\underset{\mathrm{\Gamma }}{\overset{1}{H}}|\stackrel{1}{\mathrm{\Phi }}.$$
(32)
Using (8) and (31) we can compute
$`z|\underset{\mathrm{\Gamma }}{\overset{1}{H}}|\stackrel{1}{\mathrm{\Phi }}={\displaystyle _\mathrm{\Gamma }}๐z^{}{\displaystyle \frac{V_{zz^{}}V_{z^{}u}}{u+i0z^{}}}+{\displaystyle \frac{V_z\overline{V_{\overline{z}}}}{u\mathrm{\Omega }}},`$
which does not include any factor proportional to $`\delta _\mathrm{\Gamma }(zu)`$. Therefore, equations (32) give
$`\stackrel{2}{\lambda }=0P_\mathrm{\Gamma }|\stackrel{2}{\mathrm{\Phi }}={\displaystyle \frac{1}{(u+i0\underset{\mathrm{\Gamma }}{\overset{0}{H}})}}P_\mathrm{\Gamma }\underset{\mathrm{\Gamma }}{\overset{1}{H}}|\stackrel{1}{\mathrm{\Phi }},P_d|\stackrel{2}{\mathrm{\Phi }}={\displaystyle \frac{1}{(u\underset{\mathrm{\Gamma }}{\overset{0}{H}})}}P_d\underset{\mathrm{\Gamma }}{\overset{1}{H}}|\stackrel{1}{\mathrm{\Phi }}`$
In summary, up to second order, we obtain the right eigenvector
$`|f_u`$ $`=`$ $`|u+({\displaystyle \frac{1}{(u+i0\underset{\mathrm{\Gamma }}{\overset{0}{H}})}}P_\mathrm{\Gamma }\underset{\mathrm{\Gamma }}{\overset{1}{H}}+{\displaystyle \frac{1}{(u\underset{\mathrm{\Gamma }}{\overset{0}{H}})}}P_d\underset{\mathrm{\Gamma }}{\overset{1}{H}})|u+`$ (34)
$`+({\displaystyle \frac{1}{(u+i0\underset{\mathrm{\Gamma }}{\overset{0}{H}})}}P_\mathrm{\Gamma }\underset{\mathrm{\Gamma }}{\overset{1}{H}}+{\displaystyle \frac{1}{(u\underset{\mathrm{\Gamma }}{\overset{0}{H}})}}P_d\underset{\mathrm{\Gamma }}{\overset{1}{H}})^2|u,`$
with eigenvalue $`\lambda _u=u`$. Notice that we have included the factor $`i0`$, namely we make the usual choice of outgoing solutions of scattering theory, to avoid the continuous-continuous resonances.
Equations (15)-(20) with $`\stackrel{0}{\mathrm{\Psi }}|=u|`$ and $`\stackrel{0}{\lambda }=u\mathrm{\Gamma }`$ are used to obtain the corresponding left eigenvector. Up to second order we obtain the left eigenvector
$`\stackrel{~}{f}_u|`$ $`=`$ $`u|+u|(\underset{\mathrm{\Gamma }}{\overset{1}{H}}P_\mathrm{\Gamma }{\displaystyle \frac{1}{(ui0\underset{\mathrm{\Gamma }}{\overset{0}{H}})}}+\underset{\mathrm{\Gamma }}{\overset{1}{H}}P_d{\displaystyle \frac{1}{(u\underset{\mathrm{\Gamma }}{\overset{0}{H}})}})+`$ (36)
$`u|(\underset{\mathrm{\Gamma }}{\overset{1}{H}}P_\mathrm{\Gamma }{\displaystyle \frac{1}{(ui0\underset{\mathrm{\Gamma }}{\overset{0}{H}})}}+\underset{\mathrm{\Gamma }}{\overset{1}{H}}P_d{\displaystyle \frac{1}{(u\underset{\mathrm{\Gamma }}{\overset{0}{H}})}})^2.`$
It can be easily verified that the vectors given in equations (29)-(36) satisfy, up to second order, the orthogonality relations
$$\stackrel{~}{f}_\mathrm{\Omega }|f_\mathrm{\Omega }=1,\stackrel{~}{f}_u|f_u^{}=\delta _\mathrm{\Gamma }(uu^{}),\stackrel{~}{f}_\mathrm{\Omega }|f_u=\stackrel{~}{f}_u|f_\mathrm{\Omega }=0.$$
(37)
They also verify, up to second order, the โcompleteness relationโ
$$\mathrm{\Psi }|\mathrm{\Phi }=\mathrm{\Psi }|I_\mathrm{\Gamma }|\mathrm{\Phi },I_\mathrm{\Gamma }|11|+_\mathrm{\Gamma }๐u|uu|=|f_\mathrm{\Omega }\stackrel{~}{f}_\mathrm{\Omega }|+_\mathrm{\Gamma }๐u|f_u\stackrel{~}{f}_u|.$$
(38)
For the operators $`H`$ and $`H_\mathrm{\Gamma }`$ defined in equations (3) and (8) we also have
$$\mathrm{\Psi }|H|\mathrm{\Phi }=\mathrm{\Psi }|H_\mathrm{\Gamma }|\mathrm{\Phi },H_\mathrm{\Gamma }=\lambda _\mathrm{\Omega }|f_\mathrm{\Omega }\stackrel{~}{f}_\mathrm{\Omega }|+_\mathrm{\Gamma }๐u\lambda _u|f_u\stackrel{~}{f}_u|,$$
(39)
showing that the eigenvalues and eigenvectors obtained above provide the spectral decomposition of $`H_\mathrm{\Gamma }`$. It should be stressed that the expressions for $`I_\mathrm{\Gamma }`$ and $`H_\mathrm{\Gamma }`$ given in equations (38) and (39)are meaningful only when they are used to compute $`\mathrm{\Psi }|I_\mathrm{\Gamma }|\mathrm{\Phi }`$ or $`\mathrm{\Psi }|H_\mathrm{\Gamma }|\mathrm{\Phi }`$ for vectors $`\mathrm{\Psi }`$ ($`\mathrm{\Phi }`$) for which $`\psi (\omega )=\omega |\mathrm{\Psi }`$ ($`\phi (\omega )=\omega |\mathrm{\Phi }`$) have a well defined analytic extension to the complex upper (lower) half plane.
### C The time evolution.
Let us make a first application of the above results. The probability $`P_t`$ of the state vector $`\mathrm{\Phi }_t=e^{iHt}\mathrm{\Phi }_0`$ at the time $`t`$ of being measured in the state $`\mathrm{\Psi }`$ is given by
$$P_t=\left|\mathrm{\Psi }|\mathrm{exp}(iHt)|\mathrm{\Phi }_0\right|^2=\left|\mathrm{\Psi }|\mathrm{exp}(iH_\mathrm{\Gamma }t)|\mathrm{\Phi }_0\right|^2,$$
(40)
where
$$\mathrm{\Psi }|\mathrm{exp}(iH_\mathrm{\Gamma }t)|\mathrm{\Phi }_0=\mathrm{exp}(i\lambda _\mathrm{\Omega }t)\mathrm{\Psi }|f_\mathrm{\Omega }\stackrel{~}{f}_\mathrm{\Omega }|\mathrm{\Phi }_0+_\mathrm{\Gamma }๐u\mathrm{exp}(iut)\mathrm{\Psi }|f_u\stackrel{~}{f}_u|\mathrm{\Phi }_0.$$
(41)
If the exact expressions for the generalized eigenvalues and eigenvectors were known, equations (40) and (41) would give $`P_t`$ exactly. An approximated expression will be obtained if eigenvalues and eigenvectors are given up to some finite order in the perturbation expansion.
Let us consider the survival amplitude of the unstable state $`|1`$, given by
$`1|\mathrm{exp}(iH_\mathrm{\Gamma }t)|1=\mathrm{exp}(i\lambda _\mathrm{\Omega }t)1|f_\mathrm{\Omega }\stackrel{~}{f}_\mathrm{\Omega }|1+{\displaystyle _\mathrm{\Gamma }}๐u\mathrm{exp}(iut)1|f_u\stackrel{~}{f}_u|1.`$
For a small coupling between the discrete and the continuous spectrum, i.e. if in the expression for $`\stackrel{1}{H}`$ given in equation (3) there is a small multiplicative parameter $`\epsilon `$ included in $`V_\omega `$, and using the explicit expressions of the generalized eigenvectors given in sections IIIA and IIIB, we obtain
$$1|f_\mathrm{\Omega }\stackrel{~}{f}_\mathrm{\Omega }|1=1+O(\epsilon ^2),1|f_u\stackrel{~}{f}_u|1=O(\epsilon ^2).$$
(42)
If $`\lambda _\mathrm{\Omega }`$ is computed up to second order from equation (27), we have
$$\left|1|\mathrm{exp}(iHt)|1\right|^2\mathrm{exp}(2\pi V_\mathrm{\Omega }^2t).$$
(43)
This expression is valid if $`\epsilon ^21`$ and $`\epsilon ^2t1`$, i.e. for small coupling between the discrete and the continuous spectrums and for not too large times. This is the standard approximated expression for a decay process which can be found in several text books on quantum mechanics . The deviations from exponential behavior (Zeno and Khalfin effects for shot and long times respectively), are not shown in this approximation. Concerning Zeno effect, notice that the time derivative for $`t=0`$ of expression (43), is of order $`\epsilon ^2`$. If we compute the survival probability keeping the second order terms coming from equations (42), we can obtain the Zeno effect, because the initial time derivative at $`t=0`$ turns out to be of order $`\epsilon ^4`$.
### D Tunneling through a barrier.
We now consider some particular models. Let us first consider the one dimensional rectangular well problem with the following Hamiltonian, given in coordinate representation
$$\stackrel{}{\stackrel{0}{H}}=\frac{\mathrm{}^2}{2\mu }\frac{d^2}{dx^2}+\stackrel{0}{V}(x),$$
(44)
where
$$\stackrel{0}{V}(x)=\begin{array}{c}V_0|x|>a\hfill \\ 0|x|<a.\hfill \end{array}$$
(45)
The spectral decomposition of $`\stackrel{}{\stackrel{0}{H}}`$ is well known. The discrete part of the spectrum has eigenvalues $`E_j`$ between $`0`$ and $`V_0`$. In coordinate representation, the corresponding normalized eigenvectors $`x|j`$ are of the form $`\mathrm{exp}\left(\sqrt{\frac{2\mu }{\mathrm{}^2}(V_0E_j)}x\right)`$, the minus (plus) sign corresponding to $`x>a`$ ($`x<a`$). Therefore the vector states $`x|j`$ are concentrated inside the rectangular well. We are going to consider the simple case in which the parameters $`a`$ and $`V_0`$ verify the condition $`a\sqrt{\frac{2\mu V_0}{\mathrm{}^2}}<\frac{\pi }{2}`$, where there is a single discrete eigenvalue $`E_1`$, and the corresponding eigenvector $`x|1`$ is an even function of the coordinate $`x`$. There is also a continuous spectrum with generalized eigenvalues $`E_k=\frac{\mathrm{}^2}{2\mu }k^2+V_0`$ ($`\mathrm{}<k<+\mathrm{}`$). The generalized eigenvectors $`|k`$ can be obtained as the solutions of the Lipmann-Schwinger equation. The vectors $`|k`$ and $`|1`$ form an orthonormal basis
$$k|k^{}=\delta (kk^{}),1|1=1,1|k=k|1=0,I=|11|+๐k|kk|.$$
(46)
We can now add to $`\stackrel{}{\stackrel{0}{H}}`$ the potential
$$\stackrel{}{\stackrel{1}{V}}(x)=\begin{array}{c}0|x|<b\hfill \\ V_1|x|>b,\hfill \end{array}$$
(47)
with $`ba`$ and $`0<V_1<V_0`$. Thus, a barrier of height $`V_0V_1`$ and length $`ba`$ appears. The potential $`\stackrel{}{\stackrel{1}{V}}`$ can be written in terms of the basis given in (46)
$`\stackrel{}{\stackrel{1}{V}}=|1V_{11}1|+{\displaystyle }dk^{}|1V_{1k^{}}k^{}|+{\displaystyle }dkV_{k1}|k1|+{\displaystyle }dk{\displaystyle }dk^{}|k\stackrel{~}{V}_{kk^{}}k^{}|,`$
where
$`V_{11}=V_1\left[{\displaystyle _{\mathrm{}}^b}๐x1|xx|1+{\displaystyle _{+b}^+\mathrm{}}๐x1|xx|1\right],`$
$`V_{1k}=\overline{V_{k1}}=V_1\left[{\displaystyle _{\mathrm{}}^b}๐x1|xx|k+{\displaystyle _{+b}^+\mathrm{}}๐x1|xx|k\right],`$
$`\stackrel{~}{V}_{kk^{}}=V_1\delta (kk^{})+V_{kk^{}},V_{kk^{}}=V_1{\displaystyle _b^{+b}}๐xk|xx|k^{}.`$
Due to the exponential decrease of $`x|1`$ when $`x\pm \mathrm{}`$, together with the assumption $`ba`$, it can be expected that $`V_{11}`$ and $`V_{1k}`$ be small.
The total Hamiltonian $`H=\stackrel{}{\stackrel{0}{H}}+\stackrel{}{\stackrel{1}{V}}`$ can be also given by $`H=\stackrel{0}{H}+\stackrel{1}{H}`$ where
$`\stackrel{0}{H}`$ $`=`$ $`(E_1+V_{11})|11|+{\displaystyle _{\mathrm{}}^+\mathrm{}}๐k(E_kV_1)|kk|,`$
$`\stackrel{1}{H}`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}๐k^{}|1V_{1k^{}}k^{}|+{\displaystyle _{\mathrm{}}^+\mathrm{}}๐kV_{k1}|k1|+{\displaystyle _{\mathrm{}}^+\mathrm{}}๐k{\displaystyle _{\mathrm{}}^+\mathrm{}}๐k^{}|kV_{kk^{}}k^{}|.`$
Our interest is to describe the time evolution of the unstable state $`|1`$, which is an even function in coordinate representation. Only the even part of the Hamiltonian is relevant to this process
$`\underset{even}{\overset{0}{H}}`$ $`=`$ $`(E_1+V_{11})|11|+{\displaystyle _0^+\mathrm{}}๐k(E_kV_1)|k_{even}k_{even}|,`$
$`\underset{even}{\overset{1}{H}}`$ $`=`$ $`{\displaystyle _0^+\mathrm{}}๐k^{}|1V_{1k^{}}k_{even}^{}|+{\displaystyle _0^+\mathrm{}}๐kV_{k1}|k_{even}1|+{\displaystyle _0^+\mathrm{}}๐k{\displaystyle _0^+\mathrm{}}๐k^{}|k_{even}V_{kk^{}}k_{even}^{}|,`$
where $`|k_{even}=\frac{1}{\sqrt{2}}(|k+|k)`$. If the integrals over $`^+`$ are deformed into integrals over a curve $`\mathrm{\Gamma }`$ in the lower complex half plane, and if $`V_{11}`$ and $`V_{1k}`$ are small, we can use the perturbation method of the beginning of this section. The correction to the discrete eigenvalue $`(E_1+V_{11})`$ of $`\underset{even}{\overset{0}{H}}`$ due to the interaction is
$`1|\underset{even}{\overset{1}{H}}{\displaystyle \frac{1}{E_1+V_{11}\underset{even}{\overset{0}{H}}+i0}}\underset{even}{\overset{1}{H}}|1`$
$`=`$ $`{\displaystyle _0^+\mathrm{}}{\displaystyle \frac{V_{1k}V_{k1}}{(E_1+V_{11})(\frac{\mathrm{}^2k^2}{2\mu }+V_0V_1)+i0}}๐k=`$
$`=`$ $`i\pi {\displaystyle _0^+\mathrm{}}V_{1k}V_{k1}\delta \left[(E_1+V_{11})({\displaystyle \frac{\mathrm{}^2k^2}{2\mu }}+V_0V_1)\right]๐k{\displaystyle _0^+\mathrm{}}V_{1k}V_{k1}๐ซ\left[{\displaystyle \frac{1}{(E_1+V_{11})(\frac{\mathrm{}^2k^2}{2\mu }+V_0V_1)}}\right]๐k.`$
If $`(E_1+V_{11})>(V_0V_1)`$ the corrected eigenvalue has an imaginary part $`i\frac{\pi \mu }{\mathrm{}^2\stackrel{~}{k}}V_{1\stackrel{~}{k}}V_{\stackrel{~}{k}1}`$, where
$`\stackrel{~}{k}=\sqrt{{\displaystyle \frac{2\mu }{\mathrm{}^2}}\left[(E_1+V_{11})(V_0V_1)\right]},`$
and the survival probability of the unstable state $`|1`$ is approximately given by
$`\left|1|\mathrm{exp}(iHt)|1\right|^2\mathrm{exp}\left({\displaystyle \frac{2\pi \mu }{\mathrm{}^2\stackrel{~}{k}}}V_{1\stackrel{~}{k}}V_{\stackrel{~}{k}1}\right)`$
### E Friedrichs model.
We will finally check our method with the well known solvable Friedrichs model. This model can be obtained as a special case of equation (3), with $`V_{\omega \omega ^{}}=0`$. The Hamiltonian is $`H=\stackrel{0}{H}+\stackrel{1}{H}`$, where
$$\stackrel{0}{H}=\mathrm{\Omega }|11|+_0^{\mathrm{}}d\omega \omega |\omega \omega |,\stackrel{1}{H}=_0^{\mathrm{}}d\omega \{V_\omega |\omega 1|+\overline{V_\omega }|1\omega |\},\mathrm{\Omega }^+$$
(48)
As we considered in the previous section, the variable $`\omega ^+`$ can be transformed into a variable $`z\mathrm{\Gamma }`$, being $`\mathrm{\Gamma }`$ the curve in the lower complex half plane shown in Fig.1. The corresponding Hamiltonian $`H_\mathrm{\Gamma }`$ is given by
$$H_\mathrm{\Gamma }=\underset{\mathrm{\Gamma }}{\overset{0}{H}}+\underset{\mathrm{\Gamma }}{\overset{1}{H}},\underset{\mathrm{\Gamma }}{\overset{0}{H}}=\mathrm{\Omega }|11|+_\mathrm{\Gamma }dzz|zz|,\underset{\mathrm{\Gamma }}{\overset{1}{H}}=_\mathrm{\Gamma }dzV_z|z1|+_\mathrm{\Gamma }dz\overline{V_{\overline{z}}}|1z|,$$
(49)
where $`V_z`$ and $`\overline{V_{\overline{z}}}`$ are the analytic extensions of $`V_\omega `$ and $`\overline{V_\omega }`$.
It is possible to obtain exact solutions for the right and left eigenvalue problems $`H_\mathrm{\Gamma }|\mathrm{\Phi }=\lambda |\mathrm{\Phi }`$ and $`\mathrm{\Psi }|H_\mathrm{\Gamma }=\lambda \mathrm{\Psi }|`$, where $`|\mathrm{\Phi }=|11|\mathrm{\Phi }+_\mathrm{\Gamma }๐zz|zz|\mathrm{\Phi }`$ and $`\mathrm{\Psi }|=\mathrm{\Psi }|11|+_\mathrm{\Gamma }๐z\mathrm{\Psi }|zz|`$. If for the sake of simplicity
$$\eta (\lambda )\lambda \mathrm{\Omega }_\mathrm{\Gamma }๐z\frac{V_z\overline{V_{\overline{z}}}}{\lambda z}$$
(50)
verifies $`\eta (\lambda _\mathrm{\Omega })=0`$ for just a single complex number $`\lambda _\mathrm{\Omega }`$ located in the region of the complex plane between $`^+`$ and the curve $`\mathrm{\Gamma }`$, the following spectral decomposition is obtained
$`H_\mathrm{\Gamma }=\lambda _\mathrm{\Omega }|f_\mathrm{\Omega }\stackrel{~}{f}_\mathrm{\Omega }|+{\displaystyle _\mathrm{\Gamma }}๐uu|f_u\stackrel{~}{f}_u|,`$
where
$`|f_\mathrm{\Omega }`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\eta ^{}(\lambda _\mathrm{\Omega })}}}\left\{|1+{\displaystyle _\mathrm{\Gamma }}๐z{\displaystyle \frac{V_z}{\lambda _\mathrm{\Omega }z}}|z\right\}`$ (51)
$`\stackrel{~}{f}_\mathrm{\Omega }|`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\eta ^{}(\lambda _\mathrm{\Omega })}}}\left\{1|+{\displaystyle _\mathrm{\Gamma }}๐z{\displaystyle \frac{\overline{V_{\overline{z}}}}{\lambda _\mathrm{\Omega }z}}z|\right\}`$ (52)
$`|f_u`$ $`=`$ $`|u+{\displaystyle \frac{\overline{V_{\overline{u}}}}{\eta (u+i\epsilon )}}\left\{|1+{\displaystyle _\mathrm{\Gamma }}๐z{\displaystyle \frac{V_z}{(u+i0)z}}|z\right\}`$ (53)
$`\stackrel{~}{f}_u|`$ $`=`$ $`u|+{\displaystyle \frac{V_u}{\eta (ui0)}}\left\{1|+{\displaystyle _\mathrm{\Gamma }}๐z{\displaystyle \frac{\overline{V_{\overline{z}}}}{(ui0)z}}z|\right\},u\mathrm{\Gamma }.`$ (54)
These generalized eigenvectors form a complete biorthogonal system , in the sense given in equations (37) and (38).
Using equations (25)-(29) we can obtain the approximated eigenvectors
$`|f_\mathrm{\Omega }`$ $``$ $`|1+{\displaystyle _\mathrm{\Gamma }}๐z{\displaystyle \frac{V_z}{\mathrm{\Omega }z}}|z{\displaystyle \frac{1}{2}}{\displaystyle _\mathrm{\Gamma }}๐z{\displaystyle \frac{V_z\overline{V_{\overline{z}}}}{(\mathrm{\Omega }z)^2}}|1`$ (55)
$`\stackrel{~}{f}_\mathrm{\Omega }|`$ $``$ $`1|+{\displaystyle _\mathrm{\Gamma }}๐z{\displaystyle \frac{\overline{V_{\overline{z}}}}{\mathrm{\Omega }z}}z|{\displaystyle \frac{1}{2}}{\displaystyle _\mathrm{\Gamma }}๐z{\displaystyle \frac{\overline{V_{\overline{z}}}V_z}{(\mathrm{\Omega }z)^2}}1|,`$ (56)
with the eigenvalue
$$\lambda _\mathrm{\Omega }\mathrm{\Omega }+_\mathrm{\Gamma }๐z\frac{V_z\overline{V_{\overline{z}}}}{\mathrm{\Omega }z}=\mathrm{\Omega }_0^{\mathrm{}}๐\omega ๐ซ\frac{1}{\omega \mathrm{\Omega }}\left|V_\omega \right|^2i\pi \left|V_\mathrm{\Omega }\right|^2.$$
(57)
Using equations (34) and (36) we also obtain, up to second order
$`|f_u`$ $``$ $`|u+{\displaystyle \frac{\overline{V_{\overline{u}}}}{u\mathrm{\Omega }}}|1+{\displaystyle \frac{\overline{V_{\overline{u}}}}{u\mathrm{\Omega }}}{\displaystyle _\mathrm{\Gamma }}๐z{\displaystyle \frac{V_z}{(uz+i\epsilon )}}|z`$ (58)
$`\stackrel{~}{f}_u|`$ $``$ $`u|+{\displaystyle \frac{V_u}{u\mathrm{\Omega }}}1|+{\displaystyle \frac{V_u}{u\mathrm{\Omega }}}{\displaystyle _\mathrm{\Gamma }}๐z{\displaystyle \frac{\overline{V_{\overline{z}}}}{(uzi\epsilon )}}z|,`$ (59)
with eigenvalue $`\lambda _u=u`$ ($`u\mathrm{\Gamma }`$).
In order to compare these results with the approximate eigenvalues and eigenvectors obtained in equations (56)-(59), we need an expansion of the exact expressions (54) in powers of the interaction parameter. We can use equation (50) to obtain the proper eigenvalue $`\lambda _\mathrm{\Omega }`$, satisfying $`\eta (\lambda _\mathrm{\Omega })=0`$, from the iterative equation
$`\stackrel{n+1}{\lambda _\mathrm{\Omega }}=\mathrm{\Omega }+{\displaystyle _\mathrm{\Gamma }}dz{\displaystyle \frac{V_z\overline{V_{\overline{z}}}}{\left(\stackrel{n}{\lambda _\mathrm{\Omega }}z\right)}},,n=0,1,2,\mathrm{},\stackrel{0}{\lambda }=\mathrm{\Omega }.`$
For small values of the interaction parameter, the secuence $`\stackrel{n}{\lambda _\mathrm{\Omega }}`$ converge to a single solution $`\lambda _\mathrm{\Omega }`$ of $`\eta (\lambda _\mathrm{\Omega })=0`$. Up to second order we have
$$\lambda _\mathrm{\Omega }\mathrm{\Omega }+_\mathrm{\Gamma }๐z\frac{V_z\overline{V_{\overline{z}}}}{\mathrm{\Omega }z}.$$
(60)
Up to second order in the interaction parameter, we also have
$$\frac{1}{\sqrt{\eta ^{}(\lambda _\mathrm{\Omega })}}1\frac{1}{2}_\mathrm{\Gamma }๐z\frac{V_z\overline{V_{\overline{z}}}}{(\mathrm{\Omega }z)^2},\frac{1}{\lambda _\mathrm{\Omega }z}=\frac{1}{\mathrm{\Omega }z}+O(V^2),\frac{1}{\eta (u\pm i\epsilon )}=\frac{1}{u\mathrm{\Omega }}+O(V^2)$$
(61)
Replacing (60) and (61) in (54) we obtain the approximated expressions (56)-(59), so our approximated algorithm reproduces the results of the solvable model.
## III Generalized states and singular observables.
Up to now we have only considered pure states and observables which can be represented by projectors constructed with normalized vectors. The extension of the formalism to mixed states, represented by density operators of the form $`\widehat{\rho }=_jp_j|\mathrm{\Phi }_j\mathrm{\Phi }_j|`$ ($`\mathrm{\Phi }_j|\mathrm{\Phi }_j=1`$, $`p_j0`$, $`_jp_j=1`$) and observables of the same form ($`\widehat{O}=_jO_j|\mathrm{\Psi }_j\mathrm{\Psi }_j|`$) is straightforward. However, singular observables as $`_{\omega _1}^{\omega _2}๐\omega |\omega \omega |`$, which are not operators of the form just mentioned, normally appear, and they can not be considered in the framework of the perturbation method presented in the previous section or its straightforward generalization. In this case, it is possible to consider the states as functionals acting on the set of observables represented by operators including singular terms . These generalized states are also useful to study the approach to equilibrium in the thermodynamic limit . The corresponding perturbation method is therefore a non trivial generalization of the one in section II. This generalization will be the subject of this section and, as we will see, contains new and interesting features.
Let us consider the Friedrichs model, presented in section II. The Hamiltonian is given by equation (48) and we assume, for simplicity, that $`V_\omega =\overline{V_\omega }`$.
For this model we are going to consider the class of observables which can be represented by operators of the form
$$O=O_1|11|+_0^{\mathrm{}}๐\omega O_\omega |\omega \omega |+_0^{\mathrm{}}๐\omega _0^{\mathrm{}}๐\omega ^{}O_{\omega \omega ^{}}|\omega \omega ^{}|+_0^{\mathrm{}}๐\omega O_{\omega 1}|\omega 1|+_0^{\mathrm{}}๐\omega ^{}O_{1\omega ^{}}|1\omega ^{}|$$
(62)
where $`O_1=\overline{O_1}`$, $`O_\omega =\overline{O_\omega }`$, $`O_{\omega \omega ^{}}=\overline{O_{\omega ^{}\omega }}`$, and $`O_{\omega 1}=\overline{O_{1\omega }}`$, being $`O_\omega `$, $`O_{\omega \omega ^{}}`$, $`O_{\omega 1}`$ and $`O_{1\omega ^{}}`$ regular functions of the variables $`\omega `$ and $`\omega ^{}`$. The second term of the r.h.s. of the last equation is a singular operator.
In the Heisemberg representation, the time evolution of an observable is given by <sup>*</sup><sup>*</sup>*Superoperators, i.e. operators acting on operators, will be written in blackboard bold types.
$`i{\displaystyle \frac{d}{dt}}O_t=๐O_t,๐O[H,O],`$
where $`๐`$ is the Liouville-Von Newmann superoperator.
It is convenient to introduce the notation $`|O)O`$ and to define the generalized operators
$$|1)|11|,|\omega )|\omega \omega |,|\omega \omega ^{})|\omega \omega ^{}|,|\omega 1)|\omega 1|,|1\omega ^{})|1\omega ^{}|.$$
(63)
Therefore we can write
$$|O)=O_1|1)+_0^{\mathrm{}}d\omega O_\omega |\omega )+_0^{\mathrm{}}d\omega _0^{\mathrm{}}d\omega ^{}O_{\omega \omega ^{}}|\omega \omega ^{})+_0^{\mathrm{}}d\omega O_{\omega 1}|\omega 1)+_0^{\mathrm{}}d\omega ^{}O_{1\omega ^{}}|1\omega ^{}).$$
(64)
When the states of the system are represented by the usual density operators $`\widehat{\rho }=\widehat{\rho }^{}`$, the mean value of an observable $`O`$ is given by $`O_{\widehat{\rho }}=Tr(\widehat{\rho }^{}O)`$, and any density operator has the following well known properties
$$Tr(\widehat{\rho }^{})=Tr(\widehat{\rho }^{}I)=1,Tr(\widehat{\rho }^{}[\alpha A+\beta B]=\alpha Tr(\widehat{\rho }^{}A)+\beta Tr(\widehat{\rho }^{}B),Tr(\widehat{\rho }^{}B^{})=\overline{Tr(\widehat{\rho }^{}B)},$$
(65)
where $`I`$ is the identity operator ($`I|11|+_0^{\mathrm{}}d\omega |\omega \omega |=|1)+_0^{\mathrm{}}d\omega |\omega )`$).
As we are going to consider the possibility to have more general states than the usual density operators $`\widehat{\rho }`$, we will consider generalized states as functionals $`\rho `$ acting on the space of observables defined in equation (64). The action of a state functional $`\rho `$ on an observable $`O`$ will be denoted by $`(\rho |O)`$ and it gives the mean value $`O_\rho `$ of the observable in the state $`\rho `$ ($`O_\rho (\rho |O)`$). For any state functional $`\rho `$ we assume the following properties
$$(\rho |I)=1,(\rho |\alpha A+\beta B)=\alpha (\rho |A)+\beta (\rho |B),(\rho |B^{})=\overline{(\rho |B)},$$
(66)
which are the generalization of the properties for the density operators $`\widehat{\rho }`$ given in equations (65).
For the class of observables given in equation (64) it is convenient to define the functionals $`(1|`$, $`(\omega |`$, $`(\omega \omega ^{}|`$, $`(\omega 1|`$ and $`(1\omega ^{}|`$, in such a way that
$$(1|O)=O_1,(\omega |O)=O_\omega ,(\omega \omega ^{}|O)=O_{\omega \omega ^{}},(\omega 1|O)=O_{\omega 1},(1\omega ^{}|O)=O_{1\omega ^{}}.$$
(67)
Using (64), (66) and (67) we obtain
$`(\rho |O)`$ $`=`$ $`(\rho |1)(1|O)+{\displaystyle _0^{\mathrm{}}}๐\omega (\rho |\omega )(\omega |O)+{\displaystyle _0^{\mathrm{}}}๐\omega {\displaystyle _0^{\mathrm{}}}๐\omega ^{}(\rho |\omega \omega ^{})(\omega \omega ^{}|O)+`$ (69)
$`+{\displaystyle _0^{\mathrm{}}}๐\omega (\rho |\omega 1)(\omega 1|O)+{\displaystyle _0^{\mathrm{}}}๐\omega ^{}(\rho |1\omega ^{})(1\omega ^{}|O).`$
From this equation we obtain
$`(\rho |O)=(\rho |๐|O),`$
where $`๐`$ is the identity superoperator
$$๐|1)(1|+_0^{\mathrm{}}d\omega |\omega )(\omega |+_0^{\mathrm{}}d\omega _0^{\mathrm{}}d\omega ^{}|\omega \omega ^{})(\omega \omega ^{}|+_0^{\mathrm{}}d\omega |\omega 1)(\omega 1|+_0^{\mathrm{}}d\omega ^{}|1\omega ^{})(1\omega ^{}|.$$
(70)
The generalized state functionals and observables defined in (63) and (67) satisfy the orthogonality conditions
$`(1|1)=1,(\omega |\omega ^{})=\delta (\omega \omega ^{}),`$
$`(\omega \omega ^{}|\xi \xi ^{})=\delta (\omega \xi )\delta (\omega ^{}\xi ^{}),(\omega 1|\xi 1)=\delta (\omega \xi ),(1\omega ^{}|1\xi ^{})=\delta (\omega ^{}\xi ^{}),`$
$`(1|\omega )=(1|\xi \xi ^{})=(1|\xi 1)=(1|1\xi ^{})=0,(\omega \omega ^{}|1)=(\omega \omega ^{}|\xi )=(\omega \omega ^{}|\xi 1)=(\omega \omega ^{}|1\xi ^{})=0,`$
$$(\omega 1|1)=(\omega 1|\xi )=(\omega 1|\xi \xi ^{})=(\omega 1|1\xi ^{})=0,(1\omega ^{}|1)=(1\omega ^{}|\xi )=(1\omega ^{}|\xi \xi ^{})=(1\omega ^{}|\xi 1)=0.$$
(71)
Defining the superoperators $`\stackrel{0}{๐}`$ and $`\stackrel{1}{๐}`$ by
$`\stackrel{0}{๐}O[\stackrel{0}{H},O],\stackrel{1}{๐}O[\stackrel{1}{H},O],๐=\stackrel{0}{๐}+\stackrel{1}{๐},`$
and using the generalized operators and functionals defined in equations (63) and (67), we have
$`\stackrel{0}{๐}`$ $`=`$ $`{\displaystyle }d\omega ^{}(\mathrm{\Omega }\omega ^{})|1\omega ^{})(1\omega ^{}|+{\displaystyle }d\omega (\omega \mathrm{\Omega })|\omega 1)(\omega 1|+{\displaystyle }d\omega {\displaystyle }d\omega ^{}(\omega \omega ^{})|\omega \omega ^{})(\omega \omega ^{}|,`$ (72)
$`\stackrel{1}{๐}`$ $`=`$ $`{\displaystyle }d\omega V_\omega [|\omega 1)|1\omega )](1|+{\displaystyle }d\omega V_\omega [|1\omega )|\omega 1)](\omega |+{\displaystyle }d\omega [V_\omega |1)+{\displaystyle }d\omega ^{}V_\omega ^{}|\omega ^{}\omega )](1\omega |+`$ (74)
$`+{\displaystyle }d\omega [V_\omega |1){\displaystyle }d\omega ^{}V_\omega ^{}|\omega \omega ^{})](\omega 1|+{\displaystyle }d\omega {\displaystyle }d\omega ^{}[V_\omega |1\omega ^{})V_\omega ^{}|\omega 1)](\omega \omega ^{}|.`$
From equation (72) and the orthogonality conditions (71), the generalized states and observables defined in equations (67) and (63) form a complete biorthogonal set of generalized eigenvectors of the โfreeโ superoperator $`\stackrel{0}{๐}`$
$`\stackrel{0}{๐}|1)`$ $`=`$ $`\stackrel{0}{๐}|\omega )=0,\stackrel{0}{๐}|1\omega ^{})=(\mathrm{\Omega }\omega ^{})|1\omega ^{}),\stackrel{0}{๐}|\omega 1)=(\omega \mathrm{\Omega })|\omega 1),\stackrel{0}{๐}|\omega \omega ^{})=(\omega \omega ^{})|\omega \omega ^{}),`$ (75)
$`(1|\stackrel{0}{๐}`$ $`=`$ $`(\omega |\stackrel{0}{๐}=0,(1\omega ^{}|\stackrel{0}{๐}=(\mathrm{\Omega }\omega ^{})(1\omega ^{}|,(\omega 1|\stackrel{0}{๐}=(\omega \mathrm{\Omega })(\omega 1|,(\omega \omega ^{}|\stackrel{0}{๐}=(\omega \omega ^{})(\omega \omega ^{}|,`$ (76)
If we try to construct a complete biorthogonal set of eigenvectors of the Liouville-Von Newmann superoperator $`๐=\stackrel{0}{๐}+\stackrel{1}{๐}`$, depending analytically on the small interaction parameter which we assume included in $`\stackrel{1}{๐}`$, the standard methods of perturbation theory are not applicable: the superposition of the discrete and continuous parts of the spectrum of $`\stackrel{0}{๐}`$ will produce non defined terms in the perturbation expansion.
However, this problem can be avoided if we are satisfied, as in the pure state case, to obtain information about terms of the form $`(\rho |O)`$ or $`(\rho |๐|O)`$, considering states $`\rho `$ and observables $`O`$ for which $`(\rho |\omega \omega ^{})`$, $`(\rho |\omega 1)`$, $`(\rho |1\omega ^{})`$, $`(\omega \omega ^{}|O)`$, $`(\omega 1|O)`$ and $`(1\omega ^{}|O)`$ have well defined analytic extensions $`(\rho |zz^{})`$, $`(\rho |z1)`$, $`(\rho |1z^{})`$, $`(zz^{}|O)`$, $`(z1|O)`$ and $`(1z^{}|O)`$ to complex values of $`z`$ ($`z^{}`$) on the upper (lower) complex half plane. We also consider states $`\rho `$ and observables $`O`$ for which $`(\rho |\omega )`$ and $`(\omega |O)`$ have well defined analytic extensions $`(\rho |z)`$ and $`(z|O)`$ for complex values of $`z`$ near $`^+`$ in the upper and in the lower complex half planes. The physical bases of these analyticity properties are given in papers and .
For these states and observables, the Cauchy theorem can be used in equation (69) to write
$`(\rho |O)`$ $`=`$ $`(\rho |1)(1|O)+{\displaystyle _0^{\mathrm{}}}๐\omega (\rho |\omega )(\omega |O)+{\displaystyle _\mathrm{\Gamma }}๐z{\displaystyle _{\overline{\mathrm{\Gamma }}}}๐z^{}(\rho |zz^{})(zz^{}|O)+`$ (78)
$`+{\displaystyle _\mathrm{\Gamma }}๐z(\rho |z1)(z1|O)+{\displaystyle _{\overline{\mathrm{\Gamma }}}}๐z^{}(\rho |1z^{})(1z^{}|O),`$
where $`\mathrm{\Gamma }`$ is a curve in the lower complex half plane, as shown in Fig. 1, and $`\overline{\mathrm{\Gamma }}`$ is the conjugate curve in the upper complex half plane. Therefore
$$(\rho |O)=(\rho |๐_{ext}|O),๐_{ext}|1)(1|+d\omega |\omega )(\omega |+_{\overline{\mathrm{\Gamma }}}dz_\mathrm{\Gamma }dz^{}|zz^{})(zz^{}|+_{\overline{\mathrm{\Gamma }}}dz|z1)(z1|+_\mathrm{\Gamma }dz^{}|1z^{})(1z^{}|.$$
(79)
If in the expression for $`๐_{ext}`$ given in equation (79), the term $`d\omega |\omega )(\omega |`$ is replaced by $`_{\overline{\mathrm{\Gamma }}}dz|z)(z|`$ or by $`_\mathrm{\Gamma }dz^{}|z^{})(z^{}|`$, the identity $`(\rho |O)=(\rho |๐_{ext}|O)`$ is also verified for the restricted classes of states and observables with the well defined analytic extensions defined above.
We can also define $`๐_{ext}=\underset{ext}{\overset{0}{๐}}+\underset{ext}{\overset{1}{๐}}`$, where
$`\underset{ext}{\overset{0}{๐}}`$ $``$ $`{\displaystyle _\mathrm{\Gamma }}dz^{}(\mathrm{\Omega }z^{})|1z^{})(1z^{}|+{\displaystyle _{\overline{\mathrm{\Gamma }}}}dz(z\mathrm{\Omega })|z1)(z1|+{\displaystyle _{\overline{\mathrm{\Gamma }}}}dz{\displaystyle _\mathrm{\Gamma }}dz^{}(zz^{})|zz^{})(zz^{}|,`$ (80)
$`\underset{ext}{\overset{1}{๐}}`$ $``$ $`{\displaystyle _{\overline{\mathrm{\Gamma }}}}dzV_z|z1)(1|{\displaystyle _\mathrm{\Gamma }}dz^{}V_z^{}|1z^{})(1|+{\displaystyle _\mathrm{\Gamma }}dz^{}V_z^{}|1z^{})(z^{}|{\displaystyle _{\overline{\mathrm{\Gamma }}}}dzV_z|z1)(z|+`$ (83)
$`{\displaystyle _\mathrm{\Gamma }}dz^{}V_z^{}|1)(1z^{}|+{\displaystyle _{\overline{\mathrm{\Gamma }}}}dz{\displaystyle _\mathrm{\Gamma }}dz^{}V_z|zz^{})(1z^{}|+{\displaystyle _{\overline{\mathrm{\Gamma }}}}dzV_z|1)(z1|`$
$`{\displaystyle _{\overline{\mathrm{\Gamma }}}}dz{\displaystyle _\mathrm{\Gamma }}dz^{}V_z^{}|zz^{})(z1|+{\displaystyle _{\overline{\mathrm{\Gamma }}}}dz{\displaystyle _\mathrm{\Gamma }}dz^{}V_z|1z^{})(zz^{}|{\displaystyle _{\overline{\mathrm{\Gamma }}}}dz{\displaystyle _\mathrm{\Gamma }}dz^{}V_z^{}|z1)(zz^{}|`$
The superoperators $`\underset{ext}{\overset{0}{๐}}`$ and $`\underset{ext}{\overset{1}{๐}}`$ verify $`(\rho |\stackrel{0}{๐}|O)=(\rho |\underset{ext}{\overset{0}{๐}}|O)`$ and $`(\rho |\stackrel{1}{๐}|O)=(\rho |\underset{ext}{\overset{1}{๐}}|O)`$. The generalized eigenvalues of $`\underset{ext}{\overset{0}{๐}}`$ are shown in Fig. 2. The form given in equations (80) and (83) for the superoperator $`๐_{ext}=\underset{ext}{\overset{0}{๐}}+\underset{ext}{\overset{1}{๐}}`$ can also be obtained directly from the definition $`๐_{ext}OH_{\overline{\mathrm{\Gamma }}}OOH_\mathrm{\Gamma }`$, where $`H_\mathrm{\Gamma }`$ was defined in section II.
For the right eigenvalue problem $`๐_{ext}|\mathrm{\Phi })=\lambda |\mathrm{\Phi })`$ we assume an expansion
$`|\mathrm{\Phi })=|\stackrel{0}{\mathrm{\Phi }})+|\stackrel{1}{\mathrm{\Phi }})+|\stackrel{2}{\mathrm{\Phi }})+\mathrm{},\lambda =\stackrel{0}{\lambda }+\stackrel{1}{\lambda }+\stackrel{2}{\lambda }+\mathrm{},`$
with respect to the small interaction parameter that we suppose it is contained in $`\stackrel{1}{๐}`$. The following set of equations is obtained order by order
$`(\stackrel{0}{\lambda }\underset{ext}{\overset{0}{๐}})|\stackrel{0}{\mathrm{\Phi }})`$ $`=`$ $`0`$ (84)
$`(\stackrel{0}{\lambda }\underset{ext}{\overset{0}{๐}})|\stackrel{1}{\mathrm{\Phi }})`$ $`=`$ $`(\underset{ext}{\overset{1}{๐}}\stackrel{1}{\lambda })|\stackrel{0}{\mathrm{\Phi }})`$ (85)
$`(\stackrel{0}{\lambda }\underset{ext}{\overset{0}{๐}})|\stackrel{2}{\mathrm{\Phi }})`$ $`=`$ $`(\underset{ext}{\overset{1}{๐}}\stackrel{1}{\lambda })|\stackrel{1}{\mathrm{\Phi }})\stackrel{2}{\lambda }|\stackrel{0}{\mathrm{\Phi }})`$ (87)
$`\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
$`(\stackrel{0}{\lambda }\underset{ext}{\overset{0}{๐}})|\stackrel{n}{\mathrm{\Phi }})`$ $`=`$ $`(\underset{ext}{\overset{1}{๐}}\stackrel{1}{\lambda })|\stackrel{n1}{\mathrm{\Phi }})\stackrel{2}{\lambda }|\stackrel{n2}{\mathrm{\Phi }})\mathrm{}\stackrel{n}{\lambda }|\stackrel{0}{\mathrm{\Phi }})`$ (89)
$`\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
For the left eigenvalue problem $`(\mathrm{\Psi }|๐_{ext}=\lambda (\mathrm{\Psi }|`$, the expansion
$`(\mathrm{\Psi }|=(\stackrel{0}{\mathrm{\Psi }}|+(\stackrel{1}{\mathrm{\Psi }}|+(\stackrel{2}{\mathrm{\Psi }}|+\mathrm{},\lambda =\stackrel{0}{\lambda }+\stackrel{1}{\lambda }+\stackrel{2}{\lambda }+\mathrm{},`$
gives
$`(\stackrel{0}{\mathrm{\Psi }}|(\stackrel{0}{\lambda }\underset{ext}{\overset{0}{๐}})`$ $`=`$ $`0`$ (90)
$`(\stackrel{1}{\mathrm{\Psi }}|(\stackrel{0}{\lambda }\underset{ext}{\overset{0}{๐}})`$ $`=`$ $`(\stackrel{0}{\mathrm{\Psi }}|(\underset{ext}{\overset{1}{๐}}\stackrel{1}{\lambda })`$ (91)
$`(\stackrel{2}{\mathrm{\Psi }}|(\stackrel{0}{\lambda }\underset{ext}{\overset{0}{๐}})`$ $`=`$ $`(\stackrel{1}{\mathrm{\Psi }}|(\underset{ext}{\overset{1}{๐}}\stackrel{1}{\lambda })(\stackrel{0}{\mathrm{\Psi }}|\stackrel{2}{\lambda }`$ (93)
$`\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
$`(\stackrel{n}{\mathrm{\Psi }}|(\stackrel{0}{\lambda }\underset{ext}{\overset{0}{๐}})`$ $`=`$ $`(\stackrel{n1}{\mathrm{\Psi }}|(\underset{ext}{\overset{1}{๐}}\stackrel{1}{\lambda })(\stackrel{n2}{\mathrm{\Psi }}|\stackrel{2}{\lambda }\mathrm{}(\stackrel{0}{\mathrm{\Psi }}|\stackrel{n}{\lambda }`$ (95)
$`\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
It is convenient to define the following projectors
$$_0|1)(1|+d\omega |\omega )(\omega |,_{\overline{\mathrm{\Gamma }}\mathrm{\Gamma }}_{\overline{\mathrm{\Gamma }}}dz_\mathrm{\Gamma }dz^{}|zz^{})(zz^{}|,_{\overline{\mathrm{\Gamma }}}_{\overline{\mathrm{\Gamma }}}dz|z1)(z1|,_\mathrm{\Gamma }_\mathrm{\Gamma }dz^{}|1z^{})(1z^{}|.$$
(96)
These projectors commute with $`\underset{ext}{\overset{0}{๐}}`$ and also satisfy
$`๐_{ext}=_0+_{\overline{\mathrm{\Gamma }}\mathrm{\Gamma }}+_{\overline{\mathrm{\Gamma }}}+_\mathrm{\Gamma },_0\underset{ext}{\overset{0}{๐}}=\underset{ext}{\overset{0}{๐}}_0=0.`$
Therefore $`_0`$ is the projector on the invariant space of the time evolution generated by the superoperator $`\underset{ext}{\overset{0}{๐}}`$. As in the pure state case we will compute the discrete and the continuous spectrum.
### A The discrete spectrum.
We first consider equations (84)-(89) with
$$\stackrel{0}{\lambda }=0,|\stackrel{0}{\mathrm{\Phi }})=_0|\stackrel{0}{\mathrm{\Phi }})=|1)(1|\stackrel{0}{\mathrm{\Phi }})+d\omega |\omega )(\omega |\stackrel{0}{\mathrm{\Phi }}).$$
(97)
Multiplying equation (85) by $`_0`$ and by $`_0๐_{ext}_0`$, we obtain
$`\stackrel{1}{\lambda }|\stackrel{0}{\mathrm{\Phi }})`$ $`=`$ $`_0\underset{ext}{\overset{0}{๐}}_0|\stackrel{0}{\mathrm{\Phi }}),`$ (98)
$`_0|\stackrel{1}{\mathrm{\Phi }})`$ $`=`$ $`{\displaystyle \frac{(1)}{_0\underset{ext}{\overset{0}{๐}}_0}}_0\underset{ext}{\overset{1}{๐}}|\stackrel{0}{\mathrm{\Phi }}).`$ (99)
Equation (87) gives
$`\stackrel{2}{\lambda }|\stackrel{0}{\mathrm{\Phi }})`$ $`=`$ $`\left[_0\underset{ext}{\overset{1}{๐}}{\displaystyle \frac{(1)}{_0\underset{ext}{\overset{0}{๐}}_0}}_0\underset{ext}{\overset{1}{๐}}_0\right]|\stackrel{0}{\mathrm{\Phi }}),`$ (100)
$`_0|\stackrel{2}{\mathrm{\Phi }})`$ $`=`$ $`{\displaystyle \frac{(1)}{_0\underset{ext}{\overset{0}{๐}}_0}}_0[\underset{ext}{\overset{1}{๐}}\stackrel{1}{\lambda }]|\stackrel{1}{\mathrm{\Phi }}).`$ (101)
From the definitions (96) and (83) we have $`_0\underset{ext}{\overset{0}{๐}}_0=0`$ and therefore $`\stackrel{1}{\lambda }=0`$: the degeneration of the space expanded by the projector $`_0`$ is not eliminated by first order corrections. The previous equations give no information about $`_0|\stackrel{n}{\mathrm{\Phi }})`$ ($`n=1,2,\mathrm{}`$), and we make the usual choice $`|\stackrel{n}{\mathrm{\Phi }})=_0|\stackrel{n}{\mathrm{\Phi }})`$ ($`n=1,2,\mathrm{}`$).
But the degeneration is partially removed with the second order correction. Equation (100), an eigenvalue problem for $`\stackrel{2}{\lambda }`$, gives
$$2\pi iV_\mathrm{\Omega }^2|1)(1|\stackrel{0}{\mathrm{\Phi }})2\pi iV_\mathrm{\Omega }^2|1)(\omega =\mathrm{\Omega }|\stackrel{0}{\mathrm{\Phi }})=\stackrel{2}{\lambda }|1)(1|\stackrel{0}{\mathrm{\Phi }})+\stackrel{2}{\lambda }d\omega |\omega )(\omega |\stackrel{0}{\mathrm{\Phi }}).$$
(102)
This equation has the solutions
$$\underset{d}{\overset{2}{\lambda }}=2\pi iV_\mathrm{\Omega }^2,|\underset{d}{\overset{0}{\mathrm{\Phi }}})=|1),\underset{\omega }{\overset{2}{\lambda }}=0,|\underset{\omega }{\overset{0}{\mathrm{\Phi }}})=\delta (\omega \mathrm{\Omega })|1)+|\omega ).$$
(103)
We give in appendix A the details of these calculations.
For the left eigenvalue problem with
$$\stackrel{0}{\lambda }=0,(\stackrel{0}{\mathrm{\Psi }}|=(\stackrel{0}{\mathrm{\Psi }}|_0=(\stackrel{0}{\mathrm{\Psi }}|1)(1|+d\omega (\stackrel{0}{\mathrm{\Psi }}|\omega )(\omega |,$$
(104)
equations (90)-(95) give
$`(\stackrel{0}{\mathrm{\Psi }}|\stackrel{1}{\lambda }`$ $`=`$ $`(\stackrel{0}{\mathrm{\Psi }}|(_0\underset{ext}{\overset{0}{๐}}_0),`$ (105)
$`(\stackrel{1}{\mathrm{\Psi }}|_0`$ $`=`$ $`(\stackrel{0}{\mathrm{\Psi }}|\underset{ext}{\overset{1}{๐}}_0{\displaystyle \frac{(1)}{_0\underset{ext}{\overset{0}{๐}}_0}},`$ (106)
$`(\stackrel{0}{\mathrm{\Psi }}|\stackrel{2}{\lambda }`$ $`=`$ $`(\stackrel{0}{\mathrm{\Psi }}|\left[_0\underset{ext}{\overset{1}{๐}}{\displaystyle \frac{(1)}{_0\underset{ext}{\overset{0}{๐}}_0}}_0\underset{ext}{\overset{1}{๐}}_0\right],`$ (107)
$`(\stackrel{2}{\mathrm{\Psi }}|_0`$ $`=`$ $`(\stackrel{1}{\mathrm{\Psi }}|[\underset{ext}{\overset{1}{๐}}\stackrel{1}{\lambda }]_0{\displaystyle \frac{(1)}{_0\underset{ext}{\overset{0}{๐}}_0}},`$ (109)
$`\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
Once again, the degeneration is not removed by the first order corrections, and $`\stackrel{1}{\lambda }=0`$ as a consequence of $`_0\underset{ext}{\overset{0}{๐}}_0=0`$. As there is no information about $`(\stackrel{n}{\mathrm{\Psi }}|_0`$ for $`n1`$, the usual condition $`(\stackrel{n}{\mathrm{\Psi }}|=(\stackrel{n}{\mathrm{\Psi }}|_0`$ can be imposed. For the eigenvalue problem (107) we obtain
$$2\pi iV_\mathrm{\Omega }^2(\stackrel{0}{\mathrm{\Psi }}|1)(1|2\pi iV_\mathrm{\Omega }^2(\stackrel{0}{\mathrm{\Psi }}|1)(\omega =\mathrm{\Omega }|=\stackrel{2}{\lambda }(\stackrel{0}{\mathrm{\Psi }}|1)(1|+\stackrel{2}{\lambda }d\omega (\stackrel{0}{\mathrm{\Psi }}|\omega )(\omega |.$$
(110)
This equation has the solutions
$$\underset{d}{\overset{2}{\lambda }}=2\pi iV_\mathrm{\Omega }^2,(\underset{d}{\overset{0}{\mathrm{\Psi }}}|=(1|(\omega =\mathrm{\Omega }|,\underset{\omega }{\overset{2}{\lambda }}=0,(\underset{\omega }{\overset{0}{\mathrm{\Psi }}}|=(\omega |.$$
(111)
(See appendix A for the details).
The right and left eigenvectors given in equations (103) and (111) form a complete biorthogonal system for the $`_0`$ subspace, i.e.
$`(\underset{d}{\overset{0}{\mathrm{\Psi }}}|\underset{d}{\overset{0}{\mathrm{\Phi }}})=1,(\underset{\omega }{\overset{0}{\mathrm{\Psi }}}|\underset{\omega ^{}}{\overset{0}{\mathrm{\Phi }}})=\delta (\omega \omega ^{}),(\underset{d}{\overset{0}{\mathrm{\Psi }}}|\underset{\omega ^{}}{\overset{0}{\mathrm{\Phi }}})=(\underset{\omega }{\overset{0}{\mathrm{\Psi }}}|\underset{d}{\overset{0}{\mathrm{\Phi }}})=0,`$
$`_0=|\underset{d}{\overset{0}{\mathrm{\Phi }}})(\underset{d}{\overset{0}{\mathrm{\Psi }}}|+{\displaystyle }d\omega |\underset{\omega }{\overset{0}{\mathrm{\Phi }}})(\underset{\omega }{\overset{0}{\mathrm{\Psi }}}|.`$
The degeneration of the zero eigenvalue of $`\underset{ext}{\overset{0}{๐}}`$ have been partially removed by the interaction: a proper eigenvalue $`\lambda _d2\pi iV_\mathrm{\Omega }^2`$ with a single eigenvector is obtained. However, the infinite degeneration of the eigenvalue $`\lambda _\omega =0`$ ($`\omega ^+`$) remains. This is a new feature that must be taken in consideration: the infinitely degenerated eigenvalue corresponds to the invariant states. As we see from equations (111), these invariant states are expanded, up to zero order, by the singular states $`(\underset{\omega }{\overset{0}{\mathrm{\Psi }}}|=(\omega |`$. These invariant states only appear if we introduce the singular structure for states and observables. The singular structure is also necessary for the description of the decay of the unstable discrete state. As we will see below, the imaginary eigenvalue $`\lambda _d`$ is associated with the exponential approximation of the decay process.
### B The continuous spectrum.
Let us consider the following generalized eigenvalues and eigenvectors of $`\underset{ext}{\overset{0}{๐}}`$
$$\underset{u1}{\overset{0}{\lambda }}=u\mathrm{\Omega },|\underset{u1}{\overset{0}{\mathrm{\Phi }}})=|u1),(\underset{u1}{\overset{0}{\mathrm{\Psi }}}|=(u1|,u\mathrm{\Gamma },$$
(112)
Equation (85) gives
$$(u\mathrm{\Omega }\underset{ext}{\overset{0}{๐}})|\stackrel{1}{\mathrm{\Phi }})=(\underset{ext}{\overset{1}{๐}}\stackrel{1}{\lambda })|u1).$$
(113)
When this equation is multiplied from the left by $`(v1|`$ we obtain
$`(uv)(v1|\stackrel{1}{\mathrm{\Phi }})=\stackrel{1}{\lambda }\delta _{\overline{\mathrm{\Gamma }}}(uv),`$
where $`\delta _{\overline{\mathrm{\Gamma }}}(uv)`$ is the โ$`\delta `$-functionโ on the curve $`\overline{\mathrm{\Gamma }}`$ ($`_{\overline{\mathrm{\Gamma }}}๐u\delta _{\overline{\mathrm{\Gamma }}}(uv)f(u)=f(v)`$). From this last equation we obtain $`\underset{u1}{\overset{1}{\lambda }}=0`$. The first order correction for the right eigenvector can be obtained from equation (113)
$$|\underset{u1}{\overset{1}{\mathrm{\Phi }}})=\frac{V_u}{u\mathrm{\Omega }}|1)_\mathrm{\Gamma }du^{}\frac{V_u^{}}{u^{}\mathrm{\Omega }}|uu^{}).$$
(114)
From equation (91) we obtain the first order correction of the left eigenvector with the same eigenvalue
$$(\underset{u1}{\overset{1}{\mathrm{\Psi }}}|=\frac{V_u}{u\mathrm{\Omega }}[(1|(u|]_\mathrm{\Gamma }du^{}\frac{V_u^{}}{u^{}\mathrm{\Omega }}(uu^{}|.$$
(115)
The second order correction to the eigenvalue can be obtained multiplying equation (87) by $`(v1|`$ from the left
$`(uv)(v1|\underset{u1}{\overset{2}{\mathrm{\Phi }}})`$ $`=`$ $`(v1|\underset{ext}{\overset{1}{๐}}|\stackrel{1}{\mathrm{\Phi }})\underset{u1}{\overset{2}{\lambda }}\delta _{\overline{\mathrm{\Gamma }}}(uv)=`$
$`=`$ $`[{\displaystyle _\mathrm{\Gamma }}du^{}{\displaystyle \frac{V_u^{}^2}{u^{}\mathrm{\Omega }}}\underset{u1}{\overset{2}{\lambda }}]\delta _{\overline{\mathrm{\Gamma }}}(uv)+{\displaystyle \frac{V_uV_v}{u\mathrm{\Omega }}}.`$
From this equation we obtain
$$\underset{u1}{\overset{2}{\lambda }}=_\mathrm{\Gamma }du^{}\frac{V_u^{}^2}{u^{}\mathrm{\Omega }}=_0^{\mathrm{}}d\omega \frac{V_\omega ^2}{\omega i0\mathrm{\Omega }}=i\pi V_\mathrm{\Omega }^2+_0^{\mathrm{}}d\omega V_\omega ^2๐ซ\left(\frac{1}{\omega \mathrm{\Omega }}\right).$$
(116)
Therefore, the generalized eigenvalues of $`\underset{ext}{\overset{0}{๐}}`$, represented by points $`z=u\mathrm{\Omega }`$ ($`u\overline{\mathrm{\Gamma }}`$) of the complex plane, are shifted to the upper half plane by the interaction. The corresponding eigenvalues of $`๐_{ext}`$ are the points $`z=u\mathrm{\Omega }+\underset{u1}{\overset{2}{\lambda }}`$, where $`u\overline{\mathrm{\Gamma }}`$ and $`\underset{u1}{\overset{2}{\lambda }}`$ is given by equation (116) (see Fig. 2).
We can also consider the following generalized eigenvalues and eigenvectors of $`\underset{ext}{\overset{0}{๐}}`$
$$\underset{1u^{}}{\overset{0}{\lambda }}=\mathrm{\Omega }u^{},|\underset{1u^{}}{\overset{0}{\mathrm{\Phi }}})=|1u^{}),(\underset{1u^{}}{\overset{0}{\mathrm{\Psi }}}|=(1u^{}|,u^{}\mathrm{\Gamma },$$
(117)
The perturbation method gives
$`\underset{1u^{}}{\overset{1}{\lambda }}`$ $`=`$ $`0,`$ (118)
$`\underset{1u^{}}{\overset{2}{\lambda }}`$ $`=`$ $`{\displaystyle _{\overline{\mathrm{\Gamma }}}}๐u{\displaystyle \frac{V_u^2}{\mathrm{\Omega }u}}=i\pi V_\mathrm{\Omega }^2{\displaystyle _0^{\mathrm{}}}๐\omega V_\omega ^2๐ซ\left({\displaystyle \frac{1}{\omega \mathrm{\Omega }}}\right),`$ (119)
$`|\underset{1u^{}}{\overset{1}{\mathrm{\Phi }}})`$ $`=`$ $`{\displaystyle \frac{V_u^{}}{u^{}\mathrm{\Omega }}}|1){\displaystyle _{\overline{\mathrm{\Gamma }}}}du{\displaystyle \frac{V_u}{u\mathrm{\Omega }}}|uu^{}),`$ (120)
$`(\underset{1u^{}}{\overset{1}{\mathrm{\Psi }}}|`$ $`=`$ $`{\displaystyle \frac{V_u^{}}{u^{}\mathrm{\Omega }}}[(1|(u^{}|]{\displaystyle _{\overline{\mathrm{\Gamma }}}}du{\displaystyle \frac{V_u}{u\mathrm{\Omega }}}(uu^{}|.`$ (121)
This result shows that the generalized eigenvalues of equation (117) are also shifted to the upper half plane by the interaction (see Fig. 2).
Finally, we can consider the following generalized eigenvalues and eigenvectors of $`\underset{ext}{\overset{0}{๐}}`$
$$\underset{uu^{}}{\overset{0}{\lambda }}=uu^{},|\underset{uu^{}}{\overset{0}{\mathrm{\Phi }}})=|uu^{}),(\underset{uu^{}}{\overset{0}{\mathrm{\Psi }}}|=(uu^{}|,u\overline{\mathrm{\Gamma }},u^{}\mathrm{\Gamma },$$
(122)
(see the shaded region in Fig.1). In this case the perturbed expansion gives
$`\underset{uu^{}}{\overset{1}{\lambda }}`$ $`=`$ $`\underset{uu^{}}{\overset{2}{\lambda }}=0,`$ (123)
$`|\underset{uu^{}}{\overset{1}{\mathrm{\Phi }}})`$ $`=`$ $`{\displaystyle \frac{V_u}{u\mathrm{\Omega }}}|1u^{})+{\displaystyle \frac{V_u^{}}{u^{}\mathrm{\Omega }}}|u1),`$ (124)
$`(\underset{uu^{}}{\overset{1}{\mathrm{\Psi }}}|`$ $`=`$ $`{\displaystyle \frac{V_u}{u\mathrm{\Omega }}}(1u^{}|+{\displaystyle \frac{V_u^{}}{u^{}\mathrm{\Omega }}}(u1|.`$ (125)
There is no shift of this part of the generalized spectrum by the interaction.
### C Time evolution.
In the previous subsection we obtained by a perturbation method a set of generalized eigenvalues and eigenvectors of the superoperator $`๐_{ext}`$, which reduce to the complete biorthogonal set of $`\underset{ext}{\overset{0}{๐}}`$ when the parameter of the interaction goes to zero. The generalized eigenvalues of $`๐_{ext}`$ are shown in Fig. 3. It seems reasonable to assume that the obtained generalized eigenvalues and eigenvectors of $`๐_{ext}`$, also form a complete biorthogonal set, at least for small values of the interaction parameter This fact can be explicitly verified up to second order for the generalized eigenvectors obtained in sections IIIA and IIIB.. If this is the case, the time evolution operator is
$`\mathrm{exp}(i๐_{ext}t)`$ $`=`$ $`{\displaystyle }d\omega \mathrm{exp}(i\lambda _\omega t)|\mathrm{\Phi }_\omega )(\mathrm{\Psi }_\omega |+\mathrm{exp}(i\lambda _dt)|\mathrm{\Phi }_d)(\mathrm{\Psi }_d|+{\displaystyle _{\overline{\mathrm{\Gamma }}}}du\mathrm{exp}(i\lambda _{u1}t)|\mathrm{\Phi }_{u1})(\mathrm{\Psi }_{u1}|+`$ (127)
$`+{\displaystyle _\mathrm{\Gamma }}du^{}\mathrm{exp}(i\lambda _{1u^{}}t)|\mathrm{\Phi }_{1u^{}})(\mathrm{\Psi }_{1u^{}}|+{\displaystyle _{\overline{\mathrm{\Gamma }}}}du{\displaystyle _\mathrm{\Gamma }}du^{}\mathrm{exp}(i\lambda _{uu^{}}t)|\mathrm{\Phi }_{uu^{}})(\mathrm{\Psi }_{uu^{}}|.`$
Up to second order, the eigenvalues are
$`\lambda _\omega `$ $`=`$ $`0,`$ (128)
$`\lambda _d`$ $`=`$ $`2\pi iV_\mathrm{\Omega }^2,`$ (129)
$`\lambda _{u1}`$ $`=`$ $`u\mathrm{\Omega }+i\pi V_\mathrm{\Omega }^2+{\displaystyle _0^{\mathrm{}}}๐\omega V_\omega ^2๐ซ\left({\displaystyle \frac{1}{\omega \mathrm{\Omega }}}\right),`$ (130)
$`\lambda _{1u^{}}`$ $`=`$ $`\mathrm{\Omega }u^{}+i\pi V_\mathrm{\Omega }^2{\displaystyle _0^{\mathrm{}}}๐\omega V_\omega ^2๐ซ\left({\displaystyle \frac{1}{\omega \mathrm{\Omega }}}\right),`$ (131)
$`\lambda _{uu^{}}`$ $`=`$ $`uu^{}`$ (132)
and the approximated expressions for the eigenvectors are given in the previous subsection.
The time evolution of the mean value of an observable $`O`$, given by
$`O_{\rho _t}=(\rho _t|O)=(\rho _0|\mathrm{exp}(i๐t)|O),`$
can be easily computed using (127). Let us consider the time evolution of the unstable state given by $`(\rho _0|=(1|`$. If we call by $`\epsilon `$ the small multiplicative interaction parameter included in $`V_\omega `$, we can use the explicit approximated expressions of the eigenvectors to prove that
$`(1|\mathrm{\Phi }_\omega )(\mathrm{\Psi }_\omega |1)`$ $`=`$ $`O(\epsilon ^2),`$
$`(1|\mathrm{\Phi }_d)(\mathrm{\Psi }_d|1)`$ $`=`$ $`1+O(\epsilon ^2),`$
$`(1|\mathrm{\Phi }_{u1})(\mathrm{\Psi }_{u1}|1)`$ $`=`$ $`O(\epsilon ^2),`$
$`(1|\mathrm{\Phi }_{1u^{}})(\mathrm{\Psi }_{1u^{}}|1)`$ $`=`$ $`O(\epsilon ^2),`$
$`(1|\mathrm{\Phi }_{uu^{}})(\mathrm{\Psi }_{uu^{}}|1)`$ $`=`$ $`O(\epsilon ^2),`$
and therefore the survival probability of the unstable state which is represented by the state functional $`(\rho _0|=(1|`$ (and also by the vector state $`|1`$) is approximately given by
$$(\rho _t|1)=(1|\mathrm{exp}(i๐t)|1)\mathrm{exp}(2\pi V_\mathrm{\Omega }^2t).$$
(133)
This expression for the survival probability is valid if $`\epsilon ^21`$ and $`t\epsilon ^2`$.
Taking into account that
$`(1|\mathrm{\Phi }_\omega ^{})(\mathrm{\Psi }_\omega ^{}|\omega )`$ $`=`$ $`\delta (\omega ^{}\mathrm{\Omega })\delta (\omega ^{}\omega )+O(\epsilon ^2),`$
$`(1|\mathrm{\Phi }_d)(\mathrm{\Psi }_d|\omega )`$ $`=`$ $`\delta (\omega \mathrm{\Omega })+O(\epsilon ^2),`$
$`(1|\mathrm{\Phi }_{u1})(\mathrm{\Psi }_{u1}|\omega )`$ $`=`$ $`O(\epsilon ^2),`$
$`(1|\mathrm{\Phi }_{1u^{}})(\mathrm{\Psi }_{1u^{}}|\omega )`$ $`=`$ $`O(\epsilon ^2),`$
$`(1|\mathrm{\Phi }_{uu^{}})(\mathrm{\Psi }_{uu^{}}|\omega )`$ $`=`$ $`O(\epsilon ^2),`$
for the initial unstable state $`(\rho _0|=(1|`$ we obtain
$$(\rho _t|\omega )=(1|\mathrm{exp}(i๐t)|\omega )\delta (\omega \mathrm{\Omega })\left[1\mathrm{exp}(2\pi V_\mathrm{\Omega }^2t)\right],$$
(134)
also valid if $`\epsilon ^21`$ and $`t\epsilon ^2`$. Equations (133)-(134) describes the transition of the unstable state $`(1|`$, into $`(\omega =\mathrm{\Omega }|`$.
Therefore, the obtained set of generalized eigenvectors is a useful tool for a complete description of the decay process. The pure state $`|1`$ is given in this formalism by the state functional $`(1|`$, and the standard approximated expression its the survival probability is reobtained in equation (133). But in addition we obtain in equation (134) an explicit expression for the by-products of the decay process.
## IV Conclusions.
When faced with the problem of the spectral decomposition of the Hamiltonian
$`H`$ $`=`$ $`\stackrel{0}{H}+\stackrel{1}{H},`$ (135)
$`\stackrel{0}{H}`$ $`=`$ $`\mathrm{\Omega }|11|+{\displaystyle _0^{\mathrm{}}}๐\omega \omega |\omega \omega |,\mathrm{\Omega }^+`$ (136)
$`\stackrel{1}{H}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}๐\omega V_\omega |\omega 1|+{\displaystyle _0^{\mathrm{}}}๐\omega \overline{V_\omega }|1\omega |+{\displaystyle _0^{\mathrm{}}}๐\omega {\displaystyle _0^{\mathrm{}}}๐\omega ^{}V_{\omega \omega ^{}}|\omega \omega ^{}|,`$ (137)
we find that the superposition of the continuous and the discrete part of the spectrum of $`\stackrel{0}{H}`$ makes impossible to apply the standard methods of perturbation theory.
If we restrict the class of states (observables) of the system to the set represented by vectors $`\mathrm{\Phi }`$ ($`\mathrm{\Psi }`$) for which $`\omega |\mathrm{\Phi }`$ ($`\omega |\mathrm{\Psi }`$) has a well defined analytic extension $`z|\mathrm{\Phi }`$ ($`z|\mathrm{\Psi }`$) to the lower (upper) complex half plane, it is possible to decompose the amplitude $`\mathrm{\Psi }|\mathrm{\Phi }`$ of the state $`|\mathrm{\Phi }`$ to be observed in the state $`|\mathrm{\Psi }`$ in the form
$$\mathrm{\Psi }|\mathrm{\Phi }=\mathrm{\Psi }|I_\mathrm{\Gamma }|\mathrm{\Phi },I_\mathrm{\Gamma }|11|+_\mathrm{\Gamma }๐z|zz|,$$
(138)
being $`\mathrm{\Gamma }`$ a curve in the lower complex half plane as shown in fig.1.
For the amplitude of the time evolved state $`|\mathrm{\Phi }_t=\mathrm{exp}(iHt)|\mathrm{\Phi }`$ to be observed in $`|\mathrm{\Psi }`$, we also have
$$\mathrm{\Psi }|\mathrm{\Phi }_t=\mathrm{\Psi }|\mathrm{exp}(iHt)|\mathrm{\Phi }=\mathrm{\Psi }|\mathrm{exp}(iH_\mathrm{\Gamma }t)|\mathrm{\Phi },$$
(139)
where
$`H_\mathrm{\Gamma }`$ $`=`$ $`\underset{\mathrm{\Gamma }}{\overset{0}{H}}+\underset{\mathrm{\Gamma }}{\overset{1}{H}},`$ (140)
$`\underset{\mathrm{\Gamma }}{\overset{0}{H}}`$ $``$ $`\mathrm{\Omega }|11|+{\displaystyle _\mathrm{\Gamma }}๐zz|zz|,`$ (141)
$`\underset{\mathrm{\Gamma }}{\overset{1}{H}}`$ $``$ $`{\displaystyle _\mathrm{\Gamma }}๐zV_z|z1|+{\displaystyle _\mathrm{\Gamma }}๐z\overline{V_{\overline{z}}}|1z|+{\displaystyle _\mathrm{\Gamma }}๐z{\displaystyle _\mathrm{\Gamma }}๐z^{}V_{zz^{}}|zz^{}|.`$ (142)
Starting from the right (left) eigenvectors $`|1`$ and $`|z`$ ($`1|`$ and $`z|`$), which form a complete biorthogonal set for the spectral decomposition of $`\underset{\mathrm{\Gamma }}{\overset{0}{H}}`$, it is possible to obtain a well defined perturbative algorithm to compute the corresponding right (left) eigenvectors $`|f_\mathrm{\Omega }`$ and $`|f_z`$ ($`\stackrel{~}{f}_\mathrm{\Omega }|`$ and $`\stackrel{~}{f}_z|`$), a complete biorthogonal set for the spectral decomposition of $`H_\mathrm{\Gamma }`$, i.e.
$`I_\mathrm{\Gamma }`$ $`=`$ $`|f_\mathrm{\Omega }\stackrel{~}{f}_\mathrm{\Omega }|+{\displaystyle _\mathrm{\Gamma }}๐u|f_u\stackrel{~}{f}_u|`$ (143)
$`H_\mathrm{\Gamma }`$ $`=`$ $`\lambda _\mathrm{\Omega }|f_\mathrm{\Omega }\stackrel{~}{f}_\mathrm{\Omega }|+{\displaystyle _\mathrm{\Gamma }}๐u\lambda _u|f_u\stackrel{~}{f}_u|,`$ (144)
$$\stackrel{~}{f}_\mathrm{\Omega }|f_\mathrm{\Omega }=1,\stackrel{~}{f}_u|f_u^{}=\delta _\mathrm{\Gamma }(uu^{}),\stackrel{~}{f}_\mathrm{\Omega }|f_u=\stackrel{~}{f}_u|f_\mathrm{\Omega }=0.$$
(145)
For the Hamiltonian $`H_\mathrm{\Gamma }`$ of the general form given in equations (142), the orthogonality and completeness relations (equations (145) and (143)) can be verified by explicit calculations up to any finite order in the perturbation. For the Friedrichs model (a special case of the Hamiltonian given in equation (137)), the solutions of the eigenvalue problem obtained by the perturbation alghorithm coincide with the exact solutions obtained by Sudarshan C.B.Chiu, V.Gorini . The same complex spectral decomposition for the Friedrichs model was also obtained by T.Petrosky, I.Prigogine and S.Tasaki , with a perturbative method using a โtime ordering ruleโ by which an small imaginary part is included to avoid small denominators, with a different sign according to the type of transition involved. In the perturbative algorithm presented in this paper, there is no need of additional rules to avoid the singularities due to resonances between discrete and continuous parts of the unperturbed spectrum. When states and observables are restricted to have suitable analytic properties, the deformation of the integrals over the continuous part of the spectrum to a curve in the complex lower half plane, eliminate the continuous-discrete resonances. After the expressions for the eigenvalues and eigenvectors are obtained up to de desired order, the complex contour of integration can be deformed back to the real axis. In this way, the โtime ordering ruleโ of reference can be deduced from the analytic extension properties of states and observables.
If the transition probability
$$P_t=\left|\mathrm{\Psi }|\mathrm{exp}(iHt)|\mathrm{\Phi }\right|^2=\left|\mathrm{exp}(i\lambda _\mathrm{\Omega }t)\mathrm{\Psi }|f_\mathrm{\Omega }\stackrel{~}{f}_\mathrm{\Omega }|\mathrm{\Phi }+_\mathrm{\Gamma }๐z\mathrm{exp}(izt)\mathrm{\Psi }|f_z\stackrel{~}{f}_z|\mathrm{\Phi }\right|^2,$$
(146)
is obtained through eigenvalues and eigenvectors computed up to n-th order, the necessary conditions for a good approximation of $`P_t`$ are $`\epsilon ^{n+1}1`$ and $`\epsilon ^{n+1}t1`$ where $`\epsilon `$ is the small interaction parameter. Then, even for small interactions, it is only for not too large times that we can obtain a good approximation of $`P_t`$.
An important advantage of the spectral decomposition in terms of generalized eigenvectors with complex eigenvalues is that it can be obtained with a well defined perturbation algorithm. The fact that this spectral decomposition is only feasible for states $`|\mathrm{\Phi }`$ and observables $`|\mathrm{\Psi }`$ with the analytic properties mentioned above is not a limitation. When we have to model a state or an observable from the finite amount of information coming from experimental data it is always possible to choose $`\omega |\mathrm{\Phi }`$ ($`\omega |\mathrm{\Psi }`$) with well defined analytic extensions to the lower (upper) complex half plane. Moreover, we are usually free to choose upper analytic extensions for the states and lower analytic extensions for the observables . For this choice, we have
$$P_t=\left|\mathrm{\Psi }|\mathrm{exp}(iHt)|\mathrm{\Phi }\right|^2=\left|\mathrm{exp}(i\overline{\lambda }_\mathrm{\Omega }t)\mathrm{\Phi }|\stackrel{~}{f}_\mathrm{\Omega }f_\mathrm{\Omega }|\mathrm{\Psi }+_{\overline{\mathrm{\Gamma }}}๐z\mathrm{exp}(izt)\mathrm{\Phi }|\stackrel{~}{f}_zf_z|\mathrm{\Psi }\right|^2,$$
(147)
where $`|\stackrel{~}{f}_\mathrm{\Omega }`$ and $`|\stackrel{~}{f}_z`$ ($`f_\mathrm{\Omega }|`$ and $`f_z|`$) are right (left) generalized eigenvectors of $`H_{\overline{\mathrm{\Gamma }}}`$, being $`\overline{\mathrm{\Gamma }}`$ a curve in the upper half of the complex plane.
For states and observables having both analytic extensions to the upper and lower half plane near the real axis, both expressions (146) and (147) can be used to obtain $`P_t`$. However, equation (147) would be a bad choice for $`t>0`$, as the positive imaginary part of the eigenvalues will give a well defined $`P_t`$ in terms of unbounded factors (the time evolution would include exponentially growing terms).
We also considered in this paper the class of observables represented by operators of the form
$$|O)=O_1|1)+_0^{\mathrm{}}d\omega O_\omega |\omega )+_0^{\mathrm{}}d\omega _0^{\mathrm{}}d\omega ^{}O_{\omega \omega ^{}}|\omega \omega ^{})+_0^{\mathrm{}}d\omega O_{\omega 1}|\omega 1)+_0^{\mathrm{}}d\omega ^{}O_{1\omega ^{}}|1\omega ^{}),$$
(148)
where the second term of the right hand side is a singular operator, and
$$|1)|11|,|\omega )|\omega \omega |,|\omega \omega ^{})|\omega \omega ^{}|,|\omega 1)|\omega 1|,|1\omega ^{})|1\omega ^{}|.$$
(149)
Following the work of references , and , we considered the states as functionals $`\rho `$ acting on observables $`O`$ to give the mean value of the observable in the state ($`O_\rho (\rho |O)`$). The states can be expressed as linear combinations of the functionals defined by
$$(1|O)=O_1,(\omega |O)=O_\omega ,(\omega \omega ^{}|O)=O_{\omega \omega ^{}},(\omega 1|O)=O_{\omega 1},(1\omega ^{}|O)=O_{1\omega ^{}}.$$
(150)
The Liouville-Von Newmann superoperators $`\stackrel{0}{๐}=[\stackrel{0}{H},]`$ and $`\stackrel{1}{๐}=[\stackrel{1}{H},]`$ can be expanded in terms of the generalized observables and states given in equations (149) and (150). For the Friedrichs model we obtain
$`\stackrel{0}{๐}`$ $`=`$ $`{\displaystyle }d\omega ^{}(\mathrm{\Omega }\omega ^{})|1\omega ^{})(1\omega ^{}|+{\displaystyle }d\omega (\omega \mathrm{\Omega })|\omega 1)(\omega 1|+{\displaystyle }d\omega {\displaystyle }d\omega ^{}(\omega \omega ^{})|\omega \omega ^{})(\omega \omega ^{}|,`$ (151)
$`\stackrel{1}{๐}`$ $`=`$ $`{\displaystyle }d\omega V_\omega [|\omega 1)|1\omega )](1|+{\displaystyle }d\omega V_\omega [|1\omega )|\omega 1)](\omega |+{\displaystyle }d\omega [V_\omega |1)+{\displaystyle }d\omega ^{}V_\omega ^{}|\omega ^{}\omega )](1\omega |+`$ (153)
$`+{\displaystyle }d\omega [V_\omega |1){\displaystyle }d\omega ^{}V_\omega ^{}|\omega \omega ^{})](\omega 1|+{\displaystyle }d\omega {\displaystyle }d\omega ^{}[V_\omega |1\omega ^{})V_\omega ^{}|\omega 1)](\omega \omega ^{}|.`$
Assuming suitable analytic properties for $`\rho `$ and $`O`$, the eigenvalue problem for $`๐=\stackrel{0}{๐}+\stackrel{1}{๐}`$ can be transformed into the eigenvalue problem for $`๐_{ext}`$, obtained from $`๐`$ replacing the integrals $`_0^{\mathrm{}}๐\omega `$ and $`_0^{\mathrm{}}๐\omega ^{}`$ by $`_{\overline{\mathrm{\Gamma }}}๐z`$ and $`_\mathrm{\Gamma }๐z^{}`$, being $`\overline{\mathrm{\Gamma }}`$ ($`\mathrm{\Gamma }`$) a curve in the upper (lower) complex half plane. Starting from the right and left generalized eigenvectors of $`\underset{ext}{\overset{0}{๐}}`$, a well defined perturbation algorithm can be used to obtain a complete biorthogonal set of generalized eigenvectors of $`๐_{ext}`$
$`๐_{ext}={\displaystyle }d\omega \lambda _\omega |\mathrm{\Phi }_\omega )(\mathrm{\Psi }_\omega |+\lambda _d|\mathrm{\Phi }_d)(\mathrm{\Psi }_d|+{\displaystyle _{\overline{\mathrm{\Gamma }}}}du\lambda _{u1}|\mathrm{\Phi }_{u1})(\mathrm{\Psi }_{u1}|+{\displaystyle _\mathrm{\Gamma }}du^{}\lambda _{1u^{}}|\mathrm{\Phi }_{1u^{}})(\mathrm{\Psi }_{1u^{}}|+{\displaystyle _{\overline{\mathrm{\Gamma }}}}du{\displaystyle _\mathrm{\Gamma }}du^{}\lambda _{uu^{}}|\mathrm{\Phi }_{uu^{}})(\mathrm{\Psi }_{uu^{}}|.`$
$`๐_{ext}={\displaystyle }d\omega |\mathrm{\Phi }_\omega )(\mathrm{\Psi }_\omega |+|\mathrm{\Phi }_d)(\mathrm{\Psi }_d|+{\displaystyle _{\overline{\mathrm{\Gamma }}}}du|\mathrm{\Phi }_{u1})(\mathrm{\Psi }_{u1}|+{\displaystyle _\mathrm{\Gamma }}du^{}|\mathrm{\Phi }_{1u^{}})(\mathrm{\Psi }_{1u^{}}|+{\displaystyle _{\overline{\mathrm{\Gamma }}}}du{\displaystyle _\mathrm{\Gamma }}du^{}|\mathrm{\Phi }_{uu^{}})(\mathrm{\Psi }_{uu^{}}|.`$
The generalized eigenvalues of $`\underset{ext}{\overset{0}{๐}}`$ and $`๐_{ext}`$ are shown in Figs. 2 and 3. The time evolution of the mean value for any observable $`O`$ of the form given in equation (148) can be obtained in terms of these generalized eigenvalues and eigenvectors. For example, it is possible to obtain the time evolution of the unstable state $`\rho _0=(1|`$. If $`\epsilon ^21`$ and $`t\epsilon ^2`$ (being $`\epsilon `$ the interaction parameter), we obtain
$`(\rho _t|=(1|\mathrm{exp}(i๐t)\mathrm{exp}(2\pi V_\mathrm{\Omega }^2t)(1|+[1\mathrm{exp}(2\pi V_\mathrm{\Omega }^2t)](\omega =\mathrm{\Omega }|.`$
Therefore, the obtained set of generalized eigenvectors is a useful tool for a complete description of the decay process. The pure state $`|1`$ is given in this formalism by the state functional $`(1|`$, and the standard approximated expression for its survival probability is reobtained. But in addition we obtain an explicit expression for the by-products of the decay process.
The generalized spectral decomposition is useful to describe the time evolution of the physical state functionals, because the contribution of the resonances is incorporated through the exponential dumping factors. However, it is important to understand that not all the generalized left eigenvectors of the Liouville-Von Newmann superoperator $`๐`$ involved in the spectral decomposition can have an independent physical meaning: Let us consider $`(\rho |=(\mathrm{\Psi }_\lambda |`$, where $`(\mathrm{\Psi }_\lambda |`$ is a left generalized eigenvector of $`๐`$ with eigenvalue $`\lambda 0`$. Therefore we have $`(\mathrm{\Psi }_\lambda |๐=\lambda (\mathrm{\Psi }_\lambda |`$. Acting with both sides of this equation on the identity operator $`I`$ we obtain
$`\lambda (\mathrm{\Psi }_\lambda |I)=(\mathrm{\Psi }_\lambda |๐I)=(\mathrm{\Psi }_\lambda |HIIH)=(\mathrm{\Psi }_\lambda |HH)=0.`$
If $`\lambda 0`$, $`(\mathrm{\Psi }_\lambda |I)=0`$ and the total probability condition $`(\rho |I)=1`$ is not verified by $`(\mathrm{\Psi }_\lambda |`$. Therefore, it can not be a physical state.
ACKNOWLEDGMENTS
This work was partially supported by Grant No. CI1-CT94-0004 of the European Community, Grant No. PID-0150 of CONICET (National Research Council of Argentina), Grant No. EX-198 of Buenos Aires University, Grant No. 12217/1 of Fundaciรณn Antorchas, and also a Grant from the Foundation pour la Recherche Foundamentale OLAM.
## A Second order eigenvalues for the discrete spectrum.
As we showed in section II, $`\underset{ext}{\overset{0}{๐}}`$ has a zero eigenvalue with infinite degeneration. This degeneration remains unchanged up to first order but it is partially removed in the second order approximation.
For the right eigenvalue problem, the second order corrections $`\stackrel{2}{\lambda }`$ to the eigenvalues is given by equation (102)
$$2\pi iV_\mathrm{\Omega }^2|1)(1|\stackrel{0}{\mathrm{\Phi }})2\pi iV_\mathrm{\Omega }^2|1)(\omega =\mathrm{\Omega }|\stackrel{0}{\mathrm{\Phi }})=\stackrel{2}{\lambda }|1)(1|\stackrel{0}{\mathrm{\Phi }})+\stackrel{2}{\lambda }d\omega |\omega )(\omega |\stackrel{0}{\mathrm{\Phi }}),$$
(A1)
which is equivalent to the set of equations
$`\stackrel{2}{\lambda }(1|\stackrel{0}{\mathrm{\Phi }})`$ $`=`$ $`2\pi iV_\mathrm{\Omega }^2(1|\stackrel{0}{\mathrm{\Phi }})2\pi iV_\mathrm{\Omega }^2(\omega =\mathrm{\Omega }|\stackrel{0}{\mathrm{\Phi }})`$ (A2)
$`\stackrel{2}{\lambda }(\omega |\stackrel{0}{\mathrm{\Phi }})`$ $`=`$ $`0,\omega ^+.`$ (A3)
If $`\stackrel{2}{\lambda }0`$, equations (A3) (A2) give $`(\omega |\stackrel{0}{\mathrm{\Phi }})=0`$ and $`\stackrel{2}{\lambda }(1|\stackrel{0}{\mathrm{\Phi }})=2\pi iV_\mathrm{\Omega }^2(1|\stackrel{0}{\mathrm{\Phi }})`$. Therefore
$$\underset{d}{\overset{2}{\lambda }}=2\pi iV_\mathrm{\Omega }^2,|\underset{d}{\overset{0}{\mathrm{\Phi }}})=|1).$$
(A4)
If $`\stackrel{2}{\lambda }=0`$, equation (A3) give no condition on $`(\omega |\stackrel{0}{\mathrm{\Phi }})`$, so we can choose $`(\omega |\stackrel{0}{\mathrm{\Phi }})=\delta (\omega \stackrel{~}{\omega })`$. Replacing in (A2) we obtain $`(1|\stackrel{0}{\mathrm{\Phi }})=\delta (\mathrm{\Omega }\stackrel{~}{\omega })`$, and therefore
$$\underset{\stackrel{~}{\omega }}{\overset{2}{\lambda }}=0,|\underset{\stackrel{~}{\omega }}{\overset{0}{\mathrm{\Phi }}})=|\stackrel{~}{\omega })+\delta (\stackrel{~}{\omega }\mathrm{\Omega })|1).$$
(A5)
For the left eigenvectors, the second order correction to the eigenvalues is given by equation (110)
$$2\pi iV_\mathrm{\Omega }^2(\stackrel{0}{\mathrm{\Psi }}|1)(1|2\pi iV_\mathrm{\Omega }^2(\stackrel{0}{\mathrm{\Psi }}|1)(\omega =\mathrm{\Omega }|=\stackrel{2}{\lambda }(\stackrel{0}{\mathrm{\Psi }}|1)(1|+\stackrel{2}{\lambda }d\omega (\stackrel{0}{\mathrm{\Psi }}|\omega )(\omega |,$$
(A6)
or equivalently
$`\stackrel{2}{\lambda }(\stackrel{0}{\mathrm{\Psi }}|1)`$ $`=`$ $`2\pi iV_\mathrm{\Omega }^2(\stackrel{0}{\mathrm{\Psi }}|1),`$ (A7)
$`\stackrel{2}{\lambda }(\stackrel{0}{\mathrm{\Psi }}|\omega )`$ $`=`$ $`2\pi iV_\mathrm{\Omega }^2(\stackrel{0}{\mathrm{\Psi }}|1)\delta (\omega \mathrm{\Omega }).`$ (A8)
If $`\stackrel{2}{\lambda }0`$, equation (A7) gives $`\underset{d}{\overset{2}{\lambda }}=2\pi iV_\mathrm{\Omega }^2`$ and no condition on $`(\stackrel{0}{\mathrm{\Psi }}|1)`$. If we choose $`(\stackrel{0}{\mathrm{\Psi }}|1)=1`$, equation (A8) implies $`(\underset{d}{\overset{0}{\mathrm{\Psi }}}|\omega )=\delta (\omega \mathrm{\Omega })`$. Therefore
$$\underset{d}{\overset{2}{\lambda }}=2\pi iV_\mathrm{\Omega }^2,(\underset{d}{\overset{0}{\mathrm{\Psi }}}|=(1|d\omega \delta (\omega \mathrm{\Omega })(\omega |=(1|(\omega =\mathrm{\Omega }|$$
(A9)
If $`\stackrel{2}{\lambda }=0`$, equation (A7) implies $`(\stackrel{0}{\mathrm{\Psi }}|1)=0`$, while equation (A8) gives no condition on $`(\stackrel{0}{\mathrm{\Psi }}|\omega )`$. We can choose
$$\underset{\omega }{\overset{2}{\lambda }}=0,(\underset{\omega }{\overset{0}{\mathrm{\Psi }}}|=(\omega |.$$
(A10)
We have obtained the eigenvalues and eigenvectors quoted in equations (103) and (111). The normalization constants have been chosen to satisfy the orthogonality conditions
$$(\underset{d}{\overset{0}{\mathrm{\Psi }}}|\underset{d}{\overset{0}{\mathrm{\Phi }}})=1,(\underset{\omega }{\overset{0}{\mathrm{\Psi }}}|\underset{\omega ^{}}{\overset{0}{\mathrm{\Phi }}})=\delta (\omega \omega ^{}),(\underset{d}{\overset{0}{\mathrm{\Psi }}}|\underset{\omega }{\overset{0}{\mathrm{\Phi }}})=(\underset{\omega }{\overset{0}{\mathrm{\Psi }}}|\underset{d}{\overset{0}{\mathrm{\Phi }}})=0,$$
(A11)
as can be easily verified using equations (A4), (A5), (A9) and (A10). It is also straighforward to verify that these eigenvectors form a complete set to expand the subspace generated by $`_0`$, i.e.
$$_0=|1)(1|+d\omega |\omega )(\omega |=|\underset{d}{\overset{0}{\mathrm{\Phi }}})(\underset{d}{\overset{0}{\mathrm{\Psi }}}|+d\omega |\underset{\omega }{\overset{0}{\mathrm{\Phi }}})(\underset{\omega }{\overset{0}{\mathrm{\Psi }}}|.$$
(A12)
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# The Hubbard Model, Spinless Fermions, and Ising Spins
## I Introduction
The Hubbard model was originally constructed to describe a metal-insulator transition for spin-dependent fermions in a simple way . This transition reflects the competition between potential (static) energy and kinetic energy. The model is defined on a lattice, where the potential energy consists of a chemical potential and an on-site repulsion of fermions with opposite spin. The kinetic energy is given by a nearest neighbor hopping. It turned out from a number of calculations that this model has a rich structure because of the complicated interplay of charge and spin degrees of freedom. For instance, mean-field calculations for a magnetic order parameter indicate para-, ferro- and antiferromagnetic states for the half-filled system . Thus, the magnetic properties of the model became a central subject of investigations in solid state physics.
The metal-insulator transition was discussed originally by Hubbard using self-consistent approximations , later in terms of a variational approach , and in the limit of an infinite dimensional lattice . Very interesting investigations were obtained from computer simulations which indicate an insulating phase at half filling for sufficiently strong fermion interaction . However, the detailed mechanism and the properties of the transition are not entirely clear.
To study the metal-insulator transition one can, in principle, start either from the metallic or from the insulating side. As the simplest approximations one could use non-interacting fermions on the metallic side or the local limit on the insulating side, where the hopping rate is zero. Unfortunately, neither of these starting points is very useful in order to understand the interacting Hubbard model: Non-interacting fermions are unstable against an arbitrarily weak interaction , and the local limit is completely degenerate with respect to the spin. Therefore, an arbitrarily weak hopping rate would lift the degeneracy leading to a new state which might be magnetically ordered . The basic idea of the present work is to start from the extreme insulating as well as from the extreme metallic state at low temperatures and to construct a perturbation theory without analyzing its magnetic order. The latter is a restriction which simplifies the calculations significantly because the spin degree of freedom can be ignored.
In this work a grand canonical ensemble is considered, where on average one fermion per site (half-filled system) is assumed. The non-interacting fermions as well as the static fermions (i.e. fermions without a hopping term) have a $`2^M`$-degeneracy ($`M`$ is the number of lattice sites) because each site can accomodate a fermion with spin-up or one with spin-down. Consequently, a perturbation theory around one of these limiting states is plagued by the degeneracies. For instance, a perturbation around the static state is a spontaneous hop of a fermion from any site to its nearest neighbor site. As a consequence, the fermion spontaneously creates a doubly occupied site and an empty site. The doubly occupied site may decay after some time again into two singly occupied sites. The resulting state is two-fold degenerate because of the possible two spin orientations. The unperturbed state can be an antiferromagnetic (Nรฉel) state. A hopping process at a time $`t_1`$ can exchange two neighboring fermions which leads to two pairs of neighboring fermions with parallel spins. At time $`t_2`$ the inverse hopping process can re-create the original antiferromagnetic state. Therefore, the two hopping processes are not independent. Moreover, the intermediate state between time $`t_1`$ and $`t_2`$ has the same energy as the antiferromagnetic state. This implies a constant interaction in time. Consequently, the linked cluster theorem cannot be applied, since it works only for independent clusters or clusters which interact with a decaying interaction . The central point of the present work is a concept which deals with this degeneracy.
In order to control this exponential degeneracy it is natural to eliminate one spin orientation. This can be achieved formally by integrating out one of the spin orientation in the functional integral representation of the Hubbard model. The result of this operation reveals an important structure of the effective spinless fermion model which is formally an expansion of the model in terms of the degeneracy: the expansion terms are not degenerate and the perturbation expansion can be applied independently to each of them. It turns out that the expansion is equivalent to the summation over the $`2^M`$ states of thermal Ising spins which are coupled to the spinless fermions. After an approximation which is applicable for the strongly metallic and the strongly insulating regime of the Hubbard model, the fermionic degrees of freedom can be integrated out. Thus, the physics is described by the Ising spins: The strongly metallic regime is characterized by a ferromagnetic Ising structure in which the fermions can freely move at low temperatures. On the contrary, the strongly insulating regime is characterized by an antiferromagnetic Ising structure which creates a gap for the fermions, in formal analogy to a Peierls instability.
The article is organized as follows: In Sect. II the Hubbard model is defined in a coherent state representation for a grand canonical ensemble of fermions. The static limit (no hopping) of the Hubbard is briefly discussed in Sect. III. Then in Sect. IV the integration over the spin-up component of the model is performed. The resulting model of spinless fermions, which has a complicated but instantaneous cluster interaction, is analyzed in Sect. IV.A. In Sect. V the Ising spin representation of the spinless fermion model is introduced and discussed. Finally, in Sect. V.A the weak-coupling limit and in Sect. V.B the weak hopping limit are studied. Appendices A, B, and C give details of the calculations.
## II The Hubbard Model
The Hubbard model describes fermions with spin $`\sigma =,`$ on a $`d`$-dimensional lattice $`\lambda `$. It is defined by the Hamiltonian
$`H[c_\sigma ^{}(r),c_\sigma (r)]=\overline{t}{\displaystyle \underset{r,r^{},\sigma }{}}c_\sigma ^{}(r)c_\sigma (r^{})+{\displaystyle \underset{r}{}}\left[\mu {\displaystyle \underset{\sigma }{}}c_\sigma ^{}(r)c_\sigma (r)+Uc_{}^{}(r)c_{}(r)c_{}^{}(r)c_{}(r)\right],`$where $`c_\sigma ^{}(r)`$, $`c_\sigma (r)`$ are fermion creation and annihilation operators, respectively. $`\overline{t}0`$ is the hopping rate. $`r,r^{}`$ means pairs of nearest neighbor sites on the lattice. $`\mu `$ is the chemical potential.
Using this Hamiltonian a grand canonical ensemble of fermions at the inverse temperature $`\beta `$ can be defined by the partition function, given in terms of a functional integral (coherent state representation) on a Grassmann algebra . For the latter the integration over a complex Grassmann field $`(\mathrm{\Psi }_\sigma (r,t),\overline{\mathrm{\Psi }}_\sigma (r,t))`$ is given as a linear mapping from a Grassmann algebra to the complex numbers. At a space-time point $`(r,t)`$ we have for integers $`k,l0`$
$`{\displaystyle [\overline{\mathrm{\Psi }}_\sigma (r,t)]^k[\mathrm{\Psi }_\sigma (r,t)]^l๐\mathrm{\Psi }_\sigma (r,t)๐\overline{\mathrm{\Psi }}_\sigma (r,t)}=\delta _{k,1}\delta _{l,1}.`$The partition function of the grand canonical ensemble of fermions then reads
$`Z={\displaystyle \mathrm{exp}(S)๐[\mathrm{\Psi },\overline{\mathrm{\Psi }}]}`$with the action
$$S=i\mathrm{\Delta }\underset{r,t}{}\frac{1}{i\mathrm{\Delta }}\overline{\mathrm{\Psi }}_\sigma (r,t)[\mathrm{\Psi }_\sigma (r,t)\mathrm{\Psi }_\sigma (r,t\mathrm{\Delta })]+i\mathrm{\Delta }\underset{t}{}\frac{1}{\mathrm{}}H[\overline{\mathrm{\Psi }}_\sigma (r,t),\mathrm{\Psi }_\sigma (r,t\mathrm{\Delta })]$$
(1)
and the product measure
$`๐[\mathrm{\Psi }_{},\mathrm{\Psi }_{}]={\displaystyle \underset{r,t,\sigma }{}}d\mathrm{\Psi }_\sigma (r,t)d\overline{\mathrm{\Psi }}_\sigma (r,t).`$The discrete time is used with $`t=\mathrm{\Delta },2\mathrm{\Delta },\mathrm{},\beta `$. $`\overline{\mathrm{\Psi }}_\sigma (r,t)`$ and $`\mathrm{\Psi }_\sigma (r,t)`$ are independent Grassmann fields which satisfy antiperiodic boundary conditions in time $`\mathrm{\Psi }_\sigma (r,\beta +\mathrm{\Delta })=\mathrm{\Psi }_\sigma (r,\mathrm{\Delta })`$ and $`\overline{\mathrm{\Psi }}_\sigma (r,\beta +\mathrm{\Delta })=\overline{\mathrm{\Psi }}_\sigma (r,\mathrm{\Delta })`$. For the subsequent calculations it is convenient to rename $`\mathrm{\Psi }_\sigma (r,t)\mathrm{\Psi }_\sigma (r,t+\mathrm{\Delta })`$ because then the Grassmann field appears with the same time in the Hamiltonian of the action (1).
## III The Local Limit
Neglecting the hopping term in the Hamiltonian (i.e., for $`\overline{t}=0`$), the integration in the partition function factorizes in space, and the corresponding expression can be evalueated as
$$Z=\underset{r,t}{}\mathrm{exp}[S_d(r,t)]๐[\mathrm{\Psi }_{},\mathrm{\Psi }_{}]=\underset{r}{}\underset{t}{}\mathrm{exp}[S_d(r,t)]๐[\mathrm{\Psi }_{},\mathrm{\Psi }_{}]=Z_1^M,$$
(2)
where
$`S_d={\displaystyle \underset{t,r}{}}\left\{\overline{\mathrm{\Psi }}_\sigma (r,t)\mathrm{\Psi }_\sigma (r,t+\mathrm{\Delta })\overline{\mu }\overline{\mathrm{\Psi }}_\sigma (r,t)\mathrm{\Psi }_\sigma (r,t)+i\mathrm{\Delta }U\overline{\mathrm{\Psi }}_{}(r,t)\mathrm{\Psi }_{}(r,t)\overline{\mathrm{\Psi }}_{}(r,t)\mathrm{\Psi }_{}(r,t)\right\}`$with $`\mathrm{}=1`$ and $`\overline{\mu }=1i\mathrm{\Delta }\mu `$. $`Z_1`$ is the partition function of the Hubbard model with one lattice site (static, local or atomic limit):
$`Z_1={\displaystyle \underset{t}{}\mathrm{exp}[S_d(r,t)]d\mu (r,t)}=1+2\overline{\mu }^{\beta /\mathrm{\Delta }}+[\overline{\mu }^2i\mathrm{\Delta }U]^{\beta /\mathrm{\Delta }}.`$Using the new parameters $`\mu ^{}=i\mathrm{\Delta }\mu `$, $`U^{}=i\mathrm{\Delta }U`$, and $`\beta ^{}=\beta /\mathrm{\Delta }`$, we can define the following weights, depending on the number of particles per site
$`w_0`$ $`=`$ $`1/Z_1`$ (3)
$`w_1`$ $`=`$ $`2(1\mu ^{})^\beta ^{}/Z_1`$ (4)
$`w_2`$ $`=`$ $`[(1\mu ^{})^2U^{}]^\beta ^{}/Z_1.`$ (5)
In the temperature formalism, where the time $`t`$ is replaced by the imaginary time through a Wick rotation, the weights $`w_0`$, $`w_1`$ and $`w_2`$ are statistical weights. Then the average number of particles per site is
$`n=w_1+2w_2.`$At zero temperature ($`\beta ^{}\mathrm{}`$) this gives
$`n=\{\begin{array}{cc}0\hfill & \text{if }0<\mu ^{}<1+\sqrt{1+U^{}}\hfill \\ 1\hfill & \text{if }\mu ^{}<0\text{ and }\mu ^{}(\mu ^{}1)<U^{}\hfill \\ 2\hfill & \text{otherwise}\hfill \end{array}.`$The groundstate of the expansion around the local limit is degenerated because each singly occupied site can accomodate a fermion with either spin $``$ or with spin $``$. This degeneracy must be handled with care. In particular, to obtain a unique hopping expansion one has to separate the degenerate contributions. Therefore, one cannot directly work with the hopping term as a perturbation but have to set up a perturbation theory which works with non-degenerate groundstates. This means that one has to divide the degenerate groundstate such that the perturbations remain in their corresponding non-degenerate groundstates.
## IV Integration over the spin-up field $`\mathrm{\Psi }_{}`$
The action $`S`$ can be divided into three pieces as
$`S=S_{}+S_{}+S_I`$with
$`S_\sigma ={\displaystyle \underset{t}{}}\left\{{\displaystyle \underset{r}{}}\left[\overline{\mathrm{\Psi }}_\sigma (r,t)\mathrm{\Psi }_\sigma (r,t+\mathrm{\Delta })\overline{\mu }\overline{\mathrm{\Psi }}_\sigma (r,t)\mathrm{\Psi }_\sigma (r,t)\right]\tau {\displaystyle \underset{r,r^{}}{}}\overline{\mathrm{\Psi }}_\sigma (r,t)\mathrm{\Psi }_\sigma (r^{},t)\right\}`$for $`\sigma =,`$ with $`\tau =i\mathrm{\Delta }\overline{t}`$. The interaction between the two spin orientations is given by
$`S_I=U^{}{\displaystyle \underset{r,t}{}}\overline{\mathrm{\Psi }}_{}(r,t)\mathrm{\Psi }_{}(r,t)\overline{\mathrm{\Psi }}_{}(r,t)\mathrm{\Psi }_{}(r,t).`$Now it is possible to integrate out the spin-up field $`\mathrm{\Psi }_{}`$, since the field $`\mathrm{\Psi }_{}`$ appears in $`S`$ only as a quadratic form. The integration over this Grassmann field gives a determinant
$$e^{S_{}S_I}\underset{r,t}{}d\mathrm{\Psi }_{}(r,t)d\overline{\mathrm{\Psi }}_{}(r,t)=det\left[_t+\overline{\mu }+\widehat{t}U^{}\overline{\mathrm{\Psi }}_{}\mathrm{\Psi }_{}\right],$$
(6)
where $`_t`$ is the time-shift operator
$`_t\mathrm{\Psi }(r,t)=\{\begin{array}{cc}\mathrm{\Psi }(r,t+\mathrm{\Delta })\hfill & \mathrm{\Delta }t<\beta \hfill \\ \mathrm{\Psi }(r,\mathrm{\Delta })\hfill & t=\beta \hfill \end{array}.`$The last equation is due to the antiperiodic boundary condition of the Grassmann field. This definition gives
$`(_t)^1=_t^T,det(_t)=1.`$We assume that the number of time slices $`\beta ^{}`$ is even such that $`det(_t)=det(_t)=1`$. The matrix $`\widehat{t}_{r,r^{}}=\tau `$ if $`r,r^{}`$ are nearest neighbors and zero otherwise. Expressions in the determinant which do not have a specified matrix structure are implicitly multiplied by the corresponding unit matrix. For instance, $`\overline{\mu }`$ is multiplied by the space-time unit matrix whereas $`\widehat{t}`$ is multiplied by the time-like unit matrix.
In the following subsets of the space-time lattice $`\mathrm{\Lambda }=\lambda \{\mathrm{\Delta },2\mathrm{\Delta },\mathrm{},\beta \}`$ will be considered. For a subset $`\mathrm{\Lambda }_k\mathrm{\Lambda }`$ we define the determinant of the the projected matrix $`P_kAP_k`$ as
$`det_{\mathrm{\Lambda }_k}Adet_{\mathrm{\Lambda }_k}(P_kAP_k),`$where $`P_k`$ is the projector onto $`\mathrm{\Lambda }_k`$.
### A Effective Cluster Action of Spinless Fermions
The partition function is now a functional integral of the spin-down Grassmann field
$`Z={\displaystyle e^S_{}๐et\left[_t+\overline{\mu }+\widehat{t}U^{}\overline{\mathrm{\Psi }}_{}\mathrm{\Psi }_{}\right]๐[\mathrm{\Psi }_{}]}.`$Formally, the determinant could be expressed as part of the action by using the identity $`detA=\mathrm{exp}[Tr(\mathrm{log}A)]`$. However, this would be too naive because the term
$`Tr\left[\mathrm{log}\left(_t+\overline{\mu }+\widehat{t}U^{}\overline{\mathrm{\Psi }}_{}\mathrm{\Psi }_{}\right)\right]`$has a complicated interaction of the Grassmann field in space and time. Moreover, at least for $`\overline{t}=0`$ the interaction has a long-range part in time which reflects the degeneracy of the unperturbed system. Fortunately, there is a way to avoid these difficulties: As shown in Appendix A, the determinant can be expanded in terms of the partitions $`\mathrm{\Lambda }_k\mathrm{\Lambda }`$ of the space-time lattice $`\mathrm{\Lambda }`$ as
$`det\left[_t+\overline{\mu }+\widehat{t}U^{}\overline{\mathrm{\Psi }}_{}\mathrm{\Psi }_{}\right]={\displaystyle \underset{\mathrm{\Lambda }_k\mathrm{\Lambda }}{}}det_{\mathrm{\Lambda }_k}[(\overline{\mu }+\widehat{t})_t^T]\mathrm{exp}\left[Tr_{\mathrm{\Lambda }_k}\mathrm{log}\left(\mathrm{๐}(\widehat{t}+\overline{\mu })^1U^{}\overline{\mathrm{\Psi }}_{}\mathrm{\Psi }_{}\right)\right].`$The partitions include the empty set which gives $`det_{\mathrm{}}A=1`$. This expansion is the most important step for the treatment of the Hubbard model in this work. The first consequence is that the partition function $`Z`$ is now given by an expansion in terms of $`\mathrm{\Lambda }_k`$ as $`Z=_{\mathrm{\Lambda }_k}Z_{\mathrm{\Lambda }_k}`$ with
$$Z_{\mathrm{\Lambda }_k}=det_{\mathrm{\Lambda }_k}[(\overline{\mu }+\widehat{t})_t^T]\mathrm{exp}(S_{\mathrm{\Lambda }_k})๐[\mathrm{\Psi }_{}]$$
(7)
and the action
$`S_{\mathrm{\Lambda }_k}=\overline{\mathrm{\Psi }}_{}(_t\overline{\mu }\widehat{t})\mathrm{\Psi }_{}Tr_{\mathrm{\Lambda }_k}\mathrm{log}\left(\mathrm{๐}(\widehat{t}+\overline{\mu })^1U^{}\overline{\mathrm{\Psi }}_{}\mathrm{\Psi }_{}\right).`$The second term of $`S_{\mathrm{\Lambda }_k}`$ can be expanded in powers of the Grassmann field. This yields an instantaneous cluster interaction
$`S_i={\displaystyle \underset{l1}{}}{\displaystyle \frac{U_{}^{}{}_{}{}^{l}}{l}}{\displaystyle \underset{(r,t)\mathrm{\Lambda }_k}{}}{\displaystyle \underset{r_1,\mathrm{},r_{l1}}{}}(\widehat{t}+\overline{\mu })_{r,r_1}^1\overline{\mathrm{\Psi }}_{}(r_1,t)\mathrm{\Psi }_{}(r_1,t)\mathrm{}(\widehat{t}+\overline{\mu })_{r_{l1},r}^1\overline{\mathrm{\Psi }}_{}(r,t)\mathrm{\Psi }_{}(r,t).`$Due to the identity
$$(\widehat{t}+\overline{\mu })^1=\frac{1}{\overline{\mu }}\frac{\widehat{t}}{\overline{\mu }}(\widehat{t}+\overline{\mu })^1$$
the action $`S_i`$ can also be written with $`u=U^{}(\widehat{t}+\overline{\mu })_{rr}^1`$ as
$`S_i=u{\displaystyle \underset{(r,t)\mathrm{\Lambda }_k}{}}\overline{\mathrm{\Psi }}_{}(r,t)\mathrm{\Psi }_{}(r,t)`$
$`+{\displaystyle \underset{l2}{}}{\displaystyle \frac{U_{}^{}{}_{}{}^{l}}{l}}{\displaystyle \underset{(r,t)\mathrm{\Lambda }_k}{}}{\displaystyle \underset{r_1,\mathrm{},r_{l1}}{}}\left({\displaystyle \frac{\widehat{t}}{\overline{\mu }}}(\widehat{t}+\overline{\mu })^1\right)_{r,r_1}\overline{\mathrm{\Psi }}_{}(r_1,t)\mathrm{\Psi }_{}(r_1,t)`$
$`\mathrm{}\left({\displaystyle \frac{\widehat{t}}{\overline{\mu }}}(\widehat{t}+\overline{\mu })^1\right)_{r_{l1},r}\overline{\mathrm{\Psi }}_{}(r,t)\mathrm{\Psi }_{}(r,t).`$This shows that the interaction terms with $`l2`$ do not distinguish between weak interaction ($`U^{}0`$) and weak hopping ($`\overline{t}0`$) because the expansion parameter is $`U^{}\overline{t}`$.
The non-interacting limit $`U^{}=0`$ as well as the local limit $`\widehat{t}=0`$ can be checked immediately because of a vanishing interaction. The former gives the product of two determinants (from both spin orientations)
$`Z_0=det\left(_t+\overline{\mu }+\widehat{t}\right)^2`$and the latter
$`Z=Z_1^M,`$in agreement with Eq. (2).
## V Ising Spin Statistics
It is interesting to notice that the determinant $`det_{\mathrm{\Lambda }_k}[(\overline{\mu }+\widehat{t})_t^T]`$ in $`Z_{\mathrm{\Lambda }_k}`$ determines which partition $`\mathrm{\Lambda }_k`$ contributes with non-zero weight. Since $`_t`$ is diagonal in space but off-diagonal in time, the determinant $`det_{\mathrm{\Lambda }_k}(_t^T)`$ is non-zero only if a site $`r`$ is in $`\mathrm{\Lambda }_k`$ at all times. Therefore, the contributing partitions are of the form $`\mathrm{\Lambda }_k=\lambda _k\times \{\mathrm{\Delta },2\mathrm{\Delta },\mathrm{},\beta \}`$, where $`\lambda _k`$ is a partition of the space lattice $`\lambda `$. The determinant then reads
$`det_{\mathrm{\Lambda }_k}[(\overline{\mu }+\widehat{t})_t^T]=[det_{\lambda _k}(\overline{\mu }+\widehat{t})]^\beta ^{}.`$Now it is convenient to introduce the projector $`I_{\lambda _k}`$ which projects on $`\lambda _k`$
$`(I_{\lambda _k})_r=\{\begin{array}{cc}1\hfill & \text{if }r\lambda _k\hfill \\ 0\hfill & \text{if }r\lambda _k\hfill \end{array}.`$Consequently, the determinant can be expressed as
$$[det_{\lambda _k}(\overline{\mu }+\widehat{t})]^\beta ^{}=\left[det_\lambda (\mathrm{๐}+I_{\lambda _k}[\overline{\mu }+\widehat{t}\mathrm{๐}]I_{\lambda _k})\right]^\beta ^{}.$$
(8)
The interaction can also be written in terms of the projection operator:
$`S_i={\displaystyle \underset{l1}{}}{\displaystyle \frac{U_{}^{}{}_{}{}^{l}}{l}}{\displaystyle \underset{t}{}}{\displaystyle \underset{r,r_1,\mathrm{},r_{l1}}{}}(I_{\lambda _k})_r(\widehat{t}+\overline{\mu })_{r,r_1}^1\overline{\mathrm{\Psi }}_{}(r_1,t)\mathrm{\Psi }_{}(r_1,t)\mathrm{}(\widehat{t}+\overline{\mu })_{r_{l1},r}^1\overline{\mathrm{\Psi }}_{}(r,t)\mathrm{\Psi }_{}(r,t)`$The partitions $`\lambda _k`$ enter $`Z`$ only through the projection operator $`I_{\lambda _k}`$. Eqs. (7) and (8) imply
$`Z={\displaystyle \underset{\lambda _k}{}}Z_{\lambda _k}={\displaystyle \underset{\lambda _k}{}}\left[det_\lambda (\mathrm{๐}+I_{\lambda _k}[\overline{\mu }+\widehat{t}\mathrm{๐}]I_{\lambda _k})\right]^\beta ^{}{\displaystyle \mathrm{exp}(S_{\mathrm{\Lambda }_k})๐[\mathrm{\Psi }_{}]}`$with the action
$`S_{\lambda _k}=\overline{\mathrm{\Psi }}_{}(_t\widehat{t}\overline{\mu }+uI_{\lambda _k})\mathrm{\Psi }_{}+O((U^{}\overline{t}/\overline{\mu }^2)^2).`$The projection operator $`I_{\lambda _k}`$ can be expressed as $`(I(\{S_r\}))_{rr^{}}=(1+S_r)\delta _{rr^{}}/2`$ with the Ising spin
$`S_r=\{\begin{array}{cc}1\hfill & r\lambda _k\hfill \\ 1\hfill & r\lambda _k\hfill \end{array}.`$The partitions $`\lambda _k`$ enter $`Z`$ only through the projection operator $`I_{\lambda _k}`$. Therefore, the sum over the randomly chosen partitions $`\lambda _k`$ of the space lattice is equivalent to the sum over randomly chosen Ising spins. This is a thermal Ising model with complicated interaction spin interaction in a magnetic field. The latter is implied by the missing invariance under a global transformation $`S_rS_r`$.
Neglecting the interaction $`S_i`$, which is weak for the metallic as well as for the insulating regime, one gets
$$ZZ_I=\underset{\{S_r=\pm 1\}}{}\left[det_\lambda (\mathrm{๐}+I[\overline{\mu }+\widehat{t}\mathrm{๐}]I)\right]^\beta ^{}det(_t+\overline{\mu }uI+\widehat{t}).$$
(9)
The partition function $`Z_I`$ can be used as a starting point for an approximative treatment of the Hubbard model. In the rest of this work the favoured spin structures will be analyzed, i.e., the spin configurations with maximal Boltzmann weight. Details of the calculations are given in App. B and C.
### A Weak Interaction: $`U^{}0`$
A metallic behavior is expected in this regime. Since
$`[det_{\lambda _k}(\overline{\mu }+\widehat{t})]^\beta ^{}`$is independent of $`U^{}`$, only the second determinant of $`Z_I`$ must be expanded around $`U^{}=0`$. This yields in leading order
$`det(_t+\overline{\mu }uI+\widehat{t})\mathrm{exp}[{\displaystyle \mathrm{\Theta }(|\kappa |1)\mathrm{log}(\kappa )d^dk}]\mathrm{exp}[\beta ^{}h{\displaystyle \underset{r}{}}(1+S_r)].`$with the effective magnetic field
$`h={\displaystyle \frac{u}{2}}{\displaystyle \frac{\mathrm{\Theta }(|\kappa |1)}{\kappa }d^dk}\mathrm{with}\kappa =\overline{\mu }+2\widehat{t}{\displaystyle \underset{j=1}{\overset{d}{}}}\mathrm{cos}k_j.`$Thus, the weakly interacting regime prefers a ferromagnetic state in terms of the Ising spin $`S_r=1`$. Knowing the Ising spin configuration with the highest weight, one can return to the partition function of Eq. (9) and obtains
$`Z_I[det(\overline{\mu }+\widehat{t})]^\beta ^{}det(_t+\overline{\mu }u+\widehat{t}).`$The argument of the determinant can be considered as an inverse Greenโs function
$`G=(_t+\overline{\mu }u+\widehat{t})^1.`$Thus, the effect of the ferromagnetic Ising spin structure is a constant shift of the chemical potential. Of course, the fluctuations of the Ising spins around the structure with maximal weight have to be taken into account as soon as $`\beta ^{}<\mathrm{}`$.
### B Hopping Expansion: $`\overline{t}0`$
The expansion of the first determinant of $`Z_I`$ in powers of $`\tau `$ is
$$\overline{\mu }^{\beta ^{}|\lambda _k|}\mathrm{exp}[\frac{\beta ^{}}{2\overline{\mu }^2}\underset{r,r\lambda _k}{}(\widehat{t}_{r,r^{}})^2]$$
(10)
and the expansion of the second determinant is
$$det(_t+\overline{\mu }uI+\widehat{t})\overline{\mu }^{\beta ^{}(M|\lambda |)}\mathrm{exp}\left[\frac{\beta ^{}\tau ^2}{2}\underset{r,r}{}W(V_r,V_r^{})\right].$$
(11)
Here the potential terms $`V_r,V_r^{}`$ are either $`\overline{\mu }`$ or $`\overline{\mu }u\overline{\mu }U^{}/\overline{\mu }`$ with the function
$`W(V_r,V_r^{})=\{\begin{array}{cc}0\hfill & \text{for }V_r=V_r^{}=\overline{\mu }u\hfill \\ 1/(u\overline{\mu })\hfill & \text{for }V_rV_r^{}\hfill \\ 1/\overline{\mu }^2\hfill & \text{for }V_r=V_r^{}=\overline{\mu }\hfill \end{array}.`$For $`\tau =0`$ all terms in the sum $`_\lambda Z_\lambda `$ have the same weight. This reflects the fact that there is a degenerate perturbation theory for the hopping expansion in the half-filled system due to the spin degree of freedom. For $`\tau 0`$ the second factor of (11) can be approximated for $`\beta \mathrm{}`$ ($`\beta `$ is now real) as
$$\mathrm{exp}\left[\frac{\beta ^{}\tau ^2}{2}\underset{r,r}{}W(V_r,V_r^{})\right]\mathrm{exp}\left[\beta \frac{\overline{t}^2}{2U}\underset{r,r}{}S_rS_r^{}\right].$$
(12)
According to (10) and (12) the maximal contribution to $`Z_I`$ comes from $`V_rV_r^{}`$ on nearest-neighbor sites. This corresponds with an antiferromagnet state of the Ising spin $`S_r=(1)^{r_1+\mathrm{}+r_d}`$. The partition function in Eq. (9) reads with the staggered Ising spin structure
$`Z_I2\overline{\mu }^{\beta ^{}M/2}det(_t+V_{AFM}+\widehat{t}),`$where $`V_{AFM}`$ is the staggered antiferromagnetic potential
$`V_{AFM}=\overline{\mu }{\displaystyle \frac{u}{2}}[1+(1)^{r_1+\mathrm{}+r_d}].`$As an effective Greenโs function of the fermions one can study
$`G=(_t+V_{AFM}+\widehat{t})^1.`$The staggered potential creates a gap. This can be seen, for instance, in the special case of a two-dimensional lattice. Then the eigenvalues of $`G^1`$ are
$`\overline{\mu }{\displaystyle \frac{u}{2}}\pm \sqrt{{\displaystyle \frac{u^2}{4}}+|\tau h_{12}|^2}`$with $`h_{12}=1+e^{ik_x}+e^{ik_y}+e^{ik_x+ik_y}`$ ($`\pi k_j<\pi `$). This result describes the physics of an insulator with a gap
$`U/2,`$in contrast to the homogeneous potential of the ferromagnetic Ising structure for the weak interacting case.
There can be a phase transition from weak coupling (ferromagnetic Ising system) to weak hopping (antiferromagnetic Ising system) even at finite $`\beta `$ (i.e., non-zero temperature). It is not clear, however, whether this is a first or second order transition. Numerical (Monte Carlo) simulations may reveal more details of the transition.
## VI Conclusions
Considering quantities of the Hubbard model which depend only on one spin direction, e.g., spin-down, it is possible to integrate out the other spin direction in the functional integral. This idea was carried out for the partition function $`Z`$. The latter could have been extended to a generating function for spin-down Greenโs function without modification of the procedure. The result of this integration is a functional integral which has a representation in terms of thermal Ising spins. Moreover, the fermionic interaction of the effective model has a coupling parameter $`U^{}\overline{t}`$ which becomes small in the limit of weak Hubbard interaction $`U^{}`$ as well as in the limit of low mobility of the particle (small hopping rate $`\overline{t}`$). Therefore, the effective fermionic interaction can be neglected if one is only interested in the strongly metallic and the strongly insulating regime at half filling. The advantage of this approach is that the degeneracy of the special cases $`U^{}=0`$ and $`\overline{t}=0`$, which are difficult to handle, are controlled by the coupling of $`U^{}`$ and $`\overline{t}`$ to the Ising spins: a small Hubbard interaction $`U^{}`$ creates an effective magnetic field $`hU^{}`$ which couples linearly to the Ising spin. On the other hand, a small hopping term favours an antiferromagnetic (staggered) Ising spin configuration because of the effective Ising spin interaction
$`{\displaystyle \frac{\overline{t}^2}{2U}}{\displaystyle \underset{r,r}{}}S_rS_r^{}.`$The Boltzmann weights of the effective Ising spin system were used to extract the most relevant contributions to the sum over the spin configurations.
The introduction of Ising spins by a direct decoupling of the interaction term in Eq. (1) (Hubbard-Stratonovich transformation) is a well-known approach. It leads to a dynamic Ising spin . The approach of the present work avoids the dynamics of the Ising spin. However, the price for this simplification is a cluster interaction of spinless fermions.
Although the approximations used in this work do not provide insight into the transition from the metallic to the insulating phase, a more accurate treatment of the partition function $`Z_I`$ of (9) by including fluctuations of the Ising spins, e.g., using a Monte Carlo simulation, may allow to access the transition at not too low temperatures.
Acknowledgement
This work was supported by the Sonderforschungsbereich 484.
## Appendix A
The space-time determinant on the r.h.s. of Eq. (6) can also be written as
$`det\left[_t+\overline{\mu }+\widehat{t}U^{}\overline{\mathrm{\Psi }}_{}\mathrm{\Psi }_{}\right]=det(\mathrm{๐}+A){\displaystyle \underset{\pi }{}}(1)^\pi {\displaystyle \underset{(r,t)\mathrm{\Lambda }}{}}\left[\delta _{\pi (r,t),(r,t)}+A_{r,t;\pi (r,t)}\right]`$with the matrix $`A=(\overline{\mu }+\widehat{t}U^{}\overline{\mathrm{\Psi }}_{}\mathrm{\Psi }_{})_t^T`$. The product over the lattice sites gives a sum over all subsets $`\mathrm{\Lambda }_k\mathrm{\Lambda }`$ of the space $`\mathrm{\Lambda }`$ and their complements $`\mathrm{\Lambda }_k^{}=\mathrm{\Lambda }\mathrm{\Lambda }_k`$
$`{\displaystyle \underset{\mathrm{\Lambda }_k\mathrm{\Lambda }}{}}{\displaystyle \underset{\pi }{}}(1)^\pi \left[{\displaystyle \underset{(r,t)\mathrm{\Lambda }_k}{}}A_{r,t;\pi (r,t)}\right]\left[{\displaystyle \underset{(r,t)\mathrm{\Lambda }_k^{}}{}}\delta _{\pi (r,t),(r,t)}\right].`$The Kronecker delta $`\delta _{\pi (r,t),(r,t)}`$ on $`\mathrm{\Lambda }_k^{}`$ implies $`\pi (r,t)\mathrm{\Lambda }_k`$ for $`(r,t)\mathrm{\Lambda }_k`$. Therefore, only that part of the matrix $`A`$ contributes which is projected onto $`\mathrm{\Lambda }_k`$. This implies an expansion of the determinant in terms of all partitions of the space-time lattice $`\mathrm{\Lambda }`$ as
$$det\left[_t+\overline{\mu }+\widehat{t}U^{}\overline{\mathrm{\Psi }}_{}\mathrm{\Psi }_{}\right]=\underset{\mathrm{\Lambda }_k\mathrm{\Lambda }}{}det_{\mathrm{\Lambda }_k}(P_kAP_k)\underset{\mathrm{\Lambda }_k\mathrm{\Lambda }}{}det_{\mathrm{\Lambda }_k}A$$
(13)
with $`det_{\mathrm{}}A=1`$ for an empty set $`\mathrm{\Lambda }_k`$. The projected determinant is
$`det_{\mathrm{\Lambda }_k}((\overline{\mu }+\widehat{t}U^{}\overline{\mathrm{\Psi }}_{}\mathrm{\Psi }_{})_t^T)=det_{\mathrm{\Lambda }_k}[(\overline{\mu }+\widehat{t})_t^T]det_{\mathrm{\Lambda }_k}(\mathrm{๐}(\overline{\mu }+\widehat{t})^1U^{}\overline{\mathrm{\Psi }}_{}\mathrm{\Psi }_{})),`$where the second second determinant reads
$`\mathrm{exp}\left[Tr_{\mathrm{\Lambda }_k}\mathrm{log}\left(\mathrm{๐}(\widehat{t}+\overline{\mu })^1U^{}\overline{\mathrm{\Psi }}_{}\mathrm{\Psi }_{}\right)\right].`$
## Appendix B
Expanding the second determinant of $`Z_I`$ up to first order in $`\delta \widehat{V}`$ yields
$`det(_t+\overline{\mu }+\delta \widehat{V}+\widehat{t})det(_t+\overline{\mu }+\widehat{t})\mathrm{exp}\left(Tr[\delta \widehat{V}(_t+\overline{\mu }+\widehat{t})^1]\right).`$The trace term can be written as
$`\beta ^{}h{\displaystyle \underset{r}{}}(1+S_r),`$where
$`h={\displaystyle \frac{u}{2}}{\displaystyle \frac{\mathrm{\Theta }(|\kappa |1)}{\kappa }\frac{d^dk}{(2\pi )^d}}\mathrm{with}\kappa =\overline{\mu }+2\tau {\displaystyle \underset{j=1}{\overset{d}{}}}\mathrm{cos}k_j.`$
## Appendix C
The first determinant of $`Z_I`$ can be expanded up to second order in $`\widehat{t}`$ as
$`[det_{\lambda _k}(\overline{\mu }+\widehat{t})]^\beta ^{}\overline{\mu }^{\beta ^{}|\lambda _k|}\mathrm{exp}[{\displaystyle \frac{\beta ^{}}{2\overline{\mu }^2}}{\displaystyle \underset{r,r^{}\lambda _k}{}}(\widehat{t}_{rr^{}})^2].`$The second determinant of $`Z_I`$ can be expanded in powers of $`\widehat{t}`$ up to second order
$`det(_t+\overline{\mu }+\delta \widehat{V}+\widehat{t})`$
$`det(_t+\overline{\mu }+\delta \widehat{V})\mathrm{exp}\left({\displaystyle \frac{1}{2}}Tr[\widehat{t}(_t+\overline{\mu }+\delta \widehat{V})^1\widehat{t}(_t+\overline{\mu }+\delta \widehat{V})^1]\right).`$The trace term reads
$`Tr[\widehat{t}(_t+\overline{\mu })^1\widehat{t}(_t+\overline{\mu })^1]=\beta ^{}{\displaystyle \underset{r,r^{}}{}}(\widehat{t}_{rr^{}})^2{\displaystyle _0^{2\pi }}{\displaystyle \frac{1}{V_re^{i\omega }}}{\displaystyle \frac{1}{V_r^{}e^{i\omega }}}{\displaystyle \frac{d\omega }{2\pi }},`$where the integral gives
$`{\displaystyle _0^{2\pi }}{\displaystyle \frac{1}{V_re^{i\omega }}}{\displaystyle \frac{1}{V_r^{}e^{i\omega }}}{\displaystyle \frac{d\omega }{2\pi }}=\{\begin{array}{cc}0\hfill & \text{for }V_r=V_r^{}=\overline{\mu }U^{}/\overline{\mu }\hfill \\ 1/(u\overline{\mu })\hfill & \text{for }V_rV_r^{}\hfill \\ 1/\overline{\mu }^2\hfill & \text{for }V_r=V_r^{}=\overline{\mu }\hfill \end{array}.`$Due to $`\tau 0`$, the approximation $`uU^{}/\overline{\mu }`$ can be used. Then for $`\beta ^{}\mathrm{}`$ the determinant is
$`det(_t+\overline{\mu }+\delta \widehat{V})={\displaystyle \underset{r}{}}(1+[\overline{\mu }+\delta \widehat{V}_r]^\beta ^{})\overline{\mu }^{\beta ^{}(M|\lambda _k|)},`$since $`|\overline{\mu }|>1`$ and $`|\overline{\mu }U^{}/\overline{\mu }|<1`$.
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# Dynamical Charge Susceptibility of the Spinless Falicov-Kimball Model.
## I Introduction
The Falicov-Kimball model was introduced in 1969 as a thermodynamic model for metal-insulator transitions in systems that have two different types of electronsโitinerant conduction electrons and localized $`f`$-electrons. Twenty years later, Brandt and Mielsch solved the Falicov-Kimball model exactly in the limit of infinite dimensions by using dynamical mean-field theory and the Baym-Kadanoff formalism. They also showed how to calculate static susceptibilities and found an Ising-like phase transition to a chessboard charge-density-wave phase at half-filling. Freericks then showed that the system also displayed incommensurate order and phase separation at other fillings.
The Hamiltonian of the spinless Falicov-Kimball model is
$$\widehat{H}=\frac{t^{}}{2\sqrt{d}}\underset{i,j}{}d_i^{}d_j+ฯต_f\underset{i}{}w_i\mu \underset{i}{}(d_i^{}d_i+w_i)+U\underset{i}{}d_i^{}d_iw_i,$$
(1)
where $`d_i^{}`$ $`(d_i)`$ is the spinless conduction electron creation (annihilation) operator at lattice site $`i`$ and $`w_i=0`$ or 1 is a classical variable corresponding to the localized $`f`$-electron number at site $`i`$. The hopping matrix between nearest neighbors $`i,j`$ (on a hypercubic lattice in $`d`$-dimensions, with $`d\mathrm{}`$) is $`t^{}/(2\sqrt{d})`$ with $`t^{}`$ chosen as our energy unit, $`ฯต_f`$ is the localized electron level, $`\mu `$ is the chemical potential and $`U`$ is the mutual electron repulsion when a conduction electron and a localized $`f`$-electron both occupy the same lattice site.
In this contribution we examine the $`d`$-electron dynamical charge susceptibility \[where $`n_i(\tau )=\mathrm{exp}(\tau \widehat{H})n_i(0)\mathrm{exp}(\tau \widehat{H})`$ and $`n_i=d_i^{}d_i`$\]
$$\chi (๐ช,i\nu _l)=\mathrm{Tr}T_\tau \underset{j,k}{}_0^\beta d\tau e^{i\nu _l\tau }e^{i๐ช(๐_๐ฃ๐_๐ค)}\left[\frac{e^{\beta \widehat{H}}n_j(\tau )n_k(0)}{Z}\frac{e^{\beta \widehat{H}}n_j}{Z}\frac{e^{\beta \widehat{H}}n_k}{Z}\right],$$
(2)
and its analytic continuation to the real-frequency axis. Here $`i\nu _l=2i\pi Tl`$ is the Bosonic Matsubara frequency, $`Z`$ is the partition function, and $`๐_๐ฃ`$ is the position operator for lattice site $`j`$. The Falicov-Kimball model is the simplest many-body problem that has nontrivial dynamics. The model is simple enough that an exact solution can be found for the charge susceptibility, but is complicated enough to show many-body effects. Shvaika has shown how to determine the lattice charge susceptibility from a diagrammatic analysis of the atomic problem. Here we provide an alternate derivation, stressing the Baym-Kadanoff approach (which produces a much simpler form for the vertex), and we provide numerical results to complement the formalism.
We begin with a discussion of some special symmetries of the Falicov-Kimball model. Since the total conduction electron number is a conserved quantity, we find that $`_jn_j(\tau )=_jn_j(0)`$ has no $`\tau `$-dependence. Hence, the uniform charge susceptibility (q=0) vanishes for all nonzero frequencies. Similarly, the local $`f`$-occupation is also conserved \[$`w_i(\tau )=w_i(0)`$\], since $`w_i`$ commutes with the lattice Hamiltonian. This implies that both the $`ff`$-charge susceptibility and the mixed $`df`$-charge susceptibility have no $`\tau `$-dependence and therefore are static, with no frequency dependence. Furthermore, a simple application of the chain rule and the definition of each of the susceptibilities as a derivative of the corresponding average electron filling with respect to an external field, shows that there is a discontinuity in the $`dd`$-charge susceptibility at zero frequency due to the coupling of the $`d`$ and $`f`$-electrons. This implies that there is a decoupling between the static and dynamic $`dd`$-charge susceptibility!
The solution of the Falicov-Kimball model has been outlined in detail elsewhere. Here we summarize the main points to establish our notation. The local Greenโs function at the Fermionic Matsubara frequency $`i\omega _n=i\pi T(2n+1)`$ is defined by
$$G_n=G(i\omega _n)=\mathrm{Tr}_\mathrm{d}\mathrm{Tr}_\mathrm{f}T_\tau _0^\beta d\tau e^{i\omega _n\tau }\frac{e^{\beta \widehat{H}_{at}}d(\tau )d^{}(0)S(\lambda )}{Z},$$
(3)
with
$$Z=Z_0(\mu )+e^{\beta (ฯต_f\mu )}Z_0(\mu U),$$
(4)
the atomic partition function expressed in terms of
$$Z_0(\mu )=\mathrm{Tr}_\mathrm{d}e^{\beta \widehat{H}_0}S(\lambda ),\widehat{H}_0=\mu d^{}d,$$
(5)
In the above equations, the atomic Hamiltonian $`\widehat{H}_{at}`$ is the Hamiltonian of Eq. (1) restricted to one site, with $`t^{}=0`$, and all time dependence is with respect to this atomic Hamiltonian. The evolution operator $`S(\lambda )`$ satisfies
$$S(\lambda )=\mathrm{exp}[_0^\beta d\tau _0^\beta d\tau ^{}d^{}(\tau )\lambda (\tau \tau ^{})d(\tau ^{})],$$
(6)
with $`\lambda (\tau \tau ^{})`$ a time-dependent atomic field adjusted to make the atomic Greenโs function equal to the local lattice Greenโs function. We define an effective medium by
$$G_0^1(i\omega _n)=G_n^1+\mathrm{\Sigma }_n=i\omega _n+\mu \lambda _n,$$
(7)
with $`\mathrm{\Sigma }_n`$ the local self-energy and $`\lambda _n`$ the Fourier transform of $`\lambda (\tau )`$. The trace in Eq. (3) can be evaluated directly to yield
$$G_n=w_0G_0(i\omega _n)+w_1[G_0^1(i\omega _n)U]^1,$$
(8)
with $`w_0=1w_1`$ and
$$w_1=\mathrm{exp}[\beta (ฯต_f\mu )]Z_0(\mu U)/Z.$$
(9)
The self-consistency relation needed to determine $`\lambda _n`$ and hence $`G_n`$ is to equate the local lattice Greenโs function to the atomic Greenโs function via
$$G_n=_{\mathrm{}}^{\mathrm{}}๐ฯต\frac{\rho (ฯต)}{i\omega _n+\mu \mathrm{\Sigma }_nฯต},$$
(10)
with $`\rho (ฯต)=\mathrm{exp}(ฯต^2)/\sqrt{\pi }`$ the noninteracting density of states for the infinite-dimensional hypercubic lattice.
The iterative algorithm to solve for $`G_n`$ starts with $`\mathrm{\Sigma }_n=0`$. Then Eq. (10) is used to find $`G_n`$, Eq. (7) is employed to extract the effective medium, Eq. (8) is used to find a new local Greenโs function, and then Eq. (7) is used to find the new self-energy. The algorithm is then repeated until it converges, which usually requires only about a dozen or so iterations. In this contribution, we examine the half-filled case $`\rho _d=_in_i/N=0.5`$ and $`\rho _f=_iw_i/N=0.5`$, which corresponds to $`\mu =U/2`$ and $`ฯต_f=U/2`$.
In the following section we illustrate how the Baym-Kadanoff formalism can be used to determine the dynamical charge susceptibility. Numerical results at half filling are presented in Section III and conclusions in Section IV.
## II Baym-Kadanoff formalism
The momentum-dependent susceptibility satisfies the following Dysonโs equation:
$`\chi ^{dd}(๐ช,i\omega _m,i\omega _n;i\nu _l)`$ $`=`$ $`\chi _0^{dd}(๐ช,i\omega _m;i\nu _l)\delta _{mn}`$ (11)
$``$ $`T{\displaystyle \underset{n^{}}{}}\chi _0^{dd}(๐ช,i\omega _m;i\nu _l)\mathrm{\Gamma }(i\omega _m,i\omega _n^{};i\nu _l)\chi ^{dd}(๐ช,i\omega _n^{},i\omega _n;i\nu _l),`$ (12)
and the full susceptibility is found by summing over the Fermionic Matsubara frequencies $`\chi ^{dd}(๐ช,i\nu _l)=T_{mn}\chi ^{dd}(๐ช,i\omega _m,i\omega _n;i\nu _l)`$. In Eq. (12) the bare susceptibility satisfies
$`\chi _0^{dd}(X,i\omega _m;i\nu _l)`$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle \underset{๐ค}{}}G_m(๐ค)G_{m+l}(๐ค+๐ช),`$ (13)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{1X^2}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐ฯต{\displaystyle \frac{\rho (ฯต)}{i\omega _m+\mu \mathrm{\Sigma }_mฯต}}F_{\mathrm{}}\left({\displaystyle \frac{i\omega _{m+l}+\mu \mathrm{\Sigma }_{m+l}Xฯต}{\sqrt{1X^2}}}\right),`$ (14)
with $`X(๐ช)=lim_d\mathrm{}_{i=1}^d\mathrm{cos}๐ช_๐ข/d`$ describing all of the momentum dependence of the susceptibility and
$$F_{\mathrm{}}(z)=_{\mathrm{}}^{\mathrm{}}๐ฯต\frac{\rho (ฯต)}{zฯต},$$
(15)
the Hilbert transform of the noninteracting density of states.
The dynamical susceptibility simplifies in three cases: the chessboard case \[$`๐ช=(\pi ,\pi ,\mathrm{},\pi ),`$ $`X=1`$\] where
$$\chi _0^{dd}(1,i\omega _m;i\nu _l)=\frac{G_m+G_{m+l}}{i\omega _m+i\omega _{m+l}+2\mu \mathrm{\Sigma }_m\mathrm{\Sigma }_{m+l}};$$
(16)
the local case ($`X=0`$) where
$$\chi _0^{dd}(0,i\omega _m;i\nu _l)=G_mG_{m+l};$$
(17)
and the uniform case \[$`๐ช=(0,0,\mathrm{},0),`$ $`X=1`$\] where
$$\chi _0^{dd}(1,i\omega _m;i\nu _l)=\frac{G_mG_{m+l}}{i\nu _l+\mathrm{\Sigma }_m\mathrm{\Sigma }_{m+l}}.$$
(18)
We will only be interested in these three simpler cases here. Note that in instances where the denominators of Eqs. (16) and (18) vanish, the susceptibility is evaluated by lโHรดpitalโs rule.
The calculation of the full susceptibility requires the local irreducible vertex function. The Baym-Kadanoff approach solves for the irreducible vertex function in a manner which guarantees that an approximation scheme maintains the conservation laws of the Hamiltonian. The procedure requires that an approximate self-energy, $`\mathrm{\Sigma }(๐ซ_\mathrm{๐},\tau _1,๐ซ_\mathrm{๐},\tau _2)`$ be given by the functional derivative of a free-energy functional, $`\mathrm{\Phi }`$, with respect to the full Greenโs function, $`G(๐ซ_\mathrm{๐},\tau _1,๐ซ_\mathrm{๐},\tau _2)`$. That is:
$$\mathrm{\Sigma }(๐ซ_\mathrm{๐},\tau _1,๐ซ_\mathrm{๐},\tau _2)=\frac{\delta \mathrm{\Phi }}{\delta G(๐ซ_\mathrm{๐},\tau _1,๐ซ_\mathrm{๐},\tau _2)},$$
(19)
and then differentiation of the self-energy with respect to the Greenโs function produces the irreducible vertex function
$$\mathrm{\Gamma }(๐ซ_\mathrm{๐},\tau _1,๐ซ_\mathrm{๐},\tau _2;๐ซ_\mathrm{๐},\tau _3,๐ซ_\mathrm{๐},\tau _3)=\frac{1}{T}\frac{\delta \mathrm{\Sigma }(๐ซ_\mathrm{๐},\tau _1,๐ซ_\mathrm{๐},\tau _2)}{\delta G(๐ซ_\mathrm{๐},\tau _3,๐ซ_\mathrm{๐},\tau _3)}.$$
(20)
Of course, this technique can also be used for an exact solution as we do here.
Using the dynamical-mean-field approximation for the Falicov-Kimball model in a spatially invariant system yields an exact expression for the self-energy $`\mathrm{\Sigma }=G_0^1G^1`$ \[as seen in Eq. (7)\] and we do not need to find the appropriate free-energy functional $`\mathrm{\Phi }`$. Our strategy is to calculate the Greenโs function of the effective medium, $`G_0(\tau _1,\tau _2)`$, when there is an additional time-dependent field, $`\chi (\tau )`$, which couples to the charge density, $`d(\tau )d^{}(\tau )`$. \[i.e. we add $`_0^\beta d\tau \chi (\tau )d^{}(\tau )d(\tau )`$ to the exponent of the evolution operator in Eq. (6)\]. This time-dependent field removes time-translational invariance from the system, so the Greenโs functions now depend on $`\tau _1`$ and $`\tau _2`$ separately. Using $`G_0(\tau _1,\tau _2)`$, we evaluate the self-energy as an explicit function of the full Greenโs function, $`G(\tau _1,\tau _2)`$, including terms to linear order in $`\chi (\tau )`$. We take the derivative of the self-energy with respect to the full Greenโs function to obtain the vertex function and afterwards set the field, $`\chi (\tau )`$ to zero. The calculation differs from standard approaches, in that $`\chi (\tau )`$ provides a time-dependence so that $`G_0(\tau _1,\tau _2)`$ depends separately on $`\tau _1`$ and $`\tau _2`$, not just on the difference, $`(\tau _1\tau _2)`$.
To begin, we introduce an auxiliary Greenโs function (as Brandt and Urbanek did in their calculation of the $`f`$-spectral function) in the presence of just the charge-coupled field, $`\chi `$, defined by:
$$g_{aux}(\tau _1,\tau _2;\mu )=\frac{\mathrm{Tr}_\mathrm{d}T_\tau e^{\beta \widehat{H}_0}e^{_0^\beta d\overline{\tau }\chi (\overline{\tau })d^{}(\overline{\tau })d(\overline{\tau })}d(\tau _1)d^{}(\tau _2)}{1+e^{\beta \mu }e^{_0^\beta d\overline{\tau }\chi (\overline{\tau })d^{}(\overline{\tau })d(\overline{\tau })}},$$
(21)
where $`\widehat{H}_0=\mu d^{}d`$ and we expand the time-dependent field, $`\chi `$, in a Fourier series of the Bosonic frequencies
$$\chi (\tau )=\underset{l}{}\chi _le^{i\nu _l\tau }.$$
(22)
Noting that $`d(\tau )=e^{\tau \widehat{H}_0}d(0)e^{\tau \widehat{H}_0}=e^{\tau \mu }d(0)`$, and taking into account the time ordering, we can easily perform the trace in Eq.(21), which yields (to linear order in $`\chi _l`$):
$$g_{aux}(\tau _1,\tau _2;\mu )=e^{\mu (\tau _1\tau _2)}\left[1+\underset{l0}{}\frac{\chi _l}{i\nu _l}\left(e^{i\nu _l\tau _1}e^{i\nu _l\tau _2}\right)\right]\left[\frac{\theta \left(\tau _1\tau _2\right)}{1+e^{\beta \mu }e^{\beta \chi _0}}+\frac{\theta \left(\tau _2\tau _1\right)}{1+e^{\beta \mu }e^{\beta \chi _0}}\right],$$
(23)
which depends on both $`\tau _1`$ and $`\tau _2`$ (not just $`\tau _1\tau _2`$) because of the terms with $`\chi _{l0}`$. Nevertheless, the following symmetries are readily proven:
for $`\tau _1`$ $`<\tau _2<\tau _1+\beta `$ (24)
$`g_{aux}`$ $`(\tau _1+\beta ,\tau _2;\mu )=g_{aux}(\tau _1,\tau _2;\mu ),`$ (25)
and
for $`\tau _2`$ $`<\tau _1<\tau _2+\beta `$ (26)
$`g_{aux}`$ $`(\tau _1,\tau _2+\beta ;\mu )=g_{aux}(\tau _1,\tau _2;\mu ).`$ (27)
Hence we can find the Fourier transform of $`g_{aux}(\tau _1,\tau _2;\mu )`$ in terms of the Fermionic Matsubara frequencies
$$g_{aux}(i\omega _n,i\omega _n^{};\mu )=_0^\beta d\tau _1_{\beta +\tau _1}^{\tau _1}d\tau _2e^{i\omega _n\tau _1}e^{i\omega _n^{}\tau _2}g_{aux}(\tau _1,\tau _2;\mu ).$$
(28)
The result is that $`g_{aux}(i\omega _n,i\omega _n^{};\mu )`$ contains terms that are diagonal in $`n,n^{}`$ and terms shifted off the diagonal by $`l`$ that are proportional to $`\chi _l`$:
$$g_{aux}(i\omega _n,i\omega _n^{};\mu )=\frac{\delta _{n,n^{}}}{i\omega _n+\mu }+\underset{l}{}\frac{\delta _{n+l,n^{}}\chi _l}{i\nu _l}\left(\frac{1}{i\omega _n^{}+\mu }\frac{1}{i\omega _n+\mu }\right).$$
(29)
The Greenโs function for the effective medium of the atomic problem, $`G_0(\tau ,\tau ^{})`$, is obtained as before by including a local time-dependent field, $`\lambda (\tau \tau ^{})`$, that incorporates the effects of propagation on the lattice and only depends on the time-difference:
$$G_0(\tau ,\tau ^{})=\frac{\mathrm{Tr}_\mathrm{d}T_\tau e^{\beta \widehat{H}_0}e^{_0^\beta d\overline{\tau }\chi (\overline{\tau })d^{}(\overline{\tau })d(\overline{\tau })}S(\lambda )d(\tau _1)d^{}(\tau _2)}{Z_0(\mu )}.$$
(30)
As the local field only contains diagonal frequency components, $`\lambda (i\omega _n)\delta _{nn^{}}`$, the effective-medium Greenโs function is easily obtained from the auxiliary Greenโs function through:
$$\left[G_0\right]_{n,n^{}}^1=\left[g_{aux}(\mu )\right]_{n,n^{}}^1\lambda (i\omega _n)\delta _{n,n^{}},$$
(31)
where the Greenโs function matrices are inverted in frequency coordinates, such that $`\left[G_0\right]_{n,n^{}}^1`$ represents the $`n,n^{}`$ components of the inverse of the matrix $`G_0(i\omega _n,i\omega _n^{})`$. A key aspect of our derivation of the vertex function, is that to linear order in a single frequency component of the charge-coupled field, $`\chi _l`$, the Greenโs function for the effective medium, $`G_0(i\omega _n,i\omega _n^{})`$, like the auxiliary Greenโs function, $`g_{aux}(i\omega _n,i\omega _n^{};\mu )`$, only contains components in the two diagonals given by $`\delta _{n,n^{}}`$ and $`\delta _{n+l,n^{}}`$.
The full Greenโs function is defined by:
$$G(\tau _1,\tau _2)=\frac{\mathrm{Tr}_\mathrm{f}\mathrm{Tr}_\mathrm{d}T_\tau e^{\beta \widehat{H}_{at}}e^{_0^\beta d\overline{\tau }\chi (\overline{\tau })d^{}(\overline{\tau })d(\overline{\tau })}S(\lambda )d(\tau _1)d^{}(\tau _2)}{Z},$$
(32)
where $`Z`$ is given by Eq. (4) and the time dependence of the $`d`$-operators is governed by $`\widehat{H}_{at}`$. The Greenโs function contains a trace over $`f`$-electron states, which leads to two terms:
$`G(\tau _1,\tau _2)`$ $`=`$ $`{\displaystyle \frac{Tr_dT_\tau e^{\beta \widehat{H}_0}e^{_0^\beta d\overline{\tau }\chi (\overline{\tau })d^{}(\overline{\tau })d(\overline{\tau })}S(\lambda )d(\tau _1)d^{}(\tau _2)}{Z}}`$ (34)
$`+e^{\beta (ฯต_f\mu )}{\displaystyle \frac{Tr_dT_\tau e^{\beta \widehat{H}_1}e^{_0^\beta d\overline{\tau }\chi (\overline{\tau })d^{}(\overline{\tau })d(\overline{\tau })}S(\lambda )d(\tau _1)d^{}(\tau _2)}{Z}},`$
where $`\widehat{H}_0=\mu d^{}d`$ is the atomic Hamiltonian with no $`f`$-electrons, and $`\widehat{H}_1=(U\mu )d^{}d`$ is the atomic Hamiltonian for $`d`$-electrons in the presence of one $`f`$-electron.
We can now relate the full Greenโs function to the Greenโs function of the effective medium through the following matrix equation:
$$G(i\omega _n,i\omega _n^{})=w_0G_0(i\omega _n,i\omega _n^{})+w_1\left[\left(G_0^1\right)U\right]_{n,n^{}}^1$$
(35)
where $`w_0=1w_1`$ and $`w_1`$ are given by Eq. (9) and equal the fraction of sites with an occupancy of zero and one $`f`$-electron, respectively. In the above, and all following equations, $`U`$ is the only matrix which is necessarily diagonal in frequency space.
Eq.(35) can be rearranged by multiplying on the left or right by matrices like $`G^1`$, $`G_0^1`$ and $`G_0^1U`$ to give the two following matrix equations (with indices suppressed):
$`G_0^2(U+G^1)G_0^1+(1w_1)UG^1`$ $`=`$ $`0,`$ (36)
$`G_0^2G_0^1(U+G^1)+(1w_1)UG^1`$ $`=`$ $`0.`$ (37)
Adding these two equations together and collecting terms (noting the noncommutativity and hence the ordering of the matrices) then yields the following quadratic matrix equation:
$$\left[G_0^1\frac{1}{2}(U+G^1)\right]^2\frac{1}{4}(U+G^1)^2+(1w_1)UG^1=0.$$
(38)
Substitution for $`G_0`$ by $`\mathrm{\Sigma }`$, using Dysonโs equation in the form
$$\mathrm{\Sigma }(i\omega _n,i\omega _m)=\left[G_0^1\right]_{n,m}\left[G^1\right]_{n,m},$$
(39)
yields:
$$\left[\mathrm{\Sigma }+\frac{1}{2}\left(G^1U\right)\right]^2=\frac{1}{4}\left[U^2+2(2w_11)UG^1+G^2\right],$$
(40)
where each term in $`G`$ and $`\mathrm{\Sigma }`$ is a matrix in frequency-space and $`U`$ multiplies the identity matrix. Note that in the limit $`\chi _l=0`$, Eq.(40) becomes diagonal, and reduces to the conventional quadratic expression for the self-energy in terms of $`G(i\omega _n)`$ and $`w_1`$ as first determined by Brandt and Mielsch .
The irreducible vertex function is defined in frequency space by:
$$\mathrm{\Gamma }(i\omega _n,i\omega _m;i\omega _n^{},i\omega _m^{})=_0^\beta d\tau _1_{\tau _1}^{\beta +\tau _1}d\tau _2_0^\beta d\tau _1^{}_{\tau _1^{}}^{\beta +\tau _1^{}}d\tau _2^{}e^{i\omega _n\tau _1i\omega _m\tau _2+i\omega _n^{}\tau _1^{}i\omega _m^{}\tau _2^{}}\mathrm{\Gamma }(\tau _1,\tau _2;\tau _1^{},\tau _2^{}).$$
(41)
As the vertex function is independent of the absolute time, we can change variables: $`\tau _1\tau _1\tau _2^{},\tau _2\tau _2\tau _2^{},\tau _1^{}\tau _1^{}\tau _2^{}`$ so that the $`\tau _2^{}`$ integral yields a delta function which requires $`i\omega _m^{}=i\omega _ni\omega _m+i\omega _n^{}`$. Hence, given the difference between two Matsubara frequencies, $`i\omega _mi\omega _n=i\nu _l`$ is a Bosonic frequency, the vertex function becomes
$$\mathrm{\Gamma }(i\omega _n,i\omega _n^{};i\nu _l)=\frac{1}{T}\frac{\delta \mathrm{\Sigma }(i\omega _n,i\omega _n+i\nu _l)}{\delta G(i\omega _n^{},i\omega _n^{}+i\nu _l)}.$$
(42)
The problem is greatly simplified by calculating $`\mathrm{\Sigma }_{n,m}`$ to first order in $`\chi `$, for a single frequency component, $`\chi _l`$. In this case, both $`\mathrm{\Sigma }_{n,m}`$ and $`G_{n,m}`$ contain diagonal terms, and off-diagonal terms linear in $`\chi _l`$ where $`m=n+l`$, so that:
$$\begin{array}{ccccc}\mathrm{\Sigma }_{n,m}\hfill & =\hfill & \mathrm{\Sigma }_n\delta _{n,m}\hfill & +\hfill & \overline{\mathrm{\Sigma }}_n\delta _{n+l,m},\hfill \\ G_{n,m}\hfill & =\hfill & G_n\delta _{n,m}\hfill & +\hfill & \overline{G}_n\delta _{n+l,m}.\hfill \end{array}$$
(43)
We are interested in the case where $`l0`$, as the zero-frequency response (corresponding to a shift of the chemical potential) is well-known .
Next, we calculate the diagonal and off-diagonal pieces of the self-energy. Eq.(40) simplifies from a full quadratic matrix equation, to the following coupled equations:
$`\left(\mathrm{\Sigma }_n{\displaystyle \frac{U}{2}}+{\displaystyle \frac{1}{2G_n}}\right)^2`$ $`=`$ $`{\displaystyle \frac{1}{4G_n^2}}\left\{1+2UG_n(2w_11)+U^2G_n^2\right\},`$ (44)
$`\left(\overline{\mathrm{\Sigma }}_n{\displaystyle \frac{\overline{G}_n}{2G_nG_{n+l}}}\right)\left(\mathrm{\Sigma }_n+{\displaystyle \frac{1}{2G_n}}+\mathrm{\Sigma }_{n+l}+{\displaystyle \frac{1}{2G_{n+l}}}U\right)`$ $`=`$ $`{\displaystyle \frac{\overline{G}_n}{4G_nG_{n+l}}}\left[2(12w_1)U{\displaystyle \frac{(G_n+G_{n+l})}{G_nG_{n+l}}}\right].`$ (45)
Eq.(44) contains only diagonal terms with no dependence on the finite-frequency field $`\chi _l`$, hence $`\mathrm{\Sigma }_n`$ is unchanged from its value when $`\chi _l=0`$. We concentrate on rearranging Eq.(45) to obtain the frequency-dependent response. First, we multiply both sides of Eq.(45) by
$$\left(\mathrm{\Sigma }_n\frac{U}{2}+\frac{1}{2G_n}\right)\left(\mathrm{\Sigma }_{n+l}\frac{U}{2}+\frac{1}{2G_{n+l}}\right)$$
(46)
and use Eq.(44) to replace the quadratic terms on the left, giving:
$`\left(\overline{\mathrm{\Sigma }}_n{\displaystyle \frac{\overline{G}_n}{2G_nG_{n+l}}}\right)`$ $`\left[2U(2w_11)+{\displaystyle \frac{1}{G_n}}+{\displaystyle \frac{1}{G_{n+l}}}\right]\left({\displaystyle \frac{1}{4G_n}}{\displaystyle \frac{1}{4G_{n+l}}}\right)`$ (47)
$`=`$ $`{\displaystyle \frac{\overline{G}_n}{4G_nG_{n+l}}}\left[2(12w_1)U{\displaystyle \frac{(G_n+G_{n+l})}{G_nG_{n+l}}}\right]\left(\mathrm{\Sigma }_n+{\displaystyle \frac{1}{2G_n}}\mathrm{\Sigma }_{n+l}{\displaystyle \frac{1}{2G_{n+l}}}\right).`$ (48)
The above form gives rise to considerable cancellation between the two sides, leading to the amazingly simple result:
$$\overline{\mathrm{\Sigma }}_n=\overline{G}_n\frac{\mathrm{\Sigma }_n\mathrm{\Sigma }_{n+l}}{G_nG_{n+l}}.$$
(49)
So the only terms in $`\mathrm{\Sigma }_{n,m}`$ that depend on a finite frequency component, $`i\nu _l`$, of the $`\chi `$ field are $`\mathrm{\Sigma }_{n,n+l}=\overline{\mathrm{\Sigma }}_n`$. The only terms which survive the limit, $`\chi _l0`$ after differentiation, to give the vertex function, are:
$$\mathrm{\Gamma }(i\omega _n,i\omega _n;i\nu _l)=\frac{1}{T}\frac{\delta \mathrm{\Sigma }(i\omega _n,i\omega _{n+l})}{\delta G(i\omega _n,i\omega _{n+l})}=\frac{1}{T}\frac{\delta \overline{\mathrm{\Sigma }}_n}{\delta \overline{G}_n}=\frac{1}{T}\frac{\mathrm{\Sigma }_n\mathrm{\Sigma }_{n+l}}{G_nG_{n+l}}.$$
(50)
The above simple result found from this Baym-Kadanoff analysis is identical to that found by diagrammatic calculations which give the more complicated form:
$$\mathrm{\Gamma }(i\omega _n,i\omega _n;i\nu _l)=\frac{1}{T}\frac{w_0w_1U^2}{\left(1+G_n\mathrm{\Sigma }_n\right)\left[1+G_n(\mathrm{\Sigma }_nU)\right]\left(1+G_{n+l}\mathrm{\Sigma }_{n+l}\right)\left[1+G_{n+l}(\mathrm{\Sigma }_{n+l}U)\right]+w_0w_1U^2G_nG_{n+l}}.$$
(51)
The proof of this equivalence is presented in the Appendix.
Since the vertex function is diagonal, we have a simple form for the dynamical charge susceptibility
$$\chi ^{dd}(X;i\nu _l0)=T\underset{m}{}\frac{\chi _0(X,i\omega _m;i\nu _l)}{1+\chi _0(X,i\omega _m;i\nu _l)\frac{\mathrm{\Sigma }_m\mathrm{\Sigma }_{m+l}}{G_mG_{m+l}}}.$$
(52)
In the limit $`๐ช=0`$ ($`X=1`$), substitution of $`\chi _0^{dd}`$ from Eq. (18) into Eq. (52) gives
$$\chi ^{dd}(1;i\nu _l0)=T\underset{m}{}\frac{G_mG_{m+l}}{i\nu _l}=0,$$
(53)
as expected by our symmetry arguments. Notice that $`\chi _0^{dd}`$ does not vanish; it is the vertex corrections that force the full susceptibility to vanish.
The final step in our formalism development is to perform the analytic continuation of Eq. (52) from the imaginary to the real axis. Our method is not completely rigorous, because we are unable to verify necessary analyticity arguments as described below, but we compare the direct calculation of the susceptibility on the imaginary axis to the inferred value there that is found by using the spectral formula from the real axis. In nearly all cases, we have accuracy to better than one part in 1000, which supports, a posteriori, that our technique is valid.
We will consider only the case with $`X=1`$ here, since the $`X=1`$ susceptibility is trivial and since the $`X=0`$ case is simpler and can easily be worked out by following the same steps we use for the chessboard case. We start by replacing the summation over Matsubara frequencies in Eq. (52) by three contour integrals, which are chosen as in Figure 1 to encircle all of the simple poles at the Fermionic Matsubara frequencies, and no other poles in the corresponding integrals \[in other words, the only poles contributing are the simple poles from $`f(\omega )`$ below\]. The susceptibility then becomes (for simplicity, we consider $`\nu _l>0`$ here)
$`\chi ^{dd}(1;i\nu _l)`$ $`=`$ $`{\displaystyle \frac{i}{2\pi }}{\displaystyle _{C_1}}๐\omega f(\omega ){\displaystyle \frac{\frac{G_R(\omega )+G_R(\omega +i\nu _l)}{2\omega +2\mu +i\nu _l\mathrm{\Sigma }_R(\omega )\mathrm{\Sigma }_R(\omega +i\nu _l)}}{1\frac{G_R(\omega )+G_R(\omega +i\nu _l)}{2\omega +2\mu +i\nu _l\mathrm{\Sigma }_R(\omega )\mathrm{\Sigma }_R(\omega +i\nu _l)}\frac{\mathrm{\Sigma }_R(\omega )\mathrm{\Sigma }_R(\omega +i\nu _l)}{G_R(\omega )G_R(\omega +i\nu _l)}}}`$ (54)
$`+`$ $`{\displaystyle \frac{i}{2\pi }}{\displaystyle _{C_2}}๐\omega f(\omega ){\displaystyle \frac{\frac{G_A(\omega )+G_R(\omega +i\nu _l)}{2\omega +2\mu +i\nu _l\mathrm{\Sigma }_A(\omega )\mathrm{\Sigma }_R(\omega +i\nu _l)}}{1\frac{G_A(\omega )+G_R(\omega +i\nu _l)}{2\omega +2\mu +i\nu _l\mathrm{\Sigma }_A(\omega )\mathrm{\Sigma }_R(\omega +i\nu _l)}\frac{\mathrm{\Sigma }_A(\omega )\mathrm{\Sigma }_R(\omega +i\nu _l)}{G_A(\omega )G_R(\omega +i\nu _l)}}}`$ (55)
$`+`$ $`{\displaystyle \frac{i}{2\pi }}{\displaystyle _{C_3}}๐\omega f(\omega ){\displaystyle \frac{\frac{G_A(\omega )+G_A(\omega +i\nu _l)}{2\omega +2\mu +i\nu _l\mathrm{\Sigma }_A(\omega )\mathrm{\Sigma }_A(\omega +i\nu _l)}}{1\frac{G_A(\omega )+G_A(\omega +i\nu _l)}{2\omega +2\mu +i\nu _l\mathrm{\Sigma }_A(\omega )\mathrm{\Sigma }_A(\omega +i\nu _l)}\frac{\mathrm{\Sigma }_A(\omega )\mathrm{\Sigma }_A(\omega +i\nu _l)}{G_A(\omega )G_A(\omega +i\nu _l)}}},`$ (56)
where $`f(\omega )=1/[1+\mathrm{exp}(\beta \omega )]`$ is the Fermi function and the subscript $`R`$ or $`A`$ refers to the retarded or advanced Greenโs function (or self-energy). The choices of the subscripts are so that the Greenโs functions and self-energies are analytic within regions 1, 2, or 3. Hence the contours can be deformed until they run parallel to the real axis, as shown in Figure 2. We are making an assumption that there are no other poles present within these regions, when we deform the contours. In fact, if we make an analytic continuation of just the bare susceptibility, then after continuing $`i\nu _l`$ to $`\nu +i\delta `$ we do find a pole that lies in region 2 just below the real axis (at $`\omega =\nu /2i\delta )`$, but the residue of this pole vanishes when $`\nu `$ lies on the real axis. It is more difficult to make such an analysis for the full susceptibility, so we rely instead on the comparison with the direct calculation on the imaginary axis.
When we evaluate the integrals along the lines indicated in Figure 2, we will evaluate the Fermi function at $`\omega i\nu _l`$, which we set equal to $`f(\omega )`$ before continuing the $`\nu `$ frequency. Then we can evaluate the final result, which becomes
$`\chi ^{dd}(1;\nu )={\displaystyle \frac{i}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\omega \{`$ $`f`$ $`(\omega ){\displaystyle \frac{\frac{G(\omega )+G(\omega +\nu )}{2\omega +2\mu +\nu \mathrm{\Sigma }(\omega )\mathrm{\Sigma }(\omega +\nu )}}{1\frac{G(\omega )+G(\omega +\nu )}{2\omega +2\mu +\nu \mathrm{\Sigma }(\omega )\mathrm{\Sigma }(\omega +\nu )}\frac{\mathrm{\Sigma }(\omega )\mathrm{\Sigma }(\omega +\nu )}{G(\omega )G(\omega +\nu )}}}`$ (57)
$``$ $`f(\omega \nu ){\displaystyle \frac{\frac{G^{}(\omega )+G^{}(\omega +\nu )}{2\omega +2\mu +\nu \mathrm{\Sigma }^{}(\omega )\mathrm{\Sigma }^{}(\omega +\nu )}}{1\frac{G^{}(\omega )+G^{}(\omega +\nu )}{2\omega +2\mu +\nu \mathrm{\Sigma }^{}(\omega )\mathrm{\Sigma }^{}(\omega +\nu )}\frac{\mathrm{\Sigma }^{}(\omega )\mathrm{\Sigma }^{}(\omega +\nu )}{G^{}(\omega )G^{}(\omega +\nu )}}}`$ (58)
$``$ $`[f(\omega )f(\omega \nu )]{\displaystyle \frac{\frac{G^{}(\omega )+G(\omega +\nu )}{2\omega +2\mu +\nu \mathrm{\Sigma }^{}(\omega )\mathrm{\Sigma }(\omega +\nu )}}{1\frac{G^{}(\omega )+G(\omega +\nu )}{2\omega +2\mu +\nu \mathrm{\Sigma }^{}(\omega )\mathrm{\Sigma }(\omega +\nu )}\frac{\mathrm{\Sigma }^{}(\omega )\mathrm{\Sigma }(\omega +\nu )}{G^{}(\omega )G(\omega +\nu )}}}\},`$ (59)
and we replaced the advanced functions by the complex conjugate of the retarded Greenโs function (which is valid on the real axis).
## III Numerical Results
We performed a number of different numerical calculations of the charge response. To begin, we summarize the $`d`$-electron spectral functions in the weak-coupling $`(U=0.25)`$ limit, the moderate coupling limit $`(U=1)`$ and the strong-coupling limit $`(U=4)`$, which are illustrated in Figure 3. Note that the interacting density of states is temperature independent for this model (as first shown by Van Dongen). In the weak-coupling case $`(U=0.25)`$ the density of states looks Gaussian, with only small modifications due to the interactions. At moderate coupling, $`(U=1)`$, we find that a pseudogap appears in the density of states at the chemical potential. Note that this is a correlation-induced pseudogap and is not resulting from the charge-density-wave order of the ground state, since these results are for the high-temperature homogeneous phase. Finally, a true gap develops in strong coupling $`(U=4)`$, which rapidly approaches $`U`$ as $`U`$ increases in magnitude. In Figure 4 we show plots of the static chessboard susceptibility for the three different values of $`U`$ and inset a plot of the transition temperature as a function of $`U/(U+t^{})`$. The chain-dotted lines are guides to the eye for $`T_c`$ at the three values of $`U`$. The transition temperature has a classic formโincreasing for weak-$`U`$, reaching a maximum at moderate $`U`$ and then decreasing in the strong-coupling regime. Notice how the static susceptibility diverges as we approach $`T_c`$ because we are in the thermodynamic limit. In Figure 5 we show the dynamical chessboard susceptibility at five different temperatures ranging from well above $`T_c`$ to just above $`T_c`$. Notice how there is a temperature dependence to the charge susceptibility, but how there is no indication of the phase transition seen in the static susceptibility (recall the susceptibility is discontinuous at $`\nu =0`$ and we are only plotting the continuous piece of the susceptibility). In Figure 6 we show the same plot for the local dynamical susceptibility. Its behavior is quite similar to that of the chessboard susceptibility, but with a less marked temperature dependence. The expected peak in the imaginary part at low frequencies due to quasiparticle excitations can be seen, and it has a classic โtriangularโ shape as expected in a noninteracting system (recall the local susceptibility is the average over all momentum vectors, so we expect a linear contribution at small frequency and a curved drop off at high frequencies as we reach the approximate half bandwidth). The case of moderate coupling is shown in Figures 7 and 8. Here, the chessboard susceptibility actually decreases as $`T`$ is lowered because of the pseudogap in the density of states. Once again the temperature dependence of the local susceptibility is much less than the chessboard susceptibility (and we are including data below the $`T_c`$ to the charge-density-wave instability). The development from a single low-energy peak in the real part, to two nearly overlapping peaks (corresponding to charge transfer excitations), is clear even at moderate coupling where the density of states has only a pseudogap. Note once again that these peak developments arise from correlation effects and not from an underlying charge-density-wave order because these calculations are all performed within the homogeneous phase (even when we go below $`T_c`$).
In Figures 9 and 10 we show the results at strong coupling $`U=4`$. The charge susceptibility shows a number of interesting new features in the strong-coupling regime. The charge-transfer peaks, which are easiest seen as broad peaks in the imaginary part of the susceptibility, are present at the expected locations of $`\pm U`$. At high temperatures, the system has a large peak in the real part of the susceptibility near zero frequency, which rapidly decreases as $`T`$ is lowered, and becomes unnoticable at low temperatures. The imaginary part goes from having a linear behavior at low frequency to being essentially zero. Hence, in addition to the energy scale of the order of $`U`$ (corresponding to the large charge-transfer peaks in the imaginary part of the susceptibility) there is a low-energy scale on the order of $`t^2/U`$ that determines the low-temperature evolution of the charge excitations and the energy scale of those low-energy excitations. The high-energy features are โfrozen inโ at a high temperature scale, and there is strong temperature evolution in the low-frequency regime until $`T`$ becomes smaller than the low-energy scale where the low-energy features are โfrozen in.โ We believe these low-energy charge excitations are associated with virtual hopping processes, where an electron virtually hops onto a site occupied by a localized electron, and then hops off that site. Such processes turn off at temperatures below $`t^2/U`$. Once again the dependence of the chessboard and the local susceptibilities are similar and there is no signal of the phase transition in the dynamical piece of the charge susceptibility. We find that we lose some accuracy in our calculations at strong coupling, perhaps arising from the gap in the density of states and its effect on the low-temperature charge dynamics. Another possibility is that there are additional poles and resides that need to be taken into account in our โapproximateโ analytic continuation. One illustration of the numerical difficulties is that we no longer have an accuracy of one part in a thousand (it is typically a few parts in a thousand) when we evaluate the charge susceptibility on the imaginary axis by using the spectral formula and comparing with a direct calculation on the imaginary axis. Another illustration is that at the lowest temperatures the imaginary part of the spectral function becomes slightly negative for a frequency range in the vicinity of $`0.5<\nu <2`$ (where the imaginary part is approximately $`0.0003`$), which must be a numerical artifact due to the fact that the spectral function is known to be nonnegative for positive values of the frequency.
## IV Conclusions
We have presented a nontrivial example of the Baym-Kadanoff formalism to derive the dynamical susceptibilities of the Falicov-Kimball model. The dynamical susceptibility breaks into two pieces, a static piece that reflects the coupling of the system to the charge-density-wave distortions, and a frequency-dependent piece, which is less affected by any underlying charge-density-wave order. We developed the formalism to evaluate this frequency-dependent piece of the susceptibility exactly on the imaginary axis. We also developed an analytic-continuation scheme that we were unable to establish with full rigor due to the possibility of a small contribution from neglected residues of poles that could arise within the analytic-continuation procedure, but which is nevertheless quite accurate over a wide range of parameters.
We find a number of interesting properties of this model: (1) the static and dynamical susceptibilities are decoupledโthe static susceptibility diverges at $`T_c`$, but no signal of this phase transition is seen in the dynamical susceptibility; (2) the momentum dependence of the dynamical susceptibility is not too strong around the chessboard pointโwe see little variation between the $`X=1`$ and $`X=0`$ susceptibilities, but the $`X=1`$ susceptibility vanishes by symmetry, so momentum dependence is stronger near the Brillouin zone center; and (3) the dynamical susceptibilities do possess temperature dependence even though the interacting density of states is temperature independent. The temperature dependence is most striking in the strong-coupling regime where a low-energy scale on the order of $`t^2/U`$ arises and the charge susceptibility has a strong temperature dependence near this approximate temperature.
Our application of the Baym-Kadanoff formalism provides an interesting counterpoint to the diagrammatic derivation about the atomic limit, and can be applied to other related problems where the so-called โstatic approximationโ is exact. We have not examined any effects away from half-filling or at incommensurate wavevectors here. We leave those tasks to another publication.
###### Acknowledgements.
This research was supported by the Office of Naval Research under grant number N00014-99-1-0328. We acknowledge useful discussions with T. Devereaux, D. Hess, J. Serene, and A. Shvaika.
## A Equivalence to the atomic diagrammatic approach
To compare our result with the diagrammatic approach , we must begin with Shvaikaโs result for the frequency-dependent piece of the reducible vertex function,
$`\stackrel{~}{\mathrm{\Gamma }}(`$ $`i\omega _n`$ $`,i\omega _n^{};i\nu _l0)=`$ (A1)
$`{\displaystyle \frac{1}{T}}{\displaystyle \frac{\delta _{nn^{}}}{G_nG_n^{}G_{n+l}G_{n^{}+l}}}{\displaystyle \frac{U^2w_1(1w_1)}{\left(i\omega _n+\mu \lambda _n\right)\left(i\omega _n+\mu \lambda _nU\right)\left(i\omega _{n+l}+\mu \lambda _{n+l}\right)\left(i\omega _{n+l}+\mu \lambda _{n+l}U\right)}},`$ (A2)
By making the replacement, $`i\omega _n+\mu \lambda _n=\left[G_0(i\omega _n)\right]^1=\left(G_n\right)^1+\mathrm{\Sigma }_n`$, Eq.(A2) is equivalent to:
$$\frac{1}{T}\frac{\delta _{nn^{}}w_0w_1U^2}{\left(1+G_n\mathrm{\Sigma }_n\right)\left[1+G_n(\mathrm{\Sigma }_nU)\right]\left(1+G_{n+l}\mathrm{\Sigma }_{n+l}\right)\left[1+G_{n+l}(\mathrm{\Sigma }_{n+l}U)\right]}.$$
(A3)
Solving Dysonโs equation for the irreducible vertex function, $`\mathrm{\Gamma }(i\omega _n,i\omega _n^{};i\nu _l0)`$:
$$\stackrel{~}{\mathrm{\Gamma }}(i\omega _n,i\omega _n^{};i\nu _l0)=\mathrm{\Gamma }(i\omega _n,i\omega _n^{};i\nu _l0)T\underset{n^{\prime \prime }}{}\mathrm{\Gamma }(i\omega _n,i\omega _{n^{\prime \prime }};i\nu _l0)G_{n^{\prime \prime }}G_{n^{\prime \prime }+l}\stackrel{~}{\mathrm{\Gamma }}(i\omega _{n^{\prime \prime }},i\omega _n^{};i\nu _l0),$$
(A4)
then yields Eq. (51).
In order to prove that our simple result, Eq.(50) and the diagrammatic result, Eq.(51) are equivalent to each other, we first demonstrate the existence of a useful identity. Eq.(44) can be expanded and multiplied by $`G_n`$ to yield:
$`w_1U`$ $`=`$ $`\mathrm{\Sigma }_n\left(1UG_n+\mathrm{\Sigma }_nG_n\right),`$ (A5)
$`=`$ $`\mathrm{\Sigma }_{n+l}\left(1UG_{n+l}+\mathrm{\Sigma }_{n+l}G_{n+l}\right),`$ (A6)
where the second equation is evaluated equivalently, with $`nn+l`$. Eq. (A5) is used twice to simplify the product given below:
$`\left(1+G_n\mathrm{\Sigma }_n\right)\left[1+G_n(\mathrm{\Sigma }_nU)\right]`$ $`=`$ $`1+G_n\mathrm{\Sigma }_nUG_n+w_1U,`$ (A7)
$`=`$ $`w_1U\left({\displaystyle \frac{1}{\mathrm{\Sigma }_n}}+G_n\right),`$ (A8)
and similarly for $`nn+l`$.
Hence the expression on the right of Eq. (51) is rewritten as:
$$\frac{1}{T}\frac{w_0w_1U^2}{w_1^2U^2\left(\frac{1}{\mathrm{\Sigma }_n}+G_n\right)\left(\frac{1}{\mathrm{\Sigma }_{n+l}}+G_{n+l}\right)+G_nG_{n+l}w_0w_1U^2},$$
(A9)
and multiplication of top and bottom by $`(\mathrm{\Sigma }_n\mathrm{\Sigma }_{n+l})\mathrm{\Sigma }_n\mathrm{\Sigma }_{n+l}`$ leads to:
$$\frac{1}{T}\frac{(\mathrm{\Sigma }_n\mathrm{\Sigma }_{n+l})\mathrm{\Sigma }_n\mathrm{\Sigma }_{n+l}(1w_1)}{w_1\left[\mathrm{\Sigma }_n\mathrm{\Sigma }_{n+l}+G_n\mathrm{\Sigma }_n(\mathrm{\Sigma }_n\mathrm{\Sigma }_{n+l})+G_{n+l}\mathrm{\Sigma }_{n+l}(\mathrm{\Sigma }_n\mathrm{\Sigma }_{n+l})\right]+G_nG_{n+l}\mathrm{\Sigma }_n\mathrm{\Sigma }_{n+l}(\mathrm{\Sigma }_n\mathrm{\Sigma }_{n+l})}.$$
(A10)
Now Eqs. (A5-A6) are further used to make the replacements $`G_n\mathrm{\Sigma }_n^2w_1U\mathrm{\Sigma }_n+UG_n\mathrm{\Sigma }_n`$ and $`G_{n+l}\mathrm{\Sigma }_{n+l}^2w_1U\mathrm{\Sigma }_{n+l}+UG_{n+l}\mathrm{\Sigma }_{n+l}`$ for all quadratic terms in the denominator. Numerous cancellations then lead to:
$`\mathrm{\Gamma }(i\omega _n,i\omega _n;i\nu _l0)`$ $`=`$ $`{\displaystyle \frac{1}{T}}{\displaystyle \frac{(\mathrm{\Sigma }_n\mathrm{\Sigma }_{n+l})\mathrm{\Sigma }_n\mathrm{\Sigma }_{n+l}(1w_1)}{w_1\left(\mathrm{\Sigma }_{n+l}\mathrm{\Sigma }_nG_{n+l}\mathrm{\Sigma }_{n+l}\mathrm{\Sigma }_nG_n\right)\mathrm{\Sigma }_{n+l}\mathrm{\Sigma }_nG_{n+l}+\mathrm{\Sigma }_{n+l}\mathrm{\Sigma }_nG_n}},`$ (A11)
$`=`$ $`{\displaystyle \frac{1}{T}}{\displaystyle \frac{\mathrm{\Sigma }_n\mathrm{\Sigma }_{n+l}}{G_nG_{n+l}}},`$ (A12)
which yields our final result, Eq. (50).
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# 1 Introduction
## 1 Introduction
So far in the literature, the study of null orbits has been restricted to isometries only. The groups $`G_r`$, $`r4`$, on $`N_3`$ have at least one subgroup $`G_3`$ which may act on $`N_3`$, $`N_2`$ or $`S_2`$ . For $`G_3`$ on $`S_2`$, one obtains special cases of the LRS models, $`G_r`$ admitting either a group $`G_3`$ on $`N_3`$ or a null Killing vector . The case $`G_3`$ on $`N_2`$ was studied by Barnes , the group $`G_3`$ is then of Bianchi type $`II`$ and perfect fluid solutions are excluded since the metric leads to a Ricci tensor whose Segre type is not that of a perfect fluid. Another case that has been considered in the literature is that of a $`G_3`$ on $`N_3`$, such case is subject to the condition $`R_{ab}k^ak^b=0`$ and this condition excludes perfect fluids with $`\mu +p0`$. Perfect fluid solutions cannot admit a non-twisting $`(w=0)`$ null Killing vector except if $`\mu +p=0`$. The algebraically special perfect fluid solutions with twisting null Killing vectors are treated by Wainwright and they admit an Abelian group $`G_2`$.
This paper will deal with space-times admitting a 3-dimensional Lie group of conformal motions $`C_3`$ acting on null orbits. In the beginning one could get the feeling that this kind of space-times would be empty for perfect fluid solutions, since the line element of these space-times is, by the theorem of Defrise-Carter , conformally related to one admitting a $`G_3`$ acting on null orbits and these ones, as we have pointed out above, do not admit perfect fluid solutions. But, as we will show, this is not the case, since indeed a conformal scaling does change the Ricci tensor, but there are just a few solutions.
## 2 Space-times admitting CKVs acting on null orbits
We shall concern ourselves with space-times $`(M,g)`$ that admit a three-parameter conformal group $`C_3`$ containing an Abelian two-parameter subgroup of isometries $`G_2`$, whose orbits $`S_2`$ are spacelike, diffeomorphic to $`\mathrm{I}\mathrm{R}^2`$ and admit orthogonal two-surfaces; furthermore, we shall assume that the $`C_3`$ acts transitively on null orbits $`N_3`$.
The classification of all possible Lie algebra structures for $`๐_3`$ under the previous hypothesis was given in where coordinates were adapted so that the line element associated with the metric $`g`$ can be written as
$$ds^2=e^{2F}\{dt^2+dr^2+Q[H^1(dy+Wdz)^2+Hdz^2]\},$$
(1)
where $`F`$, $`Q`$, $`H`$ and $`W`$ are all functions of $`t`$ and $`r`$ alone.
If the conformal algebra $`๐_3`$ belongs to the family (A), it was shown in that, for null conformal orbits, one can always bring $`X`$ to the form
$$X=_t+_r+X^y(y,z)_y+X^z(y,z)_z,$$
(2)
where $`X^y(y,z)`$ and $`X^z(y,z)`$ are linear functions of their arguments to be determined from the commutation relations of $`X`$ with the Killing vectors. Specializing now the conformal equations to the CKV (2) and the metric (1), for each possible case, one has the following forms for $`X`$ and the metric functions $`F`$, $`Q`$, $`H`$, and $`W`$ appearing in (1)
$`(I)`$ $`Q=q(tr),H=h(tr),W=w(tr),`$ (3)
$`X=_t+_r.`$
$`(II)`$ $`Q=q(tr),H=h(tr),W=w(tr){\displaystyle \frac{t+r}{2}},`$ (4)
$`X=_t+_r+z_y.`$
$`(III)`$ $`Q=e^{\frac{t+r}{2}}q(tr),H=e^{\frac{t+r}{2}}h(tr),W=e^{\frac{t+r}{2}}w(tr),`$ (5)
$`X=_t+_r+y_y.`$
$`(IV)`$ $`Q=e^{(t+r)}q(tr),H=h(tr),W=w(tr){\displaystyle \frac{t+r}{2}},`$ (6)
$`X=_t+_r+(y+z)_y+z_z.`$
$`(V)`$ $`Q=e^{(t+r)}q(tr),H=h(tr),W=w(tr),`$ (7)
$`X=_t+_r+y_y+z_z.`$
$`(VI)`$ $`Q=e^{(1+p)\frac{t+r}{2}}q(tr),H=e^{(1p)\frac{t+r}{2}}h(tr),W=e^{(1p)\frac{t+r}{2}}w(tr),`$ (8)
$`X=_t+_r+y_y+pz_z(p0,1).`$
$`(VII)`$ $`Q=e^{p\frac{t+r}{2}}q(tr),c=c(tr),g=g(tr),`$ (9)
$`H={\displaystyle \frac{\frac{\sqrt{4p^2}}{2}}{\sqrt{1+c^2+g^2}+c\mathrm{cos}(\sqrt{4p^2}\frac{t+r}{2})+g\mathrm{sin}(\sqrt{4p^2}\frac{t+r}{2})}},`$
$`W={\displaystyle \frac{p}{2}}+{\displaystyle \frac{\frac{\sqrt{4p^2}}{2}[c\mathrm{sin}(\sqrt{4p^2}\frac{t+r}{2})g\mathrm{cos}(\sqrt{4p^2}\frac{t+r}{2})]}{\sqrt{1+c^2+g^2}+c\mathrm{cos}(\sqrt{4p^2}\frac{t+r}{2})+g\mathrm{sin}(\sqrt{4p^2}\frac{t+r}{2})}},`$
$`X=_t+_rz_y+(y+pz)_z(p^2<4).`$
In all of these cases $`F=F(t,r)`$ and the conformal factor $`\mathrm{\Psi }`$ is given by
$$\mathrm{\Psi }=F_{,t}+F_{,r}.$$
(10)
Furthermore one can prove that family (B) cannot admit conformal Killing vectors acting on null orbits (the proof can be found in ).
Note that these results are completely independent of the Einstein field equations and therefore of the assumed energy-momentum tensor.
## 3 Perfect fluid solutions
For perfect fluid solutions the study is exhausted. For a maximal $`C_3`$, with a proper CKV, all possible solutions have been found (see for details). They correspond only to the types $`III`$ and $`VI`$, although none of them satisfies the weak and dominant energy conditions over the whole space-time manifold.
Type $`VI`$ (this includes the type $`III`$ for $`p=0`$)
We make the coordinate transformation $`u=t+r`$ and $`v=tr`$, so that we have $`h=h(v)`$ and $`q=q(v)`$. The field equations yield
$$W=0,$$
(11)
$$F=f(x)+\frac{1}{2}\frac{1+p}{1p}\mathrm{ln}h\frac{1}{2}\mathrm{ln}q,xu\frac{2}{1p}\mathrm{ln}h,$$
(12)
$$0=\left\{\frac{q_{,v}h_{,v}}{qh}+\frac{h_{,vv}}{h}\right\}\mathrm{\Sigma }_0+\left(\frac{h_{,v}}{h}\right)^2\mathrm{\Sigma }_1,$$
(13)
where
$`\mathrm{\Sigma }_0`$ $``$ $`1+p^4+4f_{,x}4pf_{,x}+4p^2f_{,x}4p^3f_{,x}+8f_{,x}^28p^2f_{,x}^2`$ (14)
$``$ $`32f_{,x}^3+32pf_{,x}^38f_{,xx}+8p^2f_{,xx}+32f_{,xx}f_{,x}32pf_{,xx}f_{,x}`$
$`\mathrm{\Sigma }_1`$ $``$ $`2+2p+2p^2+2p^316f_{,x}8pf_{,x}16p^2f_{,x}8p^3f_{,x}+32f_{,x}^2+16pf_{,x}^2`$ (15)
$`+`$ $`48p^2f_{,x}^264pf_{,x}^316f_{,xx}+16pf_{,xx}32pf_{,xx}+64pf_{,xx}f_{,x}.`$
$`h_{,v}=0`$ is excluded since the solution does not correspond to a perfect fluid. Therefore, two possibilities arise:
$`\mathrm{i})`$ $`\mathrm{\Sigma }_0=0,\mathrm{\Sigma }_1=0`$
$`\mathrm{ii})`$ $`{\displaystyle \frac{q_{,v}h_{,v}}{qh}}+{\displaystyle \frac{h_{,vv}}{h}}=a\left({\displaystyle \frac{h_{,v}}{h}}\right)^2(a=\mathrm{const}).`$
In the first case $`f_{,x}`$ must be a constant, and therefore the CKV is not proper. For the second case we have
$$\frac{q_{,v}}{q}=a\frac{h_{,v}}{h}\frac{h_{,vv}}{h_{,v}},$$
(16)
which can be integrated to give
$$q=\frac{h^a}{h_{,v}},$$
(17)
and equation (13) reduces to:
$$1=\frac{f_{,xx}[f_{,x}32(apa2p)+8(2p^2a2p+4p^2+a)]}{[4f_{,x}p1][f_{,x}^28(apa2p)+f_{,x}8(p^2+1)+aap+ap^2ap^322p^2]}.$$
(18)
It is convenient to divide the analysis into three sub-cases.
Sub-case (a): $`a=2p/(p1)`$.
Equation (18) can be readily integrated to give
$$f=\frac{p+1}{4}x\frac{(1p)^2}{p^2+1}\frac{1}{2}\mathrm{ln}|x|+c,c=\mathrm{const}.$$
(19)
We notice that for $`p=1`$ there exists a third Killing vector of the form
$$\zeta =\left(\frac{1}{2}+\frac{1}{2}\frac{h}{h_{,v}}\right)_t+\left(\frac{1}{2}\frac{1}{2}\frac{h}{h_{,v}}\right)_r+y_yz_z.$$
(20)
Sub-case (b): $`a=2/(1p`$).
When $`p=1`$ the solution is a particular case of sub-case (a). The remaining cases may now be integrated giving:
$$f=\mathrm{ln}|1e^{(1+p)x/4}|+c,c=\mathrm{const}.$$
(21)
We note that this sub-case admits the further Killing vector
$$\zeta =\left(\frac{1}{2}+\frac{1p}{4}\frac{h}{h_{,v}}\right)_t+\left(\frac{1}{2}\frac{1p}{4}\frac{h}{h_{,v}}\right)_r+\frac{1p}{2}y_y\frac{1p}{2}z_z,$$
(22)
which violates our requirement of a maximal three-dimensional conformal group $`C_3`$.
Sub-case (c): we finally consider the possibility $`a2p/(p1)`$ and $`a2/(1p)`$.
The solution of (18) is then given implicitly by
$$x=\gamma _1\mathrm{ln}|f_{,x}\beta _0|+\gamma _2\mathrm{ln}|f_{,x}\beta _+|+\gamma _3\mathrm{ln}|f_{,x}\beta _{}|$$
(23)
where
$$\beta _0=\frac{p+1}{4},\beta _\pm =\frac{2(p^2+1)\pm \sqrt{2(p^2+1)(1p)^2(a^22a+2)}}{4(apa2p)}.$$
(24)
$`\gamma _i`$, $`i=1,2,3`$ being constants.
A careful analysis of the energy conditions shows that for all cases (i.e., for all values of the parameters $`a`$ and $`p`$) the solutions can only satisfy the energy conditions over certain open domains of the manifold (see ).
Acknowledgements. This work has been supported by DGICYT Research Project No. PB94-1177.
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# From Special Lagrangian to Hermitian-Yang-Mills via Fourier-Mukai Transform
## 1 Introduction
In this note we argue that a real version of the Fourier-Mukai transform would carry supersymmetric A-cycles to B-cycles. Roughly, special Lagrangian submanifolds (plus local systems) will be mapped to holomorphic submanifolds (plus bundles over them). The approach is similar to the one in , though our emphasis is geometric, as we focus on the differential equations defining D-branes.
Our goal is to provide the basis for a geometric functor relating the categories of D-branes on opposite mirror sides: the derived category of coherent sheaves on one hand, and Fukayaโs category of Lagrangian submanifolds with local systems on the other. (In fact, we consider special Lagrangians as objects.<sup>1</sup><sup>1</sup>1One canโt yet say which formulation of these categories will best fit the physics. The Donaldson-Uhlenbeck-Yau theorem relates stable bundles to solutions of the Hermitian-Yang-Mills equations. Analogously, one hopes that Lagrangians are equivalent to special Lagrangians up to Hamiltonian deformations. We are a long way away from proving such theorems, however. The recent preprint of R. Thomas investigates these two issues as moment map problems related by mirror symmetry. ) At this point, providing a physical interpretation of the derived category is premature โ even though recent discussions of the role of the brane-anti-brane tachyon offer glimpses โ so we content ourselves with thinking of D-branes as vector bundles over holomorphic submanifolds, where โsubmanifoldโ may mean the entire space.
We assume, following , that the mirror manifold has a description as a dual torus fibration, blithely ignoring singular fibers for now. The procedure is now quite simple : a point on a torus determines a line bundle on a dual torus. A section of a torus fibration then determines a family of line bundles on the mirror side. These can fit together to define a bundle. We will show how, in this idealized situation, the differential equations on the two sides are related under this transformation. We also show that the Chern-Simons action of an A-cycle equals the holomorphic Chern-Simons action of its transform, even off-shell. Vafaโs version of mirror symmetry with bundles is then verified in this setting.
It will be interesting, though perhaps quite formidable, to generalize this procedure by relaxing some of our assumptions and extending the setting to include singular fibers and more general objects in the derived category.
Acknowledgements: We gratefully acknowledge helpful conversations with Richard Thomas, who has also obtained similar structures for moduli spaces of supersymmetric cycles in his recent preprint (which he has kindly shared with us). N.C. Leung is supported by NSF grant DMS-9803616 and S.-T. Yau is supported by NSF grant DMS-9803347.
## 2 Supersymmetric A- and B-cycles
The authors of consider the supersymmetric p-brane action and determine the conditions for preserving supersymmetry (BPS).<sup>2</sup><sup>2</sup>2In this section, we follow what has become standard notation. Let us remark, however, that the IIA string theory, on which one normally considers the A-model, naturally has even-dimensional (hence type-B) branes, so-named because their compositions depend on the complex structure. IIB string theory has odd-dimensional (type-A) branes, whose composition depends only on the symplectic structure. They show that there are two kinds of supersymmetric cycles $`(C,L)`$ on a Calabi-Yau threefold $`M`$ where $`C`$ is a (possibly singular and with multiplicity) submanifold of $`M`$ and $`L`$ is a complex line bundle over $`C`$ together with a $`U\left(1\right)`$ connection $`D_A`$. Let us denote the Kรคhler form (resp. holomorphic volume form) on the Calabi-Yau threefold by $`\omega `$ (resp. $`\mathrm{\Omega }`$).
The type-A supersymmetric cycle is when $`C`$ is a special Lagrangian submanifold of $`M`$ and the curvature $`F_A`$ of $`D_A`$ vanishes,
$$F_A=0,$$
namely $`D_A`$ is a unitary flat connection. In the presence of a background $`B`$-field (an element of $`H^2(M,๐/๐)`$), $`F_A`$ should be replaced by $`F_AB,`$ where $`B`$ is understood to be pulled back to the submanifold. We take $`B=0`$ in this paper. Recall that a Lagrangian submanifold $`C`$ is called โspecialโ if when restricting to $`C`$ we have
$$\mathrm{I}m\mathrm{\Omega }=\mathrm{tan}\theta \mathrm{R}e\mathrm{\Omega },$$
for some constant $`\theta .`$ Or equivalently, $`\mathrm{I}me^{i\theta }\mathrm{\Omega }=0`$.
The type-B cycle is when $`C`$ is a complex submanifold of $`M`$ of dimension $`n`$ and the curvature two form $`F_A`$ of $`D_A`$ satisfies following conditions:
$`F_A^{0,2}`$ $`=`$ $`0,`$
$`\mathrm{I}me^{i\theta }\left(\omega +F_A\right)^n`$ $`=`$ $`0.`$
The first equation says that the (0,1) component of the connection determines a holomorphic structure on $`L`$. The second equation is called the deformed Hermitian-Yang-Mills equation and it is equivalent to the following equation,
$$\mathrm{I}m\left(\omega +F_A\right)^n=\mathrm{tan}\theta \mathrm{R}e\left(\omega +F_A\right)^n.$$
For example when $`C`$ is the whole Calabi-Yau manifold $`M`$ of dimension three then the second equation says $`F\omega ^2/2F^3/6=\mathrm{tan}\theta \left[\omega ^3/6\left(F^2/2\right)\omega \right]`$.
## 3 Fourier-Mukai Transform of A- and B-cycles
In this section we explain the Fourier-Mukai transform of supersymmetric cycles. The gist of the story is that, assuming mirror pairs are mirror torus fibrations, each point of a Lagrangian submanifold lies in some fiber โ hence defines a bundle over the dual fiber. When done in families and with connections, we get a bundle with connection on the mirror, and the differential equations defining A-cycles map to those which define B-cycles on the mirror. Recall that the base of the fibration itself โ the zero graph โ should be dual to the six-brane with zero connection. Multi-sections are dual to higher-rank bundles, and are discussed in section 3.2. Other cases appear in 3.3.
We assume that the $`m`$ dimensional Calabi-Yau mirror pair $`M`$ and $`W`$ have dual torus fibrations. To avoid the difficulties of singular fibers and unknown Calabi-Yau metrics, we will only consider a neighborhood of a smooth special Lagrangian torus and also assume the Kรคhler potential $`\varphi `$ on $`M`$ to be $`T^m`$-invariant (see for example p.20 of ). This is the semi-flat assumption of . Notice that the Lagrangian fibrations on $`M`$ and $`W`$ are in fact special.
Therefore, let $`\varphi (x^j,y^j)=\varphi \left(x^j\right)`$. ($`y`$ is the coordinate for the fiber and $`x`$ for the base $`B`$ of the fibration on $`M`$. The holomorphic coordinates on $`M`$ are $`z^j=x^j+iy^j`$โs.) As studied by Calabi, the Ricci tensor vanishes and $`\mathrm{\Omega }=dz^1\mathrm{}.dz^m`$ is covariant constant if and only if $`\varphi `$ satisfies a real Monge-Ampรจre equation
$$det\frac{^2\varphi }{x^ix^j}=const.$$
The Ricci-flat Kรคhler metric and form are
$`g`$ $`=`$ $`{\displaystyle \underset{i,j}{}}{\displaystyle \frac{^2\varphi }{x^ix^j}}\left(dx^idx^j+dy^idy^j\right),`$
$`\omega `$ $`=`$ $`{\displaystyle \frac{i}{2}}{\displaystyle \underset{i,j}{}}{\displaystyle \frac{^2\varphi }{x^ix^j}}dz^id\overline{z}^j\text{,}`$
(henceforth we sum over repeated indices). Notice that $`\mathrm{\Omega }\overline{\mathrm{\Omega }}`$ is a constant mulitple of $`\omega ^m`$ and it is direct consequence of the real Monge-Ampรจre equation.
Also note from the form of the metric $`g`$ that $`M`$ is locally isometric to the tangent bundle of $`B`$ with its metric induced from the metric $`_{i,j}\frac{^2\varphi }{x^ix^j}dx^idx^j`$ on $`B`$. If we use the metric on $`B`$ to identify its tangent bundle with its cotangent bundle, then the above symplectic form $`\omega `$ is just the canonical symplectic form $`dpdq`$ on the cotangent bundle.
We can view the universal cover of $`M`$ either as $`TB`$ with the standard complex structure, or as $`T^{}B`$ with the standard symplectic structure. A solution of the real Monge-Ampรจre equation is used to determine the symplectic structure in the former case and to determine the metric structure, and therefore the complex structure in the latter case.
### 3.1 Transformation of a section
We will construct the transform for a special Lagrangian exhibited as a section of the fibration, i.e. a graph over the base.
Recall that a section of $`T^{}B`$ is Lagrangian with respect to the standard symplectic form if and only if it is a closed one form, and hence locally exact. Therefore (or by calculation), a graph $`y(x)`$ in $`M`$ is Lagrangian with respect to $`\omega `$ if and only if $`\frac{}{x^j}(y^l\varphi _{lk})=\frac{}{x^k}(\varphi _{lj}y^l),`$ where $`\varphi _{ij}=\frac{^2\varphi }{x^ix^j}`$, from which we get
$$y^j=\varphi ^{jk}\frac{f}{x^k}$$
for some function $`f`$ (locally), where $`\varphi ^{jk}`$ is the inverse matrix of $`\varphi _{jk}`$.
Now $`dz^j=dx^j+idy^j`$ and on $`C`$ we have $`dy^j=\varphi ^{jl}\left(\frac{^2f}{x^lx^k}\varphi ^{pq}\varphi _{lkp}\frac{f}{x^q}\right)dx^k`$. Therefore $`dz^j=\left(\delta _{jk}+i\varphi ^{jl}\left(\frac{^2f}{x^lx^k}\varphi ^{pq}\varphi _{lkp}\frac{f}{x^q}\right)\right)dx^k`$ over $`C`$. Notice that if we write $`g=\varphi _{jk}dx^jdx^k`$ as the Riemannian metric on the base, then the Christoffel symbol for the Levi-Civita connection is $`\mathrm{\Gamma }_{lk}^q=\varphi ^{pq}\varphi _{lkp}.`$ Therefore $`Hess\left(f\right)=\left(\frac{^2f}{x^lx^k}\varphi ^{pq}\varphi _{lkp}\frac{f}{x^q}\right)dx^ldx^k.`$ Hence
$`dz^1\mathrm{}dz^m|_C`$ $`=`$ $`det\left(I+ig^1Hess\left(f\right)\right)dx^1\mathrm{}dx^m`$
$`=`$ $`det\left(g\right)^1det\left(g+iHess\left(f\right)\right)dx^1\mathrm{}dx^m,`$
so the special Lagrangian condition (with phase) $`\mathrm{I}m(dz^1\mathrm{}.dz^m)|_C=\mathrm{tan}\theta Re(dz^1\mathrm{}.dz^m)|_C`$ becomes
$$\mathrm{I}mdet\left(g+iHess\left(f\right)\right)=\left(\mathrm{tan}\theta \right)\mathrm{R}edet\left(g+iHess\left(f\right)\right).$$
From these data, we want to construct a $`U\left(1\right)`$ connection over the mirror manifold $`W`$ which satisfies the deformed Hermitian-Yang-Mills equation. The dual manifold $`W`$ is constructed by replacing each torus fiber $`T`$ in $`M`$ by the dual torus $`\stackrel{~}{T}=Hom(T,S^1)`$. If we write the dual coordinates to $`y^1,\mathrm{},y^m`$ as $`\stackrel{~}{y}_1,\mathrm{},\stackrel{~}{y}_m`$, then the dual Riemannian metric on $`W`$ is obtained by taking the dual metric on each dual torus fiber $`\stackrel{~}{T}`$:
$$\stackrel{~}{g}=\underset{i,j}{}\left(\varphi _{ij}dx^idx^j+\varphi ^{ij}d\stackrel{~}{y}_id\stackrel{~}{y}_j\right).$$
We need to understand the complex structure and the symplectic structure on $`W`$ (see for example and ). First we rewrite $`\stackrel{~}{g}`$ as follows,
$$\stackrel{~}{g}=\underset{i,j}{}\varphi ^{ij}\left(\left(\mathrm{\Sigma }_k\varphi _{ik}dx^k\right)\left(\mathrm{\Sigma }_l\varphi _{jl}dx^l\right)+d\stackrel{~}{y}_id\stackrel{~}{y}_j\right).$$
Notice that $`d\left(\mathrm{\Sigma }_k\varphi _{jk}dx^k\right)=0`$ because $`\varphi _{jkl}`$ is symmetric with respect to interchanging the indexes. Therefore there exist functions $`\stackrel{~}{x}_j=\stackrel{~}{x}_j\left(x\right)`$โs such that $`d\stackrel{~}{x}_j=\mathrm{\Sigma }_k\varphi _{jk}dx^k`$ locally โ then $`\frac{\stackrel{~}{x}_j}{x^k}=\varphi _{jk}`$ โ and we obtain
$$\stackrel{~}{g}=\underset{i,j}{}\varphi ^{ij}\left(d\stackrel{~}{x}_id\stackrel{~}{x}_j+d\stackrel{~}{y}_id\stackrel{~}{y}_j\right).$$
So we can use $`\stackrel{~}{z}_j=\stackrel{~}{x}_j+i\stackrel{~}{y}_j`$โs as complex coordinates on $`W`$. It is easy to check that the corresponding symplectic form is given by
$$\stackrel{~}{\omega }=\frac{i}{2}\underset{i,j}{}\varphi ^{ij}d\stackrel{~}{z}_id\overline{\stackrel{~}{z}}_j.$$
Moreover the covariant constant holomorphic m-form on $`W`$ is given by
$$\stackrel{~}{\mathrm{\Omega }}=d\stackrel{~}{z}_1\mathrm{}d\stackrel{~}{z}_m.$$
Again, as a direct consequence of $`\varphi `$ being a solution of the real Monge-Ampรจre equation, $`\stackrel{~}{\mathrm{\Omega }}\overline{\stackrel{~}{\mathrm{\Omega }}}`$ is a constant multiple of $`\stackrel{~}{\omega }^m`$.
###### Remark 1
The mirror manifold $`W`$ can be interpreted as the moduli space of special Lagrangian tori together with flat $`U(1)`$ connections over them (see ). It is because the dual torus parametrizes isomorphism classes of flat $`U(1)`$ connections on the original torus. It can be checked directly that the $`L^2`$ metric, i.e. the Weil-Petersson metric, on this moduli space $`W`$ coincides with our $`\stackrel{~}{g}`$ above.
In general, the relevent metric on the moduli space $`W`$ is given by a two-point function computed via a path integral, wihch includes instanton contributions from holomorphic disks bounding the special Lagrangian torus fibers. However, for our local Calabi-Yau $`M`$ such holomorphic disks do not exist. This is because $`M`$ is homotopic to any one of its fibers; but any such holomorphic disk would define a non-trivial relative homology class. Therefore our metric $`\stackrel{~}{g}`$ coincides with the physical metric on the moduli space $`W`$.
###### Remark 2
We note the symmetry between $`g`$ (resp. $`\omega `$) and $`\stackrel{~}{g}`$ (resp. $`\stackrel{~}{\omega }`$). For one can write $`\varphi ^{ij}`$ as the second derivative of some function $`\stackrel{~}{\varphi }`$ with respect to the $`\stackrel{~}{x}_j`$โs. Simply write $`x^j=x^j\left(\stackrel{~}{x}\right)`$, then $`\frac{x^j}{\stackrel{~}{x}_k}=\varphi ^{jk}=\frac{x^k}{\stackrel{~}{x}_j}`$ and therefore $`x^j=\frac{\mathrm{\Phi }}{\stackrel{~}{x}_j}`$ for some function, $`\mathrm{\Phi },`$ and it is easy to check that $`\stackrel{~}{\varphi }=\mathrm{\Phi }`$.
On each torus fiber, we have canonical isomorphisms $`T=\mathrm{Hom}(\stackrel{~}{T},S^1)=\mathrm{Hom}(\pi _1(\stackrel{~}{T}),S^1),`$ therefore a point $`y=(y^1,\mathrm{},y^m)`$ in $`T`$ defines a flat connection $`D_y`$ on its dual $`\stackrel{~}{T\text{.}}`$ This is the real Fourier-Mukai transform. Explicitly, we have
$$\begin{array}{cc}g_y:& \stackrel{~}{T}i\left(/\right)=S^1\\ & \stackrel{~}{y}i_{j=1}^my^j\stackrel{~}{y}_j,\end{array}$$
and $`D_y=d+A=d+idg_y=d+i\mathrm{\Sigma }y^jd\stackrel{~}{y}_j.`$
In fact we get a torus family of one-forms, since $`y`$ (hence $`A`$) has $`x`$\- (or $`\stackrel{~}{x}`$-) dependence. Namely, we obtain a $`U\left(1\right)`$ connection on $`W`$,
$$D_A=d+i\underset{j}{}y^jd\stackrel{~}{y}_j.$$
Its curvature two form is given by,
$$F_A=dA=\underset{k,j}{}i\frac{y^j}{\stackrel{~}{x}_k}d\stackrel{~}{x}_kd\stackrel{~}{y}_j.$$
In particular
$$F_A^{2,0}=\frac{1}{2}\underset{j,k}{}\left(\frac{y^k}{\stackrel{~}{x}_j}\frac{y^j}{\stackrel{~}{x}_k}\right)d\stackrel{~}{z}_jd\stackrel{~}{z}_k.$$
Therefore, that $`D_A`$ is integrable, i.e. $`F_A^{0,2}=0`$, is equivalent to the existence of $`f=f\left(\stackrel{~}{x}\right)`$ such that $`y^j=\frac{f}{\stackrel{~}{x}_j}=\varphi ^{jk}\frac{f}{x^j}`$ because of $`d\stackrel{~}{x}_j=\mathrm{\Sigma }_k\varphi _{jk}dx^k`$. Namely, the cycle $`CM`$ must be Lagrangian. Now
$$\frac{y^j}{\stackrel{~}{x}_k}=\frac{^2f}{\stackrel{~}{x}_j\stackrel{~}{x}_k}\text{.}$$
In terms of the $`x`$ variable, this is precisely the Hessian of $`f,`$ as discussed above. Therefore the cycle $`CM`$ being special is equivalent to
$$\mathrm{I}m\left(\stackrel{~}{\omega }+F_A\right)^m=\left(\mathrm{tan}\theta \right)Re\left(\stackrel{~}{\omega }+F_A\right)^m.$$
For a general type-A supersymmetric cycle in $`M`$, we have a special Lagrangian $`C`$ in $`M`$ together with a flat $`U\left(1\right)`$ connection on it. Since as before, $`C`$ is expressed as a section of $`\pi :MB`$ and is given by $`y^j=\varphi ^{jk}\frac{f}{x^k}`$, a flat $`U\left(1\right)`$ connection on $`C`$ can be written in the form $`d+ide=d+i\mathrm{\Sigma }\frac{e}{x^k}dx^k`$ for some function $`e=e\left(x\right)`$. Recall that the transformation of $`C`$ alone is the connection $`d+i\mathrm{\Sigma }y^jd\stackrel{~}{y}_j`$ over $`W`$. When the flat connection on $`C`$ is also taken into account, then the total transformation becomes
$`D_A`$ $`=`$ $`d+i\mathrm{\Sigma }y^jd\stackrel{~}{y}_j+ide`$
$`=`$ $`d+i\mathrm{\Sigma }\varphi ^{jk}{\displaystyle \frac{f}{x^k}}d\stackrel{~}{y}_j+i\mathrm{\Sigma }{\displaystyle \frac{e}{\stackrel{~}{x}^j}}d\stackrel{~}{x}_j.`$
Here we have composed the function $`e\left(x\right)`$ with the coordinate transformation $`x=x\left(\stackrel{~}{x}\right).`$ Notice that the added term $`\mathrm{\Sigma }\frac{e}{\stackrel{~}{x}^j}d\stackrel{~}{x}_j`$ is exact and therefore the curvature form of this new connection is the same as the old one. In particular $`D_A`$ satisfies
$`F_A^{0,2}`$ $`=`$ $`0,`$
$`\mathrm{I}me^{i\theta }\left(\omega +F\right)^m`$ $`=`$ $`0,`$
so is a supersymmetric cycle of type-B in $`W`$. By the same reasoning, we can couple with $`C`$ a flat connection on it of any rank and we would still obtain a non-Abelian connection $`D_A`$ on $`W`$ satisfying the above equations.
In conclusion, the transform of a type-A supersymmetric section in $`M`$ is a type-B supersymmetric $`2m`$-cycle in $`W`$.
###### Remark 3
The real Fourier-Mukai transform we discussed above exchanges the symplectic and complex aspects of the two theories.
On the A-cycle side, Donaldson and Hitchin introduce a symplectic form on the space of of maps $`Map(C,M)`$ as follows. If $`\upsilon `$ is a fixed volume form on a three manifold $`C`$, then $`_Cev^{}\omega \upsilon `$ is a symplectic form on $`Map(C,M)`$ and it equips with a Hamiltonian action by the group of volume preserving diffeomorphisms of $`C`$. The zero of the corresponding moment map is precisely the Lagrangian condition on $`fMap(C,M)`$.
If one restricts to the infinite-dimensional complex submanifold of $`Map(C,M)`$ consisting of those $`f`$ which satisfy $`f^{}\mathrm{\Omega }=v`$, then the symplectic quotient is the moduli space of $`A`$-cycles.
On the B-cycle side, we consider the pre-symplectic form $`\mathrm{I}m\left[_W\left(\stackrel{~}{\omega }+๐ฝ\right)^m\right]^{\left[2\right]}`$ on the space of connections $`๐\left(W\right)`$ (see section 4.3 or compare ). This form is preserved by the group of gauge transformations and the corresponding moment map equation is the deformed Hermitian-Yang-Mills equations.
If one restricts to the complex submanifold of $`๐\left(W\right)`$ consisting of those connections which define a holomorphic structure on the bundle, then the symplectic quotient is the moduli space of $`B`$-cycles.
Notice that the above real Fourier-Mukai transform exchanges the moment map condition on one side to the complex condition on the other side. Such exchanges of symplectic and complex aspects are typical in mirror symmetry. We expect this continues to hold true in general and not just for special Lagrangian sections in the semi-flat case.
###### Remark 4
Unlike the Hermitian-Yang-Mills equation, solutions to the fully nonlinear deformed equations may not be elliptic. On the other hand, the special Lagrangian equation is always elliptic since its solutions are calibrated submanifolds. Nevertheless, deformed Hermitian-Yang-Mills connections obtained from the real Fourier-Mukai transformation as above are always elliptic.
### 3.2 Transformation of a multi-section
When $`C`$ is a multi-section of $`\pi :MB,`$ the situation is more complicated. For one thing, holomorphic disks may bound $`C`$ โ a situation which cannot occur in the section case (see section 4.2). Here we propose to look at a line bundle on finite cover of $`W`$ as the transform B-cycle. To begin we assume that $`C`$ is smooth, $`\pi :CB`$ is a branched cover of degree $`r`$ (as in algebraic geometry) and $`_j_k\varphi =\delta _{jk}`$ for simplicity. Away from ramification locus, $`C`$ determines $`r`$ unitary connections on $`W`$ locally, and each satisfies the above equation. One might be tempted to take the diagonal connection on their direct sum, so that this $`U\left(r\right)`$ connection satisfies a non-Abelian analogue of the equation. However, such a connection cannot be defined across ramification locus because of monodromy, which can interchange different summands of the diagonal connection.
In fact, as mentioned in , it is still open (even in physics) to find or derive a non-Abelian analogue of the deformed Hermitian-Yang-Mills equations via string theory (though there are some natural guesses). To remedy this problem, we would instead construct a $`U\left(1\right)`$ connection over a degree $`r`$ cover of $`W`$. This finite cover $`\widehat{\pi }:\widehat{W}W`$ is constructed via the following Cartesian product diagram
$$\begin{array}{ccc}\widehat{W}=C\times _BW& \stackrel{\widehat{\pi }}{}& W\\ & & \\ C& \stackrel{๐}{}& B.\end{array}$$
Notice that $`\widehat{W}`$ is a smooth manifold because $`C`$ is smooth and $`\pi :CB`$ is a branched cover. Moreover, the construction as before determines a smooth unitary connection over $`\widehat{W}`$ satisfying the above equation (with $`\omega `$ replaced by $`\widehat{\pi }^{}\omega `$).
We remark that employing this construction is similar to the use of isogenies needed to define the categorical isomorphism which proves Kontsevichโs conjecture in the case of the elliptic curve . The point is that multi-sections transforming to higher-rank can be handled via single sections giving line bundles, by imposing functoriality after pushing forward under finite covers.
Even though this connection over $`\widehat{W}`$ satisfies $`F^{0,2}=0`$ in a suitable sense, $`\widehat{W}`$ is not a complex manifold.
### 3.3 Transformation for more general cycles
In the previous sections, we considered only sections or multi-sections on $`M`$ and obtain holomorphic bundles on $`W`$ which satisfied the deformed Hermitian-Yang-Mills equation. Notice that, for a Calabi-Yau threefold $`M`$, a multi-section of $`\pi :MB`$ can be characterized as special Lagrangian cycle $`C`$ whose image under $`\pi `$ is of dimension three โ namely, the whole $`B`$. If the image has dimension zero, then the Lagrangian is a torus fiber plus bundle and its dual is the point (0-brane) it represents on the corresponding dual torus. This is the basis for the conjecture of . Here we are going to look at the other cases, that is the dimension of the image of $`C`$ under $`\pi `$ is either (i) one or (ii) two. For simplicity we shall only look at the flat case, namely $`M=T^6=B\times F`$ where both $`B`$ and $`F`$ are flat three dimensional Lagrangian tori.
Case (i), when $`dim\pi \left(C\right)=1`$. The restriction of $`\pi `$ to $`C`$ express $`C`$ as the total space of an one-parameter family of surfaces. In fact we will see that this is a product family of an affine $`T^1`$ in $`B`$ with an affine $`T^2`$ in $`F`$.
As before we denote the coordinates of $`B`$ (resp. $`F`$) by $`x^1,x^2,x^3`$ (resp. $`y^1,y^2,y^3`$). Without loss of generality, we can assume that $`\pi \left(C\right)`$ is locally given by $`x^2=f\left(x^1\right)`$ and $`x^3=g\left(x^1\right)`$. Moreover the surface in $`C`$ over any such point is determined by $`y^3=h(x^1,y^1,y^2)`$. In particular, $`C`$ is parametrized by $`x^1,y^1`$ and $`y^2`$ locally.
The special condition $`\mathrm{I}m\left(dz^1dz^2dz^3\right)|_C=0`$ implies that $`h`$ is independent of $`x^1`$. Namely the surface family $`C`$ is indeed a product subfamily of $`M=F\times B`$. Now the Lagrangian conditions read as follow:
$`1+{\displaystyle \frac{dg}{dx^1}}{\displaystyle \frac{h}{y^1}}`$ $`=`$ $`0,`$
$`{\displaystyle \frac{df}{dx^1}}+{\displaystyle \frac{dg}{dx^1}}{\displaystyle \frac{h}{y^2}}`$ $`=`$ $`0.`$
These imply that,
$`f=x^2`$ $`=`$ $`ax^1+\alpha `$
$`g=x^3`$ $`=`$ $`bx^1+\beta `$
$`h=y^3`$ $`=`$ $`{\displaystyle \frac{1}{b}}y^1{\displaystyle \frac{a}{b}}y^2+{\displaystyle \frac{\gamma }{b}}\text{.}`$
By analyticity of special Lagrangians, these parametrizations hold true on the whole $`C`$. We can therefore express $`C`$ as the product $`C_B\times C_F`$ with $`C_FF`$ being a two torus and $`C_BB`$ being a circle. Here
$`C_B`$ $`=`$ $`\left\{(x^1,x^2,x^3)=(1,a,b)x^1+(0,\alpha ,\beta )\right\},`$
$`C_F`$ $`=`$ $`\{(y^1,y^2,y^3):(y^1,y^2,y^3)(1,a,b)=\gamma \}.`$
Since our primary interest is when $`C`$ is a closed subspace of $`M`$, this implies that both $`a`$ and $`b`$ are rational numbers and $`C`$ is a three torus which sits in $`M=T^6`$ as a totally geodesic flat torus. To describe a supersymmetric cycle, we also need a $`U\left(1\right)`$ flat connection $`D_A`$ over $`C`$. If we parametrize $`C`$ by coordinate functions $`x^1,y^2`$ and $`y^3`$ as above, then we have
$$D_A=d+i\left(\stackrel{~}{\gamma }dx^1+\stackrel{~}{\alpha }dy^2+\stackrel{~}{\beta }dy^3\right)\text{,}$$
for some real numbers $`\stackrel{~}{\alpha },\stackrel{~}{\beta }`$ and $`\stackrel{~}{\gamma }`$.
Now let us define the transformation of $`(C,L)`$. First, the mirror of $`M=B\times F`$ equals $`W=B\times \stackrel{~}{F}`$ where $`\stackrel{~}{F}`$ is the dual three torus to $`F`$. A point $`(x_1,x_2,x_3,\stackrel{~}{y}_1,\stackrel{~}{y}_2,\stackrel{~}{y}_3)=(x,\stackrel{~}{y})W`$ (note $`x=\stackrel{~}{x}`$ here, by flatness of the metric) lies in the mirror $`(\stackrel{~}{C},\stackrel{~}{L})`$ of the SUSY cycle $`(C,L)`$ if and only if the flat connection which is obtained by the restriction of $`D_A`$ to $`x\times C_F`$ twisted by the one form $`i\stackrel{~}{y}_jdy^j`$ is in fact trivial. That is, $`\stackrel{~}{\alpha }dy^2+\stackrel{~}{\beta }dy^3+\stackrel{~}{y}_jdy^j=0.`$ Using the equation $`(y^1,y^2,y^3)(1,a,b)=\gamma `$, we obtain that
$`\stackrel{~}{y}_2`$ $`=`$ $`a\stackrel{~}{y}_1\stackrel{~}{\alpha },`$
$`\stackrel{~}{y}_3`$ $`=`$ $`b\stackrel{~}{y}_1\stackrel{~}{\beta }.`$
Or equivalently, $`\stackrel{~}{C}=C_B\times C_{\stackrel{~}{F}}`$ with $`C_{\stackrel{~}{F}}=\left\{(\stackrel{~}{y}_1,\stackrel{~}{y}_2,\stackrel{~}{y}_3)=(1,a,b)\stackrel{~}{y}_1(0,\stackrel{~}{\alpha },\stackrel{~}{\beta })\right\}`$ and $`C_B=\left\{(x_1,x_2,x_3)=(1,a,b)x_1+(0,\alpha ,\beta )\right\}`$. In particular $`\stackrel{~}{C}`$ is a holomorphic curve in $`W`$. The last step would be to determine the $`U\left(1\right)`$ connection $`D_{\stackrel{~}{A}}`$ on $`\stackrel{~}{C}`$. By essentially the same argument as above and the fact that there is no transformation along the base direction, we obtain
$$D_{\stackrel{~}{A}}=d+i\left(\stackrel{~}{\gamma }dx_1\gamma d\stackrel{~}{y}_1\right).$$
Now the deformed Hermitian-Yang-Mills equation $`\mathrm{I}m\left(\omega +F_{\stackrel{~}{A}}\right)^n=0`$ is equivalent to $`F_{\stackrel{~}{A}}=0,`$ which is obviously true for the above connection $`D_{\stackrel{~}{A}}`$ on $`\stackrel{~}{C}`$.
Case (ii), when $`dim\pi \left(C\right)=2.`$ The restriction of $`\pi `$ to $`C`$ express $`C`$ as the total space of a two-parameter family of circles. As before, let us write the parametrizing surface $`S`$ in $`B`$ as $`x^3=f(x^1,x^2)`$ and the one dimensional fiber over any point of it by $`y^2=g(x^1,x^2,y^1)`$ and $`y^3=h(x^1,x^2,y^1)`$. Namely $`C`$ is parametrized by $`x^1,x^2`$ and $`y^1`$ locally.
Now the Lagrangian condition would imply that each fiber is an affine circle in $`T^3=F`$, among other things. On the other hand, the special condition $`\mathrm{I}m\left(dz^1dz^2dz^3\right)|_C=0`$ implies that the surface $`ST^3=B`$ satisfies a Monge-Ampรจre equation:
$$det\left(^2f\right)=0.$$
We would like to perform a transformation on $`C`$ which would produce a complex surface in $`W`$ together with a holomorphic bundle on it with a Hermitian-Yang-Mills connection. Notice that the deformed Hermitian-Yang-Mills equation in complex dimension two is the same as the Hermitian-Yang-Mills equation. To transform $`C`$ in this case is not as straight forward as before because the equation governing the family of affine circles is more complicated. However the above equation should imply that $`f`$ is an affine function which would then simplify the situation a lot.
## 4 Correspondence of Moduli Spaces
Vafa has argued that the topological open string theory describing strings ending on an A-cycle is equivalent to the topological closed-string model on a Calabi-Yau with a bundle.<sup>3</sup><sup>3</sup>3In order to get a bundle, we consider only sections or multi-sections. Equating the effective string-field theories leads to the conjecture that the ordinary Chern-Simons theory on an A-cycle be equivalent to the holomorphic Chern-Simons theory on the transform B-cycle. Gopakumar and Vafa have verified equality of the partition functions for the dual resolutions of the conifold . Further, all structures on the moduli spaces of branes must be equivalent.
In the previous section, we used the Fourier-Mukai transform in the semi-flat case to identify moduli spaces of A- and B-cycles as a set. In this section, we extend our analysis to the Chern-Simons functional for connections which are not necessarily flat or integrable. Then we study and relate various geometric objects on the moduli spaces related by the transform. In particular, we verify Vafaโs conjecture in this case.
### 4.1 Chern-Simons functionals
In this section, we show the equivalence of the relevant Chern-Simons functionals for corresponding pairs of supersymmetric cycles in our semi-flat case. In fact, we will do this โoff-shell,โ meaning that the equivalence holds even for connections which do not satisfy the flatness or integrability conditions, respectively. The argument is essentially the one given on pp. 4-5 of .
So, instead of flat connections, we consider a general $`U\left(1\right)`$ connection $`d+A`$ on the special Lagrangian section $`C`$ in $`M`$, then the above transform will still produce a connection on $`W`$ which might no longer be integrable. In real dimension three, flat connections of any rank can be characterized as those connections which are critical points of the Chern-Simons functional,
$$CS\left(A\right)=_CTr\left(AdA+\frac{2}{3}A^3\right)\text{.}$$
To be precise, one would need to impose boundary condition or growth condition for $`A`$ because $`C`$ is not a closed manifold.
There is also a complexified version of Chern-Simons for any holomorphic bundle $`E`$ on a Calabi-Yau threefold $`W`$ with holomorphic three form $`\stackrel{~}{\mathrm{\Omega }}`$ ( ). Namely, if $`๐`$ is a Hermitian connection on $`E`$ which might not be integrable, then the holomorphic Chern-Simons functional is given by
$$CS_{hol}\left(๐\right)=_WTr\stackrel{~}{\mathrm{\Omega }}\left(๐\overline{}๐+\frac{2}{3}\left(๐\right)^3\right)\text{.}$$
Notice that $`CS_{hol}\left(๐\right)`$ depends only on the $`(0,1)`$ component of $`๐`$. As in the real case, $`๐`$ is a critical point for the holomorphic Chern-Simons if and only if $`F_๐^{0,2}=0`$, that is an integrable connection.
As argued in , the holomorphic Chern-Simons theory on $`W`$ is conjectured to be mirror to the usual Chern-Simons theory on $`CM,`$ with instanton corrections given by holomorphic disks on $`M`$ with boundary lying on $`C`$ (as we will see in the next section, there are no such instantons in our setting). In fact, we can directly check that the Fourier-Mukai transform not only sends flat connections on $`C`$ to integrable connections on $`W`$, but it preserves the Chern-Simons functional for an arbitrary connection on $`C`$ which is not necessarily flat. Moreover, this holds true for connections of any rank over $`C`$.
We consider $`C=\left\{y^j=\varphi ^{jk}\frac{f}{x^k}\right\}M`$ as in section 3.1, and now $`d+A=d+ie_k\left(x\right)dx^k`$ is an arbitrary rank $`r`$ unitary connection on $`C`$. The real Fourier-Mukai transform of $`C`$ alone (resp. $`C`$ with the above connection) is the connection $`๐_0=d+i\varphi ^{jk}\frac{f}{x^k}d\stackrel{~}{y_j}`$ (resp. $`๐=d+i\left(\varphi ^{jk}\frac{f}{x^k}d\stackrel{~}{y_j}+e_k\varphi ^{jk}d\stackrel{~}{x_j}\right)`$) on $`W`$.
Recall that $`\left(๐_0\right)^{0,1}`$ determines a holomorphic structure on a bundle over $`W,`$ which we use as the background $`\overline{}`$ operator in defining the holomorphic Chern-Simons functional. That is, $`CS_{hol}\left(๐\right)=CS_{hol}(๐,๐_0)`$. In fact if we vary $`\overline{}`$ continuously among holomorphic bundles, the holomorphic Chern-Simons functional remains the same.
Now,
$$CS_{hol}(๐,๐_0)=_WTr\stackrel{~}{\mathrm{\Omega }}\left(\left(\overline{}\frac{1}{2}\varphi ^{jk}\frac{f}{x^k}d\overline{\stackrel{~}{z_j}}\right)+\frac{2}{3}\left(\right)^3\right),$$
where $`=\left(๐๐_0\right)^{0,1}=\frac{i}{2}e_k\varphi ^{jk}d\overline{\stackrel{~}{z_j}}`$. Now
$`{\displaystyle _W}Tr\stackrel{~}{\mathrm{\Omega }}\varphi ^{jk}{\displaystyle \frac{f}{x^k}}๐\overline{\stackrel{~}{z_j}}`$ $`=`$ $`\left(const\right){\displaystyle _W}Tr\stackrel{~}{\mathrm{\Omega }}{\displaystyle \frac{f}{\stackrel{~}{x}^j}}๐\overline{\stackrel{~}{z_j}}`$
$`=`$ $`\left(const\right){\displaystyle Tr\left(\epsilon ^2\right)๐f๐\stackrel{~}{y}_1๐\stackrel{~}{y}_2๐\stackrel{~}{y}_3}`$
$`=`$ $`0.`$
Here $`\epsilon =\frac{i}{2}e_k\varphi ^{jk}d\stackrel{~}{x}_k`$ is a matrix-valued one form, and therefore $`Tr\left(\epsilon ^2\right)=0`$. Using the fact that $`\overline{}=\frac{i}{2}\frac{}{\stackrel{~}{x}^l}\left(e_k\varphi ^{jk}\right)d\overline{\stackrel{~}{z_l}}d\overline{\stackrel{~}{z_j}}`$, we have
$`CS_{hol}(๐,๐_0)`$ $`=`$ $`\left(const\right){\displaystyle _W}Tr\stackrel{~}{\mathrm{\Omega }}\left(\overline{}+{\displaystyle \frac{2}{3}}\left(\right)^3\right)`$
$`=`$ $`\left(const\right){\displaystyle _W}Tr\left(\epsilon d\epsilon +{\displaystyle \frac{2}{3}}\epsilon ^3\right)๐\stackrel{~}{y}_1๐\stackrel{~}{y}_2๐\stackrel{~}{y}_3`$
$`=`$ $`\left(const\right){\displaystyle _x}Tr\left(AdA+{\displaystyle \frac{2}{3}}A^3\right)`$
$`=`$ $`\left(const\right)CS\left(A\right).\mathrm{}`$
When the dimension of $`M`$ is odd but bigger than three, the Fourier-Mukai transform still preserves the (holomorphic) Chern-Simons functional even though their critical points are no longer flat (or integrable) connection. Instead the Euler-Lagrange equation is $`\left(F_A\right)^n=0`$ (or $`\left(F_๐^{0,2}\right)^n=0`$) where $`dim_{}M=m=2n+1`$.
### 4.2 Graded tangent spaces
Transforming A-cycles to B-cycles is only the first step in understanding mirror symmetry with branes . The next step would be to analyze the correspondence between the moduli spaces of cycles (branes). In the next two sections, we identify the graded tangent spaces and the holomorphic $`m`$ forms on the two moduli spaces of supersymmetric cycles. Generally, this would involve holomorphic disk instanton contributions for the A-cycles (analgously to the usual mirror symmetry A-model), but in our simplified setting we now show these are absent.<sup>4</sup><sup>4</sup>4The following argument is a variation of the one given on pp. 25-26 of .
Let $`D`$ be a holomorphic disk whose boundary lies in the special Lagrangian section $`C`$. Since we are in the local case, $`C`$ is homeomorphic to a ball and we can find a closed disk $`D^{}C`$ with $`D=D^{}`$. Now using the assumption that $`C`$ is a section and $`\pi _2\left(T^m\right)=0`$, the closed surface $`DD^{}`$ is contractible in $`M`$. Therefore $`_{DD^{}}\omega =0`$ by Stokes theorem. Now $`_D^{}\omega =0`$ because $`D^{}`$ lies inside a Lagrangian and $`_D\omega >0`$ because it is the area of $`D`$. This is a contradiction. Hence there are no such holomorphic disks on $`M`$.
First we discuss the graded tangent spaces of the moduli of A- and B- cycles. After that we verify that the real Fourier-Mukai transform does preserve them in the semi-flat case.
For the A side, the tangent space to the moduli of special Lagrangians can be identified with the space of closed and co-closed one forms (we called such forms harmonic). This is proved by McLean . We denote it by $`H^1(C,)`$. If $`D_A`$ is a flat $`U\left(r\right)`$ connection on a bundle $`E`$ over $`C`$, then the tangent space of the moduli of such connections at $`D_A`$ can be identified with the space of harmonic one forms with valued in $`ad\left(E\right)`$. We denote it by $`H^1(C,ad\left(E\right))`$. When $`r=1`$ the spaces $`H^1(C,ad\left(E\right))`$ and $`iH^1(C,)`$ are the same. For $`r`$ bigger than one, it is expected that $`H^1(C,ad\left(E\right))`$ is the tangent space at the non-reduced point $`rC`$: If there is a family of special Lagrangians in $`M`$ converging to $`C`$ with multiplicity $`r`$, then it should determine a flat $`U\left(r\right)`$ connection on an open dense set in $`C`$. This connection would extend to the whole $`C`$ if those special Lagrangians in the family are branched covers of $`C`$ in $`T^{}C`$. Then the tangent of this moduli space at $`rC`$ should be $`H^1(C,ad\left(E\right))`$. It is useful to verify this statement.
The graded tangent spaces are defined to be $`_kH^k(C,)`$, or more generally $`_kH^k(C,ad\left(E\right))`$, the space of harmonic $`k`$ forms and those with coefficient in $`ad\left(E\right)`$.
On the B side, a cycle is a $`U\left(r\right)`$ connection $`๐_A`$ on a bundle $``$ over $`W`$ whose curvature $`_A`$ satisfies $`_A^{0,2}=0`$ and $`\mathrm{I}me^{i\theta }\left(\stackrel{~}{\omega }+_A\right)^m=0`$. If we replace the deformed Hermitian-Yang-Mills equation by the non-deformed one, then a tangent vector to this moduli space can be identified with an element $``$ in $`\mathrm{\Omega }^{0,1}(W,ad\left(\right))`$ satisfying $`\overline{}=0`$ and $`\stackrel{~}{\omega }^{m1}=0`$. The second equation is equivalent to $`\overline{}^{}=0`$. That is $``$ is a $`\overline{}`$-harmonic form of type $`(0,1)`$ on $`W`$ with valued in $`ad\left(\right)`$. The space of such $``$ equals the sheaf cohomology $`H^1(W,End\left(E\right))`$ by Dolbeault theorem, provided $`W`$ is compact.
It is not difficult to see that a tangent vector $``$ to the moduli of B-cycles is a deformed $`\overline{}`$-harmonic form in the following sense:<sup>5</sup><sup>5</sup>5If the rank of $`E`$ is bigger than one, then we need to symmetrize the product in the second equation, as done in .
$`\overline{}`$ $`=`$ $`0,`$
$`\mathrm{I}me^{i\theta }\left(\stackrel{~}{\omega }+_A\right)^{m1}`$ $`=`$ $`0.`$
In general a differential form $``$ of type $`(0,q)`$ is called a deformed $`\overline{}`$-harmonic form (compare ) if it satisfies
$`\overline{}`$ $`=`$ $`0,`$
$`\mathrm{I}me^{i\theta }\left(\stackrel{~}{\omega }+_A\right)^{mq}`$ $`=`$ $`0.`$
We denote this space as $`\stackrel{~}{H}^q(W,End\left(\right))`$. When the connection $`๐_A`$ and the phase angle $`\theta `$ are both trivial, a deformed $`\overline{}`$-harmonic form is just an ordinary $`\overline{}`$-harmonic form. It is useful to know if there is always a unique deformed $`\overline{}`$-harmonic representative for each coholomology class in $`H^q(W,End\left(\right))`$. One might want to require that $`\stackrel{~}{\omega }+_A`$ is positive to ensure ellipticity of the equation.
As argued in , mirror symmetry with bundles leads to an identification between $`H^q(C,ad\left(E\right))`$ and $`\stackrel{~}{H}^q(W,End\left(\right))`$ for each $`q.`$<sup>6</sup><sup>6</sup>6In , the author uses $`H^k(W,End\left(\right))`$ instead of $`\stackrel{~}{H}^k(W,End\left(\right))`$. This can be verified in our situation as follows. For simplicity we assume that the phase angle $`\theta `$ is zero and $`E`$ is a line bundle.
First we need to define the transformation from a degree $`q`$ form on $`CM`$ to one on $`W`$. We need one transformation $`\mathrm{\Phi }`$ corresponding to deformations of special Lagrangians and another $`\mathrm{\Psi }`$ corresponding to deformations of its flat unitary bundle.
$`\mathrm{\Omega }^q(C,ad\left(E\right))`$ $``$ $`\mathrm{\Omega }^{0,q}(W,End\left(\right))`$
$`B_1+iB_2`$ $``$ $`\mathrm{\Phi }\left(B_1\right)+i\mathrm{\Psi }\left(B_2\right)`$
If $`B`$ is a $`q`$-form on the section $`CM`$, in the coordinate system of $`x`$โs we write $`B=\mathrm{\Sigma }b_{j_1\mathrm{}j_q}\left(x\right)dx^{j_1}\mathrm{}dx^{j_q}`$. It suffices to define the transformations for $`B=dx^j,`$ by naturality. The first (resp. second) transformation of $`dx^j`$ is given by $`\mathrm{\Phi }\left(B\right)=\left(\varphi ^{jk}d\stackrel{~}{y}_k\right)^{0,1}`$ (resp. $`\mathrm{\Psi }\left(B\right)=\left(dx^j\right)^{0,1}=\left(\varphi ^{jk}d\stackrel{~}{x}_k\right)^{0,1}`$). When $`q`$ equals one, these transformations are compatible with our identification of moduli spaces of A- and B-cycles.
Now let $`B=\mathrm{\Sigma }b_j\left(x\right)dx^j`$ be any one-form on $`CM`$. Then we have $`=\mathrm{\Phi }\left(B\right)=\left(\mathrm{\Sigma }b_j\varphi ^{jk}d\stackrel{~}{y}_k\right)^{0,1}=\frac{i}{2}\mathrm{\Sigma }b_j\left(\stackrel{~}{x}\right)\varphi ^{jk}d\overline{\stackrel{~}{z_k}}`$. Therefore $`\overline{}=\frac{i}{2}\mathrm{\Sigma }\frac{}{\stackrel{~}{x}_l}\left(b_j\left(\stackrel{~}{x}\right)\varphi ^{jk}\right)d\overline{\stackrel{~}{z_l}}d\overline{\stackrel{~}{z_k}}`$ and its vanishing is clearly equivalent to $`dB=0`$ under the coordinate change $`\frac{}{\stackrel{~}{x}_l}=\varphi ^{lk}\frac{}{x^k}`$. It is also easy to see that this equivalence between $`dB=0`$ and $`\overline{}=0`$ holds true for any degree $`q`$ form too.
Our main task is to show that $`d^{}B=0`$ if and only if $`\mathrm{I}m\left(\stackrel{~}{\omega }+_A\right)^{mq}=0`$ for any degree $`q`$ form $`B`$ on $`CM`$ and $`=\mathrm{\Phi }\left(B\right)`$. Notice that, by type considerations, the latter condition is the same as $`\mathrm{I}md\left[\left(\stackrel{~}{\omega }+_A\right)^{mq}\right]=0`$.
Let $`B`$ be any degree $`q`$ form on the special Lagrangian $`C`$. Using the symplectic form $`\omega `$, we obtain a $`q`$-vector field $`v_B,`$ i.e. a section of $`\mathrm{\Lambda }^q\left(T_M\right)`$. Then using arguments as in or , we have
$$B=\pm \iota _{v_B}\mathrm{I}m\mathrm{\Omega }\text{.}$$
Therefore $`d^{}B=0`$ if and only if $`\mathrm{I}md\left(\iota _{v_B}\mathrm{\Omega }\right)=0`$ on $`C`$. If we write $`B=\mathrm{\Sigma }b_{j_1\mathrm{}j_q}\left(x\right)dx^{j_1}\mathrm{}dx^{j_q}`$ then $`v_B=\mathrm{\Sigma }b_{j_1\mathrm{}j_q}\varphi ^{j_1k_1}\mathrm{}\varphi ^{j_qk_q}\frac{}{y^{k_1}}\mathrm{}\frac{}{y^{k_q}}`$. So on $`C`$, we have
$`\iota _{v_B}\mathrm{\Omega }`$ $`=`$ $`\iota _{v_B}\left(dx^1+idy^1\right)\mathrm{}\left(dx^m+idy^m\right)`$
$`=`$ $`\pm \mathrm{\Sigma }b_{j_1\mathrm{}j_q}\varphi ^{j_1k_1}\mathrm{}\varphi ^{j_qk_q}i^q{\displaystyle \underset{lk_i}{}}\left(dx^l+idy^l\right)`$
$`=`$ $`\pm \mathrm{\Sigma }b_{j_1\mathrm{}j_q}\varphi ^{j_1k_1}\mathrm{}\varphi ^{j_qk_q}i^q{\displaystyle \underset{lk_i}{}}\left(dx^l+i{\displaystyle \frac{}{x^p}}\left(\varphi ^{lk}{\displaystyle \frac{f}{x^k}}\right)dx^p\right).`$
Now we consider the corresponding $`=\mathrm{\Phi }\left(B\right)`$ over $`W`$. Explicitly we have
$$=i^q\mathrm{\Sigma }b_{j_1\mathrm{}j_q}\varphi ^{j_1k_1}\mathrm{}\varphi ^{j_qk_q}d\overline{\stackrel{~}{z}}_{_{k_1}}\mathrm{}d\overline{\stackrel{~}{z}}_{_{k_q}}.$$
Therefore we consider the form $`\left(\stackrel{~}{\omega }+_A\right)^{mq}`$ of type $`(mq,m)`$ on $`W`$:
$`\left(\stackrel{~}{\omega }+_A\right)^{mq}`$
$`=`$ $`\left(\mathrm{\Sigma }\left(\varphi ^{jk}+i{\displaystyle \frac{^2f}{\stackrel{~}{x}_j\stackrel{~}{x}_k}}\right)d\stackrel{~}{z}_jd\overline{\stackrel{~}{z}}_k\right)^{mq}\left(i^q\mathrm{\Sigma }b_{j_1\mathrm{}j_q}\varphi ^{j_1k_1}\mathrm{}\varphi ^{j_qk_q}d\overline{\stackrel{~}{z}}_{_{k_1}}\mathrm{}d\overline{\stackrel{~}{z}}_{_{k_q}}\right).`$
After the coordinate transformation $`\frac{}{\stackrel{~}{x}_l}=\varphi ^{lk}\frac{}{x^k}`$, it is now easy to see that $`\mathrm{I}md\left[\left(\stackrel{~}{\omega }+_A\right)^{mq}\right]=0`$ if and only if $`\mathrm{I}md\left(\iota _{v_B}\mathrm{\Omega }\right)=0`$. That is, $`B`$ is a co-closed form on $`C`$ of degree $`q`$. So $`\mathrm{\Phi }`$ carries a harmonic $`q`$-form on $`C`$ to a deformed harmonic $`(0,q)`$-form on $`W`$. Now if $`E`$ is a higher rank vector bundle over $`C`$, then the only changes we need in the proof are $`B`$ and $`v_B`$ now have valued in $`ad\left(E\right)`$ and we would replace the exterior differentiation by covariant differentiation and also we need to symmetrize the product in the $`\overline{}`$-harmonic equation. The proof of the equivalence for $`\mathrm{\Psi }`$ is similar, and we omit it.
Therefore, we have proved that $`B_1+iB_2`$ is a harmonic form of degree $`q`$ over $`C`$ if and only if $`\mathrm{\Phi }\left(B_1\right)+i\mathrm{\Psi }\left(B_2\right)`$ is a $`\overline{}`$-harmonic form of degree $`(0,q)`$ over $`W`$. In particular $`\mathrm{\Phi }+i\mathrm{\Psi }`$ maps $`H^q(C,ad\left(E\right))`$ to $`\stackrel{~}{H}^q(W,End\left(\right))`$.
### 4.3 Holomorphic $`m`$-forms on moduli spaces
The moduli spaces of A-cycles and B-cycles on a Calabi-Yau $`m`$-fold have natural holomorphic $`m`$-forms. As explained in , these $`m`$-forms, which inclde holomorphic disk instanton corrections on the A-cycle side, can be identified with physical correlation functions derived from the Chern-Simons partition function. Under Vafaโs version of the mirror conjecture with bundles, these partition functions and correlators should be the same for any mirror pair $`M`$ and $`W`$, at least in dimension three. In this section we recall the definitions of the holomorphic $`m`$-forms and verify this equality our semi-flat case.
First we define a degree-$`m`$ closed form on $`Map(C,M)`$ by $`_Cev^{}\omega ^m`$ where $`\omega `$ is the Kรคhler form on $`M`$ and $`C\times Map(C,M)\stackrel{ev}{}M`$ is the evaluation map. For simplicity we will pretend $`C`$ is a closed manifold, otherwise suitable boundary condition is required. If $`v`$ is a normal vector field along a Lagrangian immersion $`fMap(C,M)`$, then $`v`$ determines an one form $`\eta _v`$ on $`C`$. At $`fMap(C,M)`$ we have $`_Cev^{}\omega ^m(v_1,\mathrm{},v_m)=_C\eta _{v_1}\mathrm{}\eta _{v_m}`$
Next we need to incorporate flat connections on $`C`$ into the picture. We denote $`๐\left(C\right)`$ the affine space of connections on $`C`$. On $`C\times ๐\left(C\right)`$ there is a naturally defined universal connection $`๐ป`$ and curvature $`๐ฝ`$ (see for example ). With respect to the decomposition of two forms on $`C\times ๐\left(C\right)`$ as $`\mathrm{\Omega }^2\left(C\right)+\mathrm{\Omega }^1\left(C\right)\mathrm{\Omega }^1\left(๐\right)+\mathrm{\Omega }^2\left(๐\right),`$ we write $`๐ฝ=๐ฝ^{2,0}+๐ฝ^{1,1}+๐ฝ^{0,2}`$. Then for $`(x,A)C\times ๐\left(C\right)`$ and $`uT_xC`$, $`B\mathrm{\Omega }^1(C,End\left(E\right)),`$ we have $`๐ฝ^{2,0}(x,A)=F_A`$, $`๐ฝ^{1,1}(x,A)(u,B)=B\left(u\right)`$ and $`๐ฝ^{0,2}=0`$. We consider the following complex valued closed $`m`$ form on $`Map(C,M)\times ๐\left(C\right)`$,
$${}_{A}{}^{}\mathrm{\Omega }=_CTr\left(ev^{}\omega +๐ฝ\right)^m\text{.}$$
Since the tangent spaces of the moduli of special Lagrangian and the moduli of flat $`U\left(1\right)`$ connections can both be identified with the space of harmonic one forms<sup>7</sup><sup>7</sup>7If the rank is greater than one, these harmonic forms will take values in the corresponding local system.. A tangent vector of the moduli space $`{}_{A}{}^{}\left(M\right)`$ is a complex harmonic one form $`\eta +i\mu `$. Then $`{}_{A}{}^{}\mathrm{\Omega }`$ is given explicitly as follows
$${}_{A}{}^{}\mathrm{\Omega }(\eta _1+i\mu _1,\mathrm{},\eta _m+i\mu _m)=_CTr\left(\eta _1+i\mu _1\right)\mathrm{}\left(\eta _m+i\mu _m\right).$$
On the $`W`$ side we have universal connection and curvature on the space of connections $`๐\left(W\right)`$ as before and we have the following complex valued closed $`m`$ form on $`๐\left(W\right)`$,
$${}_{B}{}^{}\mathrm{\Omega }=_W\stackrel{~}{\mathrm{\Omega }}Tr๐ฝ^m\text{.}$$
As before, $`{}_{B}{}^{}\mathrm{\Omega }`$ descends to a closed $`m`$ form on $`{}_{B}{}^{}\left(W\right)`$, the moduli space of holomorphic bundles, or equivalently B-cycles, on $`W`$. It is conjectured by Vafa () that under mirror symmetry, these two forms $`{}_{A}{}^{}\mathrm{\Omega }`$ and $`{}_{B}{}^{}\mathrm{\Omega }`$ are equivalent after instanton correction by holomorphic disks.
In our case, where $`M`$ is semi-flat and $`C`$ is a section, there is no holomorphic disk. Also the real Fourier-Mukai transform gives mirror cycle. Now we can verify Vafaโs conjecture in this situation. Namely $`{}_{A}{}^{}\mathrm{\Omega }`$ and $`{}_{B}{}^{}\mathrm{\Omega }`$ are preserved under the real Fourier-Mukai transform.
For a closed one form on $`C`$ which represents an infinitesimal variation of a $`A`$-cycle in $`M`$, we can write it as $`d\eta +id\mu `$ for some function $`\eta `$ and $`\mu `$ in $`x`$ variables. Under the above real Fourier-Mukai transform, the corresponding infinitesimal variation of the mirror $`B`$-cycle is $`\delta ๐=i\left(\frac{\eta }{\stackrel{~}{x}_j}d\stackrel{~}{x_j}+\frac{\mu }{\stackrel{~}{x}_j}d\stackrel{~}{y_j}\right)`$. Therefore its $`(0,1)`$ component is $`\left(\delta ๐\right)^{0,1}=\frac{i}{2}\left(\frac{\eta }{\stackrel{~}{x}_j}+i\frac{\mu }{\stackrel{~}{x}_j}\right)\left(d\stackrel{~}{x_j}id\stackrel{~}{y_j}\right)`$. So
$`{}_{B}{}^{}\mathrm{\Omega }(\delta ๐_1,\mathrm{},\delta ๐_m)`$ $`=`$ $`{\displaystyle _W}\stackrel{~}{\mathrm{\Omega }}Tr๐ฝ^m(\delta ๐_1,\mathrm{},\delta ๐_m)`$
$`=`$ $`{\displaystyle _W}\stackrel{~}{\mathrm{\Omega }}\left[\delta ๐_1\mathrm{}\delta ๐_m\right]_{sym}`$
$`=`$ $`\left(const\right){\displaystyle _W}\mathrm{\Pi }_j\left(d\eta _j+id\mu _j\right)๐\stackrel{~}{y_1}\mathrm{}๐\stackrel{~}{y_m}`$
$`=`$ $`\left(const\right)^{}{\displaystyle _C}\mathrm{\Pi }_j\left(d\eta _j+id\mu _j\right)`$
$`=`$ $`\left(const\right)_A^{}\mathrm{\Omega }(d\eta _1+id\mu _1,\mathrm{},d\eta _m+id\mu _m).`$
Hence we are done. Note that the same argument also work for higher-rank flat unitary bundles over $`C`$.
## 5 B-cycles in $`M`$ are A-cycles in $`T^{}M`$
We now show that a B-cycle in $`M`$ can also be treated as a special Lagrangian cycle in the cotangent bundle $`X=T^{}M\stackrel{๐}{}M`$. We include this observation for its possible relevance to the recovery of โclassicalโ mirror symmetry from the version with branes, as outlined briefly in . The reader is be warned that the $`2n`$-form that we use for the special condition on $`X`$ may not be the most natural one.
Recall that cotangent bundle of $`M`$, or any manifold, carries a natural symplectic form $`\vartheta =\mathrm{\Sigma }dx^kdu_k+\mathrm{\Sigma }dy^kdv_k`$ where $`x`$โs and $`y`$โs are local coordinates on $`M`$ and $`u`$โs and $`v`$โs are the dual coordinates in $`T^{}M`$. Moreover the conormal bundle of submanifold $`C`$ in $`M`$ is a Lagrangian submanifold with respect to this symplectic form on $`X`$. There are a couple other natural closed two-forms on $`X`$: (i) the pullback of Kรคhler form from $`M`$, namely $`\pi ^{}\omega `$ and (ii) the canonical holomorphic symplectic form $`\vartheta _{hol}`$ via the identification between $`T^{}M`$ and $`\left(T^{}M\right)^{1,0}`$. In term of local holomorphic coordinates $`z^1,\mathrm{},z^n`$ on $`M`$ we have
$`\pi ^{}\omega `$ $`=`$ $`i\mathrm{\Sigma }g_{j\overline{k}}dz^jd\overline{z}^k,`$
$`\vartheta _{hol}`$ $`=`$ $`\mathrm{\Sigma }dz^jdw_j\text{.}`$
Here $`z^j=x^j+iy^j`$ and $`w_j=u_j+iv_j`$. We define the $`2n`$ form $`\mathrm{\Theta }`$ on $`X`$ using a combination of $`\pi ^{}\omega `$ and $`\vartheta _{hol}`$:
$$\mathrm{\Theta }=\left(\pi ^{}\omega +\mathrm{I}m\vartheta _{hol}\right)^n\text{.}$$
Notice that $`\mathrm{\Theta }\overline{\mathrm{\Theta }}=\vartheta ^{2n}`$ and the restriction of $`\mathrm{\Theta }`$ to the zero section is a constant multiple of the volume form on $`M`$
A Lagrangian submanifold $`S`$ in $`X`$ is called special Lagrangian if the restriction of $`\mathrm{\Theta }`$ to $`S`$ satisfies $`\mathrm{I}m\mathrm{\Theta }=\mathrm{tan}\theta \mathrm{R}e\mathrm{\Theta }`$ for some phase angle $`\theta `$. Equivalently $`\mathrm{I}me^{i\theta }\mathrm{\Theta }`$ vanishes on $`S`$.
Next we consider a Hermitian line bundle $`L`$ over $`M`$. Let $`D_A`$ be a Hermitian integrable connection on $`L`$, that is $`F_A^{2,0}=0`$. With respect to a holomorphic trivialization of $`L`$, we can write $`D_A=+\overline{}+\varphi `$ locally for some real valued function $`\varphi (z,\overline{z})`$. This determines a Lagrangian submanifold $`S=\{w_j=\frac{\varphi }{z^j}:j=1,\mathrm{},n\}`$ in $`X`$ with respect to $`\vartheta `$. Notice that the definition of $`S`$ depends on the holomorphic trivialization of $`L`$. The restriction $`\vartheta _{hol}`$ to $`S`$ equals
$`\vartheta _{hol}|_S`$ $`=`$ $`\mathrm{\Sigma }dz^jd\left({\displaystyle \frac{\varphi }{z^j}}\right)`$
$`=`$ $`\mathrm{\Sigma }dz^j\left({\displaystyle \frac{^2\varphi }{z^jz^k}}dz^k+{\displaystyle \frac{^2\varphi }{z^j\overline{z}^k}}d\overline{z}^k\right)`$
$`=`$ $`\mathrm{\Sigma }{\displaystyle \frac{^2\varphi }{z^j\overline{z}^k}}dz^jd\overline{z}^k.`$
This form is pure imaginary because $`\varphi `$ is a real valued function. Therefore the restriction of $`\mathrm{\Theta }`$ to $`S`$ equals
$$\mathrm{\Theta }|_S=\left[\mathrm{\Sigma }\left(ig_{j\overline{k}}+\frac{^2\varphi }{z^j\overline{z}^k}\right)dz^jd\overline{z}^k\right]^n=\left(\omega +F\right)^n.$$
Therefore $`S`$ is a special Lagrangian in $`X`$ if and only if $`(L,D_A)`$ satisfies the deformed Hermitian-Yang-Mills equation on $`M`$.
Naichung Conan Leung, School of Mathematics, University of Minnesota, Minneapolis, MN 55455. (LEUNG@MATH.UMN.EDU)
Shing-Tung Yau, Department of Mathematics, Harvard University, Cambridge, MA 02138. (YAU@MATH.NWU.EDU)
Eric Zaslow, Department of Mathematics, Northwestern University, 2033 Sheridan Road, Evanston, IL 60208. (ZASLOW@MATH.NORTHWESTERN.EDU)
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# Extracting Weak Phase Information from ๐ตโ๐โโข๐โ Decays
\[
UdeM-GPP-TH-00-72
IMSc-00/05/17
## Abstract
We describe a new method for extracting weak, CP-violating phase information, with no hadronic uncertainties, from an angular analysis of $`BV_1V_2`$ decays, where $`V_1`$ and $`V_2`$ are vector mesons. The quantity $`\mathrm{sin}^2(2\beta +\gamma )`$ can be cleanly obtained from the study of decays such as $`B_d^0(t)D^\pm \rho ^{}`$, $`D^\pm a_1^{}`$, $`D^0`$$`(`$$`)`$ $`K^0`$$`(`$$`)`$ , etc. Similarly, one can use $`B_s^0(t)D_s^\pm K^{}`$ to extract $`\mathrm{sin}^2\gamma `$. There are no penguin contributions to these decays. It is possible that $`\mathrm{sin}^2(2\beta +\gamma )`$ will be the second function of CP phases, after $`\mathrm{sin}2\beta `$, to be measured at $`B`$-factories.
\]
One of the most important open questions in particle physics is the origin of CP violation. According to the standard model (SM), CP violation is due to the presence of a nonzero complex phase in the Cabibbo-Kobayashi-Maskawa (CKM) quark mixing matrix. This explanation can be tested in the $`B`$ system. By measuring CP-violating rate asymmetries in $`B`$ decays, one can extract $`\alpha `$, $`\beta `$ and $`\gamma `$, the three interior angles of the unitarity triangle . The measured values of these angles may be consistent with SM predictions, or they may indicate the presence of physics beyond the SM. Hopefully, it is this latter scenario which will be realized.
The reason that $`B`$ decays are such a useful tool is that the CP angles can be obtained without hadronic uncertainties. The usual technique is to consider a final state $`f`$ to which both $`B^0`$ and $`\overline{B^0}`$ can decay. Because of $`B^0`$$`\overline{B^0}`$ mixing, CP violation then comes about due to an interference between the amplitudes $`B^0f`$ and $`B^0\overline{B^0}f`$. In the early days of the field, it was thought that the CP angles could be easily measured in $`B_d^0(t)\pi ^+\pi ^{}`$ ($`\alpha `$), $`B_d^0(t)\mathrm{\Psi }K_S`$ ($`\beta `$), and $`B_s^0(t)\rho K_S`$ ($`\gamma `$). However, it soon became clear that things would not be so easy: the presence of penguin amplitudes makes the extraction of $`\alpha `$ from $`B_d^0(t)\pi ^+\pi ^{}`$ quite difficult, and completely spoil the measurement of $`\gamma `$ in $`B_s^0(t)\rho K_S`$. Even in the gold-plated mode $`B_d^0(t)\mathrm{\Psi }K_S`$, penguin contributions limit the precision with which $`\beta `$ can be measured to about 2%. In part because of this, a great deal of work was then done developing new methods to cleanly obtain the CP angles from a wide variety of final states.
One class of final states that was considered consists of two vector mesons, $`V_1V_2`$. Because the final state does not have a well-defined orbital angular momentum, $`V_1V_2`$ cannot be a CP eigenstate. This then implies that, even if both $`B^0`$ and $`\overline{B^0}`$ can decay to the final state $`V_1V_2`$, one cannot extract a CP phase cleanly. However, this situation can be remedied with the help of an angular analysis . By examining the decay products of $`V_1`$ and $`V_2`$, one can measure the various helicity components of the final state. Since each helicity state corresponds to a state of well-defined CP, an angular analysis allows one to use $`BV_1V_2`$ decays to obtain one of the CP phases cleanly. Thus, for example, the angle $`\beta `$ can be extracted from the decay $`B_d^0(t)\mathrm{\Psi }K^{}`$: each helicity state of $`\mathrm{\Psi }K^{}`$ can be treated in the same way as $`\mathrm{\Psi }K_S`$.
In this paper, we show that the angular analysis is more powerful than has previously been realized. Due to the interference between the different helicity states, there are enough independent measurements that one can obtain weak phase information from the decays of $`B^0`$ and $`\overline{B^0}`$ to any common final state $`f`$. Furthermore, contrary to other methods, it is not necessary to measure the branching ratios of both $`B^0f`$ and $`\overline{B^0}f`$. This is important for final states such as $`D^\pm \rho ^{}`$, in which one of the two decay amplitudes is considerably smaller than the other one.
We consider a final state $`f`$, consisting of two vector mesons, to which both $`B^0`$ and $`\overline{B^0}`$ can decay. We assume further that only one weak amplitude contributes to $`B^0f`$ and $`\overline{B^0}f`$. We write the helicity amplitudes as follows:
$`A_\lambda Amp(B^0f)_\lambda `$ $`=`$ $`a_\lambda e^{i\delta _\lambda ^a}e^{i\varphi _a},`$ (1)
$`A_\lambda ^{}Amp(\overline{B^0}f)_\lambda `$ $`=`$ $`b_\lambda e^{i\delta _\lambda ^b}e^{i\varphi _b},`$ (2)
$`\overline{A}_\lambda ^{}Amp(B^0\overline{f})_\lambda `$ $`=`$ $`b_\lambda e^{i\delta _\lambda ^b}e^{i\varphi _b},`$ (3)
$`\overline{A}_\lambda Amp(\overline{B^0}\overline{f})_\lambda `$ $`=`$ $`a_\lambda e^{i\delta _\lambda ^a}e^{i\varphi _a},`$ (4)
where the helicity index $`\lambda `$ takes the values $`\{0,,\}`$. In the above, $`\varphi _{a,b}`$ and $`\delta _\lambda ^{a,b}`$ are the weak and strong phases, respectively.
Using CPT invariance, the total decay amplitudes can be written as
$`๐=Amp(B^0f)=A_0g_0+A_{}g_{}+iA_{}g_{},`$ (5)
$`\overline{๐}=Amp(\overline{B^0}\overline{f})=\overline{A}_0g_0+\overline{A}_{}g_{}i\overline{A}_{}g_{},`$ (6)
$`๐^{}=Amp(\overline{B^0}f)=A_0^{}g_0+A_{}^{}g_{}iA_{}^{}g_{},`$ (7)
$`\overline{๐}^{}=Amp(B^0\overline{f})=\overline{A}_0^{}g_0+\overline{A}_{}^{}g_{}+i\overline{A}_{}^{}g_{},`$ (8)
where the $`g_\lambda `$ are the coefficients of the helicity amplitudes written in the linear polarization basis. The $`g_\lambda `$ depend only on the angles describing the kinematics .
With the above equations, the time-dependent decay rate for $`B^0(t)f`$ can be written as
$`\mathrm{\Gamma }(B^0(t)f)=e^{\mathrm{\Gamma }t}{\displaystyle \underset{\lambda \sigma }{}}`$ $`(\mathrm{\Lambda }_{\lambda \sigma }+\mathrm{\Sigma }_{\lambda \sigma }\mathrm{cos}(\mathrm{\Delta }Mt)`$ (10)
$`\rho _{\lambda \sigma }\mathrm{sin}(\mathrm{\Delta }Mt))g_\lambda g_\sigma .`$
Thus, by performing a time-dependent study and angular analysis of the decay $`B^0(t)f`$, one can measure the observables $`\mathrm{\Lambda }_{\lambda \sigma }`$, $`\mathrm{\Sigma }_{\lambda \sigma }`$ and $`\rho _{\lambda \sigma }`$. In terms of the helicity amplitudes $`A_0,A_{},A_{}`$, these can be expressed as follows:
$`\mathrm{\Lambda }_{\lambda \lambda }={\displaystyle \frac{|A_\lambda |^2+|A_\lambda ^{}|^2}{2}},`$ $`\mathrm{\Sigma }_{\lambda \lambda }={\displaystyle \frac{|A_\lambda |^2|A_\lambda ^{}|^2}{2}},`$ (11)
$`\mathrm{\Lambda }_i=\mathrm{Im}(A_{}A_i^{}A_{}^{}A_{i}^{}{}_{}{}^{}),`$ $`\mathrm{\Lambda }_0=\mathrm{Re}(A_{}A_0^{}+A_{}^{}A_{0}^{}{}_{}{}^{}),`$ (12)
$`\mathrm{\Sigma }_i=\mathrm{Im}(A_{}A_i^{}+A_{}^{}A_{i}^{}{}_{}{}^{}),`$ $`\mathrm{\Sigma }_0=\mathrm{Re}(A_{}A_0^{}A_{}^{}A_{0}^{}{}_{}{}^{}),`$ (13)
$`\rho _i=\mathrm{Re}({\displaystyle \frac{q}{p}}[A_{}^{}A_i^{}+A_i^{}A_{}^{}]),`$ $`\rho _{}=\mathrm{Im}\left({\displaystyle \frac{q}{p}}A_{}^{}A_{}^{}\right),`$ (14)
$`\rho _0=\mathrm{Im}\left({\displaystyle \frac{q}{p}}[A_{}^{}A_0^{}+A_0^{}A_{}^{}]\right),`$ $`\rho _{ii}=\mathrm{Im}\left({\displaystyle \frac{q}{p}}A_i^{}A_i^{}\right),`$ (15)
where $`i=\{0,\}`$. In the above, $`q/p=\mathrm{exp}(2i\varphi _M)`$, where $`\varphi _M`$ is the weak phase present in $`B^0`$$`\overline{B^0}`$ mixing.
Similarly, the decay rate for $`B^0(t)\overline{f}`$ is given by
$`\mathrm{\Gamma }(B^0(t)\overline{f})=e^{\mathrm{\Gamma }t}{\displaystyle \underset{\lambda \sigma }{}}`$ $`(\overline{\mathrm{\Lambda }}_{\lambda \sigma }+\overline{\mathrm{\Sigma }}_{\lambda \sigma }\mathrm{cos}(\mathrm{\Delta }Mt)`$ (17)
$`\overline{\rho }_{\lambda \sigma }\mathrm{sin}(\mathrm{\Delta }Mt))g_\lambda g_\sigma .`$
The expressions for the observables $`\overline{\mathrm{\Lambda }}_{\lambda \sigma }`$, $`\overline{\mathrm{\Sigma }}_{\lambda \sigma }`$ and $`\overline{\rho }_{\lambda \sigma }`$ are similar to those given in Eq. (15), with the replacements $`A_\lambda \overline{A}_\lambda ^{}`$ and $`A_\lambda ^{}\overline{A}_\lambda `$.
With the above expressions for the various amplitudes, we now show how to extract weak phase information using the above measurements. First, we note that
$$\mathrm{\Lambda }_{\lambda \lambda }=\overline{\mathrm{\Lambda }}_{\lambda \lambda }=\frac{(a_\lambda ^2+b_\lambda ^2)}{2},\mathrm{\Sigma }_{\lambda \lambda }=\overline{\mathrm{\Sigma }}_{\lambda \lambda }=\frac{(a_\lambda ^2b_\lambda ^2)}{2}.$$
(18)
Thus, one can determine the magnitudes of the amplitudes appearing in Eqs. (1)โ(4), $`a_\lambda ^2`$ and $`b_\lambda ^2`$. However, we must stress that, in fact, knowledge of $`b_\lambda ^2`$ will not be necessary within our method. This is important since some final states have $`b_\lambda a_\lambda `$, and so the determination of $`b_\lambda ^2`$ would be very difficult.
Next, we have
$`\mathrm{\Lambda }_i=\overline{\mathrm{\Lambda }}_i=b_{}b_i\mathrm{sin}(\delta _{}\delta _i+\mathrm{\Delta }_i)a_{}a_i\mathrm{sin}(\mathrm{\Delta }_i),`$ (19)
$`\mathrm{\Sigma }_i=\overline{\mathrm{\Sigma }}_i=b_{}b_i\mathrm{sin}(\delta _{}\delta _i+\mathrm{\Delta }_i)a_{}a_i\mathrm{sin}(\mathrm{\Delta }_i),`$ (20)
where $`\mathrm{\Delta }_i\delta _{}^a\delta _i^a`$ and $`\delta _\lambda \delta _\lambda ^b\delta _\lambda ^a`$. Using Eq. (20) one can solve for $`a_{}a_i\mathrm{sin}\mathrm{\Delta }_i`$. We will see that this is the only combination needed to cleanly extract weak phase information.
The coefficients of the $`\mathrm{sin}(\mathrm{\Delta }mt)`$ term, which can be obtained in a time-dependent study, can be written as
$$\rho _{\lambda \lambda }=\pm a_\lambda b_\lambda \mathrm{sin}(\varphi +\delta _\lambda ),\overline{\rho }_{\lambda \lambda }=\pm a_\lambda b_\lambda \mathrm{sin}(\varphi \delta _\lambda ),$$
(21)
where the sign on the right hand side is positive for $`\lambda =,0`$ and negative for $`\lambda =`$. In the above, we have defined the CP phase $`\varphi 2\varphi _M+\varphi _b\varphi _a`$. These quantities can be used to determine
$$2b_\lambda \mathrm{cos}\delta _\lambda =\pm \frac{\rho _{\lambda \lambda }+\overline{\rho }_{\lambda \lambda }}{a_\lambda \mathrm{sin}\varphi },2b_\lambda \mathrm{sin}\delta _\lambda =\pm \frac{\rho _{\lambda \lambda }\overline{\rho }_{\lambda \lambda }}{a_\lambda \mathrm{cos}\varphi }.$$
(22)
Similarly, the terms involving interference of different helicities are given as
$`\rho _i`$ $`=`$ $`a_{}b_i\mathrm{cos}(\varphi +\delta _i\mathrm{\Delta }_i)a_ib_{}\mathrm{cos}(\varphi +\delta _{}+\mathrm{\Delta }_i),`$ (23)
$`\overline{\rho }_i`$ $`=`$ $`a_{}b_i\mathrm{cos}(\varphi \delta _i+\mathrm{\Delta }_i)a_ib_{}\mathrm{cos}(\varphi \delta _{}\mathrm{\Delta }_i).`$ (24)
Putting all the above information together, we are now in a position to extract the weak phase $`\varphi `$. Using Eq. (22), the expressions in Eq. (24) can be used to yield
$`\rho _i+\overline{\rho }_i=\mathrm{cot}\varphi a_ia_{}\mathrm{cos}\mathrm{\Delta }_i\left[{\displaystyle \frac{\rho _{ii}+\overline{\rho }_{ii}}{a_i^2}}{\displaystyle \frac{\rho _{}+\overline{\rho }_{}}{a_{}^2}}\right]`$ (25)
$`a_ia_{}\mathrm{sin}\mathrm{\Delta }_i\left[{\displaystyle \frac{\rho _{ii}\overline{\rho }_{ii}}{a_i^2}}+{\displaystyle \frac{\rho _{}\overline{\rho }_{}}{a_{}^2}}\right],`$ (26)
$`\rho _i\overline{\rho }_i=\mathrm{tan}\varphi a_ia_{}\mathrm{cos}\mathrm{\Delta }_i\left[{\displaystyle \frac{\rho _{ii}\overline{\rho }_{ii}}{a_i^2}}{\displaystyle \frac{\rho _{}\overline{\rho }_{}}{a_{}^2}}\right]`$ (27)
$`a_ia_{}\mathrm{sin}\mathrm{\Delta }_i\left[{\displaystyle \frac{\rho _{ii}+\overline{\rho }_{ii}}{a_i^2}}+{\displaystyle \frac{\rho _{}+\overline{\rho }_{}}{a_{}^2}}\right].`$ (28)
Now, we already know most of the quantities in the above two equations: (i) $`\rho _{\lambda \sigma }`$ and $`\overline{\rho }_{\lambda \sigma }`$ are measured quantities, (ii) the $`a_\lambda ^2`$ are determined from the relations in Eq. (18), and (iii) $`a_ia_{}\mathrm{sin}\mathrm{\Delta }_i`$ is obtained from Eq. (20). Thus, the above two equations involve only two unknown quantities โ $`\mathrm{tan}\varphi `$ and $`a_ia_{}\mathrm{cos}\mathrm{\Delta }_i`$ โ and can easily be solved (up to a sign ambiguity in each of these quantities). In this way $`\mathrm{tan}^2\varphi `$ (or, equivalently, $`\mathrm{sin}^2\varphi `$) can be obtained from the angular analysis.
Note that this method relies on the measurement of the interference terms between different helicities. However, we do not actually require that all three helicity components of the amplitude be used. In fact, one can use observables involving any two of largest helicity amplitudes. In the above description, one could have chosen โ$`\mathrm{\hspace{0.17em}0}`$โ instead of โ$``$โ or โ$`0`$โ.
We now turn to specific applications of this method. Consider first the situation in which the final state is a CP eigenstate, $`f=\pm \overline{f}`$. In this case, the parameters of Eqs. (1)โ(4) satisfy $`a_\lambda =b_\lambda `$, $`\delta _\lambda ^a=\delta _\lambda ^b`$ (which implies that $`\delta _\lambda =0`$), and $`\varphi _a=\varphi _b`$ (so that $`\varphi 2\varphi _M+2\varphi _b`$). As described above, $`a_\lambda ^2`$ can be obtained from Eq. 18. But now the measurement of $`\rho _{\lambda \lambda }`$ \[Eq. (21)\] directly yields $`\mathrm{sin}\varphi `$. In fact, this is the conventional way of using the angular analysis to measure the weak phases: each helicity state separately gives clean CP-phase information. Thus, when $`f`$ is a CP eigenstate, nothing is gained by including the interference terms.
Of course, in general, final states that are CP โ eigenstates will all receive penguin contributions at some level. Thus, these states violate our assumption that only one weak amplitude contributes to $`B^0f`$ and $`\overline{B^0}f`$. The only quark-level decays which do not receive penguin contributions are $`\overline{b}\overline{c}u\overline{d},\overline{u}c\overline{d}`$, as well as their Cabibbo-suppressed counterparts, $`\overline{b}\overline{c}u\overline{s},\overline{u}c\overline{s}`$. These are, in fact, the types of decays for which our method is most useful, and we will give meson-level examples of each of these below.
Consider first the decays $`B_d^0/\overline{B_d^0}D^{}\rho ^+,D^+\rho ^{}`$ (which correspond to $`\overline{b}\overline{c}u\overline{d},\overline{u}c\overline{d}`$ at the quark level). In this case we have $`\varphi _M=\beta `$, $`\varphi _a=0`$ and $`\varphi _b=\gamma `$, so that $`\varphi =2\beta \gamma `$. The method described above allows one to extract $`\mathrm{sin}^2(2\beta +\gamma )`$ from an angular analysis of the final state $`D^\pm \rho ^{}`$.
In Ref. , Dunietz pointed out that $`\mathrm{sin}^2(2\beta +\gamma )`$ could, in principal, be obtained from measurements of $`B_d^0(t)D^{}\pi ^\pm `$. He used the method of Ref. , which requires the accurate measurement of the quantity $`\mathrm{\Gamma }(\overline{B_d^0}D^{}\pi ^+)/\mathrm{\Gamma }(B_d^0D^{}\pi ^+)`$. This ratio is essentially $`|V_{ub}V_{cd}^{}/V_{cb}^{}V_{ud}|^24\times 10^4`$. Obviously, it will be very difficult to measure this tiny quantity with any precision, which creates a serious barrier to carrying out Dunietzโs method in practice.
On the other hand, our method does not suffer from this problem. In our notation \[Eqs. (1)โ(4)\], the rate $`\mathrm{\Gamma }(\overline{B_d^0}D^{}\rho ^+)`$ is proportional to $`b_\lambda ^2`$. However, as we have already emphasized in the discussion following Eq. (18), a determination of this quantity is not needed to extract $`\mathrm{sin}^2(2\beta +\gamma )`$ using the angular analysis: none of the observables or combinations required for the analysis are proportional to $`b_\lambda ^2`$. Thus, we avoid the practical problems present in Dunietzโs method.
One disadvantage of the final states $`D^\pm \rho ^{}`$ is that the two decay amplitudes are very different in size (hence the small value of $`b_\lambda `$). This results in a very small CP-violating asymmetry whose size is approximately $`|V_{ub}V_{cd}^{}/V_{cb}^{}V_{ud}|2\%`$. Since the total number of $`B`$โs required to make the measurement is inversely proportional to the square of the asymmetry $`A_f`$, $`N_B1/(BR(B_d^0f)A_f^2)`$, this is a potential problem, even though the branching ratio for the decay $`B_d^0D^{}\rho ^+`$ is quite large, roughly 1%.
One can avoid the problem of a small asymmetry by instead using the Cabibbo-suppressed decays $`B_d^0\overline{D}^0K^0,D^0K^0`$ and $`\overline{B_d^0}D^0\overline{K}^0,\overline{D}^0\overline{K}^0`$ (corresponding to the quark-level decays $`\overline{b}\overline{c}u\overline{s},\overline{u}c\overline{s}`$) . (Here it is assumed that both $`K^0`$ and $`\overline{K}^0`$ decay to the same state $`K_S\pi ^0`$.) In this case the two amplitudes are much more equal in size, leading to a large asymmetry of about $`|V_{ub}V_{cs}^{}/V_{cb}^{}V_{us}|40\%`$. The disadvantage, of course, is that the branching ratios for such Cabibbo-suppressed decays are much smaller than those for $`B_d^0/\overline{B_d^0}D^\pm \rho ^{}`$. We estimate that $`B(B_d^0\overline{D}^0K^0)\lambda ^2B(B_d^0\mathrm{\Psi }K^0)=7\times 10^5`$, which, when combined with $`B(K^0K_S\pi ^0)=1/3`$, yields a net branching ratio of about $`2\times 10^5`$. Even though this branching ratio is quite a bit smaller than that for $`B_d^0D^{}\rho ^+`$, the much larger asymmetry makes up for it. We see that the measurement of $`\mathrm{sin}^2(2\beta +\gamma )`$ using $`B_d^0(t)D^0(\text{})K^0(\text{})`$ requires roughly the same number of $`B`$โs as if $`B_d^0(t)D^\pm \rho ^{}`$ were used.
Of course, this leads to an important question: just how many $`B`$โs are needed for such a measurement? The CLEO collaboration has already performed an angular analysis of $`B_d^0D^{}\rho ^+`$ with a sample of $`197\pm 15`$ events , and has been able to measure some of the interference terms. (Of course, since they do not have an asymmetric collider, it is not possible for them to measure the $`\mathrm{sin}(\mathrm{\Delta }Mt)`$ terms.) In addition, in our method it is necessary to tag the decaying $`B_d^0`$/$`\overline{B_d^0}`$. Taking the tagging efficiency to be about 30% , and using the above values for the branching ratios and asymmetries, we estimate the total number of $`B`$โs required to measure $`\mathrm{sin}^2(2\beta +\gamma )`$ using our method to be roughly $`10^8`$. This number may be reduced if it is possible to combine the various final states ($`D^\pm \rho ^{}`$, $`D^\pm a_1^{}`$, $`D^0`$$`(`$$`)`$ $`K^0`$$`(`$$`)`$ , etc.). We therefore conclude that this measurement will probably be possible at a first-generation $`B`$-factory, though it may take several years of data accumulation.
In fact, the extraction of $`\mathrm{sin}^2(2\beta +\gamma )`$ may well turn out to be the second clean measurement to be made at $`B`$-factories ($`\mathrm{sin}2\beta `$ will clearly be measured first via $`B_d^0(t)\mathrm{\Psi }K_S`$). As discussed above, the angle $`\alpha `$ cannot be obtained cleanly from $`B_d^0(t)\pi ^+\pi ^{}`$ due to the presence of penguin contributions. This difficulty can be resolved with the aid of an isospin analysis , but this technique requires measuring the branching ratio for $`B_d^0\pi ^0\pi ^0`$, which may be quite small. It is also possible to extract $`\alpha `$ with no hadronic uncertainties using a Dalitz-plot analysis of $`B_d^0(t)\pi ^+\pi ^{}\pi ^0`$ decays . Here the idea is to isolate the resonant contributions from intermediate $`\rho \pi `$ states, to which certain isospin relations apply. However, one has to be sure that the non-resonant contributions are well-understood, which requires some theoretical input. In any case, it is estimated that this measurement will take roughly six years to complete. As for the angle $`\gamma `$, the original suggestion for measuring it cleanly involved the decays $`B^\pm D^0K^\pm ,\overline{D^0}K^\pm ,D_{CP}^0K^\pm `$ . However, it was subsequently shown that this type of analysis runs into problems because it is virtually impossible to tag the flavor of the final-state $`D`$-meson , and so one cannot distinguish $`B^\pm D^0K^\pm `$ from $`B^\pm \overline{D^0}K^\pm `$ decays. One can still obtain $`\gamma `$ cleanly by studying decays such as $`B^+(K^+\pi ^{})_DK^+`$ and $`B^+(K^+\rho ^{})_DK^+`$, along with their CP-conjugates, but this requires many more $`B`$โs, so that it is unlikely such measurements can be carried out in the first round of $`B`$-factory experiments. Finally, there has been much work recently looking at the possibilities for extracting $`\gamma `$ from $`B\pi K`$ decays . However, all of these methods use flavor $`SU(3)`$ symmetry, and so rely heavily on theoretical input. In view of all of this, it is thus quite conceivable that the second clean extraction of CP phases at $`B`$ factories will be the measurement of $`\mathrm{sin}^2(2\beta +\gamma )`$ using the method described in this paper.
Note that the measurement of $`\mathrm{sin}^2(2\beta +\gamma )`$ may turn out to be very useful in looking for physics beyond the SM. If new physics is present, it will affect the CP asymmetries principally through its contributions to $`B^0`$$`\overline{B^0}`$ mixing . The most straightforward way of searching for this new physics is to consider two distinct decay modes which, in the SM, probe the same CP angle. A discrepancy between the two values would be clear evidence of physics beyond the SM. For example, the angle $`\gamma `$ can be measured using rate asymmetries in $`B^\pm `$ decays as described above ($`B^\pm DK^\pm `$ ), or in $`B_s^0`$/$`\overline{B_s^0}`$ decays ($`B_s^0(t)D_s^\pm K^{}`$ or $`B_s^0(t)D_s^\pm K^{}`$ \[see below\]). If there is new physics in $`B_s^0`$$`\overline{B_s^0}`$ mixing, with new phases, one will obtain different values of $`\gamma `$ from these two systems. Unfortunately, as argued above, it will be difficult to use $`B^\pm `$ decays to obtain $`\gamma `$, at least in the short term, so that we will not have two independent values of $`\gamma `$ to compare. However, this is where the measurement of $`\mathrm{sin}^2(2\beta +\gamma )`$ will be useful: using the value of $`2\beta `$ as measured in $`B_d^0(t)\mathrm{\Psi }K_S`$, one can obtain $`\gamma `$, up to discrete ambiguities. If none of these values of $`\gamma `$ coincide with those given by the measurement of $`\mathrm{sin}^2\gamma `$ in the $`B_s`$ system, this will be a clear signal of new physics.
Finally, one can also consider $`B_s^0`$ and $`\overline{B_s^0}`$ decays corresponding to the quark-level decays $`\overline{b}\overline{c}u\overline{d},\overline{u}c\overline{d}`$, or $`\overline{b}\overline{c}u\overline{s},\overline{u}c\overline{s}`$. The most promising processes are the Cabibbo-suppressed decay modes $`B_s^0/\overline{B_s^0}D_s^\pm K^{}`$. Here $`\varphi _M=0`$, so that the quantity $`\mathrm{sin}^2\gamma `$ can be extracted from the angular analysis of $`B_s^0(t)D_s^\pm K^{}`$. This is therefore a new method of obtaining the CP phase $`\gamma `$. Note that $`\mathrm{sin}^2\gamma `$ can also be obtained from a measurement of $`B_s^0(t)D_s^\pm K^{}`$ using a different method . The advantage of our method is that the branching ratios are likely to be larger. On the other hand, one must also perform an angular analysis, which is likely to require more $`B`$โs. We therefore conclude that the two methods will probably be of equal difficulty experimentally. Thus, this gives two independent ways of extracting $`\mathrm{sin}^2\gamma `$ from similar final states.
In summary, we have presented a new method of using the angular analysis of $`BV_1V_2`$ decays to extract weak, CP-violating phases with no hadronic uncertainties. Its most useful application involves the quark-level decays $`\overline{b}\overline{c}u\overline{d},\overline{u}c\overline{d}`$, and $`\overline{b}\overline{c}u\overline{s},\overline{u}c\overline{s}`$. We have shown that the quantity $`\mathrm{sin}^2(2\beta +\gamma )`$ can be cleanly obtained from the study of the decays $`B_d^0(t)D^\pm \rho ^{}`$, $`D^\pm a_1^{}`$, $`D^0`$$`(`$$`)`$ $`K^0`$$`(`$$`)`$ , etc. Similarly, $`\mathrm{sin}^2\gamma `$ can be extracted from $`B_s^0(t)D_s^\pm K^{}`$. In all of these cases, there are no penguin contributions to the decays. Finally, we have argued that, due to difficulties with other methods of measuring CP phases, $`\mathrm{sin}^2(2\beta +\gamma )`$ may well be the second clean measurement, after $`\mathrm{sin}2\beta `$, which will be made at $`B`$-factories.
N.S. and R.S. thank D.L. for the hospitality of the Universitรฉ de Montrรฉal, where part of this work was done. The work of D.L. was financially supported by NSERC of Canada.
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# Anomalous tricritical behaviour in the coil-globule transition of a single polymer chain
We investigate a model of self-avoiding walk exhibiting a first-order coil-globule transition. This first-order nature, unravelled through the coexistence of distinct coil and globule populations, has observable consequences on the scaling properties. A thorough analysis of the size dependence of the mean radius of gyration evidences a breakdown of the plain tricritical scaling behaviour. In some regimes, anomalous exponents are observed in the transition region and logarithmic corrections arise along the coexistence curve. PACS numbers: 05.70.Fh, 61.41.+e 64.60.Kw, 64.90.+b
Experimental observations as well as theoretical studies for a specific model suggest that the coil-globule transition of an isolated self-avoiding walk of finite size $`N`$ might exhibit first-order features. A strong evidence is the coexistence in some range of temperatures of distinct coil and globule populations. Its theoretical signature is a bimodal shape of the order parameter distribution $`P_N^{(\beta )}(t)`$, where $`t`$ is some conformational parameter of the single chain, such that the transition is achieved through an exchange of weight between the two peaks as temperature varies. We here investigate the observable consequences of this first-order feature on the scaling behaviour of the chain in the transition region. Although one recovers a second-order transition in the infinite-size limit, one may wonder whether the standard tricritical picture of the $`\mathrm{\Theta }`$-point still applies.
In previous studies, scaling arguments supplemented with Monte Carlo simulations of a self-avoiding chain on a cubic lattice provided an exact analytical expression for the entropic contribution $`P_N^{(\beta =0)}(t)`$ to the distribution $`P_N^{(\beta )}(t)`$ of the order parameter $`t`$, defined as:
$$t=\rho ^{1/(\nu d1)}=\left(N/r^3\right)^{5/4}$$
(1)
where $`r`$ is the radius of gyration of the conformation and $`\rho =N/r^d`$ its mean density (here $`d=3`$ and $`\nu =3/5`$ is the Flory exponent). We then modelled the energy of the isolated chain as $`U=NJt`$. Scaling analysis detailed in previous studies led to introduce rescaled variables:
$$\widehat{t}=tN^{1/n}\mathrm{with}n2$$
(2)
$$\widehat{\tau }=\tau N^\varphi \mathrm{with}\varphi =1\frac{1}{n}\mathrm{and}\tau =1\frac{\theta }{T}$$
(3)
where $`\theta `$ is a rough estimate of the transition temperature. We obtained the following expression for the equilibrium distribution:
$$\widehat{P}_N(\widehat{\tau },\widehat{t})=\frac{h(\widehat{\tau },\widehat{t})e^{A^{}N^{q(11/n)}\widehat{t}^q}}{_c(N,\widehat{\tau })}$$
(4)
where
$$h(\widehat{\tau },\widehat{t})=\widehat{t}^ce^{A\widehat{\tau }\widehat{t}B\widehat{t}^n}$$
(5)
is the scale-invariant contribution. The factor $`_c(N,\widehat{\tau })`$ ensures that $`_0^{\mathrm{}}\widehat{P}_N(\widehat{\tau },\widehat{t})๐\widehat{t}=1`$. Let us emphasize the key role of the factor $`\widehat{t}^c`$: for $`c<1`$, this factor, hence the function $`h(\widehat{\tau },\widehat{t})`$, are not integrable in $`\widehat{t}=0`$. It forbids to take the infinite-size limit in $`\widehat{P}_N(\widehat{\tau },\widehat{t})`$ and to use the infinite-size system as a reference system. This factor $`\widehat{t}^c`$ outweights the coil region and it will induce an observable breakdown of the scale invariance. More generally, we introduce for any real $`z`$:
$$_z(N,\widehat{\tau })=_0^{\mathrm{}}\widehat{t}^ze^{A\widehat{\tau }\widehat{t}B\widehat{t}^n}e^{A^{}N^{q(11/n)}\widehat{t}^q}๐\widehat{t}$$
(6)
An arbitrary moment simply writes:
$$<\widehat{t}^\alpha >=\frac{_{\alpha +c}(N,\widehat{\tau })}{_c(N,\widehat{\tau })}$$
(7)
Due to the factor $`\widehat{t}^z`$, the scaling behaviour of $`_z(N,\widehat{\tau })`$ strongly depends on the sign of $`z+1`$. For $`z+1>0`$, it writes:
$$_z\{\begin{array}{cc}\frac{1}{\widehat{\tau }^{z+1}}\hfill & \text{for}\widehat{\tau }1\hfill \\ & \\ \left(\frac{A\widehat{\tau }}{nB}\right)^{\frac{z}{n1}}e^{(n1)B\left(\frac{A\widehat{\tau }}{nB}\right)^{\frac{n}{n1}}}\hfill & \text{for}\widehat{\tau }<0,|\widehat{\tau }|1\hfill \end{array}$$
(8)
In the case where $`z+1>0`$, $`_z(N,\widehat{\tau })`$ is thus asymptotically scale-invariant, depending only on the rescaled variable $`\widehat{\tau }=N^\varphi \tau `$ but no more on $`N`$.
For $`z+1<0`$, two regimes occur on each side of a borderline $`\widehat{\tau }_{coex}(N,z)`$, as shown on Figure 1:
$$_z\{\begin{array}{cc}N^{(z+1)(11/n)}\hfill & \text{above}\hfill \\ & \\ \left(\frac{A\widehat{\tau }}{nB}\right)^{\frac{z}{n1}}e^{(n1)B\left(\frac{A\widehat{\tau }}{nB}\right)^{\frac{n}{n1}}}\hfill & \text{below}\hfill \end{array}$$
(9)
Here arises the above-mentionned scale-invariance breakdown. It originates in the bimodal shape of $`\widehat{t}^{zc}\widehat{P}_N(\widehat{\tau },\widehat{t})`$: when $`\widehat{\tau }`$ decreases, a โcoilโ peak is quite abruptly replaced by a โglobuleโ peak after a coexistence regime precisely located around the curve $`\widehat{\tau }_{coex}(N,z)`$. Writing that the two peaks of $`\widehat{t}^{zc}\widehat{P}_N(\widehat{\tau },\widehat{t})`$ coexist with equal weights (which requires $`z+1<0`$) yields the following implicit equation for the borderline $`\widehat{\tau }_{coex}(N,z)`$:
$$\varphi (1+z)\mathrm{ln}N=(n1)B\left(\frac{A\widehat{\tau }}{nB}\right)^{\frac{n}{n1}}+\mathrm{ln}๐_z+\frac{z\mathrm{ln}|\widehat{\tau }|}{n1}$$
(10)
where $`๐_z`$ is some numerical constant. $`\widehat{\tau }_{coex}(N,z)`$ could have been obtained straightforwardly by writing the balance between the two expressions of $`_z`$ given in Eq.(9). For large $`N`$, $`\widehat{\tau }_{coex}(N,z)(\mathrm{ln}N)^\varphi `$ for any value $`z<1`$. The borderline $`\widehat{\tau }_{coex}(N,z)`$ thus separates the domain in $`(\widehat{\tau },\mathrm{ln}N)`$-space where the globule peak is overwhelming, leading to a scale-invariant expression for $`_z`$, and the domain where the coil peak is overwhelming. In the latter domain, the behaviour of $`_z`$ is ruled by the coil peak, being controlled only by the size $`N`$ (independent of $`\widehat{\tau }`$ at fixed $`N`$). Note that we are here reasoning on the peaks of $`\widehat{t}^{zc}\widehat{P}_N(\widehat{\tau },\widehat{t})`$ so that the two peaks cannot be exactly identified with the coil and the globule phases, unless $`z=c`$.
In the special instance where $`z=c`$, $`\widehat{\tau }_{coex}(N,c)`$ is then the โtrueโ coexistence line where the coil and globule populations are equally weighted; it writes more explicitly (denoting $`\theta =\theta (\mathrm{})`$):
$$\theta (N)\theta =\frac{\theta \widehat{\tau }_{coex}(N,c)}{N^\varphi }\left(\frac{\mathrm{ln}N}{N}\right)^\varphi $$
(11)
Let us note that the value $`n=2`$ is consistent with the predicted value $`\varphi =1/2`$ for the crossover exponent, since here $`\varphi =11/n`$.
Observable quantities are related to the moments of the distribution $`\widehat{P}_N(\widehat{\tau },\widehat{t})`$. Recalling that $`r=N^{1/d}t^{1/d\nu }=N^{\nu _\theta }\widehat{t}^{1/d\nu }`$ where:
$$\nu _\theta =\frac{1}{d}+\frac{1}{n}\left(\nu \frac{1}{d}\right)$$
(12)
we introduce a generalized mean radius of gyration as:
$$R_\alpha =N^{1/d}<t^\alpha >^{\frac{\nu 1/d}{\alpha }}=N^{\nu _\theta }<\widehat{t}^\alpha >^{\frac{\nu 1/d}{\alpha }}$$
(13)
for any nonzero real $`\alpha `$. Scaling behaviour of the moment $`<\widehat{t}^\alpha >`$ follows from the behaviour of $`_c`$ and $`_{\alpha +c}`$ (see Eq.(7)). Of special interest are the values $`\alpha =1`$, as $`R_1`$ is related to the mean value $`<t>`$, and $`\alpha =2(\nu 1/d)`$, leading to the standard mean radius of gyration $`R_G=<r^2>^{1/2}`$, that can be measured through scattering experiments. For fixed $`\tau <0`$, one recovers the foreseeable globule scaling law $`R_\alpha N^{1/d}`$ (here $`d=3`$) whereas for fixed $`\tau >0`$, one recovers the coil scaling law $`R_\alpha N^\nu `$, whatever $`\alpha `$ is. In between, in the transition region $`|\tau |0`$, $`\widehat{\tau }`$ finite, various scaling regimes are observed, depending on the sign of $`1+c`$ and $`1+c+\alpha `$. The different cases are sketched on Figure 2.
If $`c>1`$ and $`\alpha +c>1`$ (case 1), the ratio $`_{\alpha +c}/_c`$ is scale-invariant; it follows that $`R_\alpha `$ exhibits the standard tricritical behaviour:
$$R_\alpha N^{\nu _\theta }f(\widehat{\tau })=N^{\nu _\theta }f(N^\varphi \tau )$$
(14)
with
$$f(\widehat{\tau })\{\begin{array}{cc}|\widehat{\tau }|^{\left(\frac{1}{n1}\right)\left(\nu \frac{1}{d}\right)}\hfill & \mathrm{for}\widehat{\tau }<0,|\widehat{\tau }|1\hfill \\ & \\ \widehat{\tau }^{(\nu 1/d)}\hfill & \mathrm{for}\widehat{\tau }>0,|\widehat{\tau }|1\hfill \end{array}$$
(15)
For negative and large $`\widehat{\tau }=N^\varphi \tau `$, one recovers the usual scaling behaviour $`R_\alpha N^{1/d}`$; this crossover reflects in the hyperscaling relation:
$$\nu _\theta \varphi \left(\nu \frac{1}{d}\right)\left(\frac{1}{n1}\right)=\frac{1}{d}$$
(16)
whereas for positive and large $`\widehat{\tau }=N^\varphi \tau `$, one recovers the usual scaling behaviour $`R_\alpha N^\nu `$; this crossover reflects in the hyperscaling relation:
$$\nu _\theta +\varphi \left(\nu \frac{1}{d}\right)=\nu $$
(17)
As soon as $`c1`$ or $`\alpha +c1`$ (cases 2 to 5), scale-invariance breakdown and anomalous behaviour are observed. The scaling law for $`R_\alpha `$ still writes:
$$R_\alpha N^{\nu _{eff}}f(\widehat{\tau })=N^{\nu _{eff}}f(N^\varphi \tau )$$
(18)
But the exponent $`\nu _{eff}`$ and the scaling function $`f`$ change when passing across the curves $`\widehat{\tau }_{coex}(N,c)`$ and $`\widehat{\tau }_{coex}(N,\alpha +c)`$. It is clear that $`\varphi =11/n`$ is the crossover exponent. In any cases (2 to 5), both $`_{\alpha +c}`$ and $`_c`$ are scale-invariant in the leftmost region. For $`\widehat{\tau }<0`$ with $`|\widehat{\tau }|`$ enough large, any observable $`R_\alpha `$ thus exhibits the scaling behaviour yet encountered in case 1:
$$R_\alpha N^{\nu _\theta }|\widehat{\tau }|^{\left(\frac{1}{n1}\right)\left(\nu \frac{1}{d}\right)}$$
(19)
For $`|\widehat{\tau }|1`$, one recovers the usual scaling behaviour $`R_\alpha N^{1/d}`$, for any $`\alpha `$, according to the hyperscaling relation given in Eq.(16).
In cases 2 to 5, the key point is the emergence of an intermediate region where either $`_{\alpha +c}`$, either $`_c`$, or both behave as $`\widehat{\tau }`$-independent powers of $`N`$. This modifies the exponent $`\nu _\theta `$ into an anomalous exponent:
$$\nu _\theta ^{}(\alpha )=\nu _\theta \varphi \left(\frac{1+c}{\alpha }\right)\left(\nu \frac{1}{d}\right)$$
(20)
observed for $`\alpha >0`$ or:
$$\nu _\theta ^{\prime \prime }(\alpha )=\nu _\theta +\varphi \left(\frac{1+c}{\alpha }\right)\left(\nu \frac{1}{d}\right)$$
(21)
observed for $`\alpha <0`$. The axis $`\widehat{\tau }=0`$ plays no particular role. In cases 2 and 3, only one curve exists, respectively $`\widehat{\tau }_{coex}(N,\alpha +c)`$ and $`\widehat{\tau }_{coex}(N,c)`$. On its right side (i.e. above it), $`R_\alpha `$ behaves respectively as $`N^{\nu _\theta ^{}}f_2(\widehat{\tau })`$ and $`N^{\nu _\theta ^{\prime \prime }}f_3(\widehat{\tau })`$ where $`f_2`$ and $`f_3`$ are complicated (but independent of $`N`$) functions of $`\widehat{\tau }`$. Both reduce to a power law for $`\widehat{\tau }1`$:
$$f_2(\widehat{\tau })\widehat{\tau }^{\left(\frac{c+1}{\alpha }\right)\left(\nu \frac{1}{d}\right)}f_3(\widehat{\tau })\widehat{\tau }^{\left(\frac{\alpha +c+1}{\alpha }\right)\left(\nu \frac{1}{d}\right)}$$
(22)
Using definitions (20) and (21), we obtain hyperscaling relations:
$$\nu _\theta ^{}+\varphi \left(\frac{\alpha +c+1}{\alpha }\right)\left(\nu \frac{1}{d}\right)=\nu $$
(23)
and
$$\nu _\theta ^{\prime \prime }\varphi \left(\frac{c+1}{\alpha }\right)\left(\nu \frac{1}{d}\right)=\nu $$
(24)
ensuring that the expected scaling behaviour $`R_\alpha N^\nu `$ is recovered as soon as $`\widehat{\tau }1`$. Note that the mean order parameter $`<t>`$ and the associated mean radius of gyration $`R_1`$ belong to case 3 ($`c=1.13`$ and $`\alpha =1`$).
In the last cases (4 and 5), the anomalous exponent is observed in the intermediate but unbounded region located between $`\widehat{\tau }_{coex}(N,c)`$ and $`\widehat{\tau }_{coex}(N,\alpha +c)`$. It is still equal to $`\nu _\theta ^{}(\alpha )`$ for $`\alpha >0`$ and $`\nu _\theta ^{\prime \prime }(\alpha )`$ for $`\alpha <0`$. In these cases, the exponent $`\nu `$ is observed in the rightmost region, even if $`\widehat{\tau }`$ is not large with respect to 1 (even negative). Note that the mean radius of gyration $`<r^2>^{1/2}`$ belongs to case 4 ($`c=1.13`$ and $`\alpha +c=1.66`$).
The remarkable features unravelled here are:
(i) the exchange of weight between the coil phase and the globule phase is far sharper with respect to the variation of $`\widehat{\tau }`$ when it occurs through the coexistence of two peaks than when it occurs through the shift of a simple peak. Accordingly, a sharp crossover is observed in the behaviour of $`_z`$ if $`z+1<0`$, located in a narrow stripe of width $`(\mathrm{log}N)^\varphi `$ around $`\widehat{\tau }_{coex}(N,z)`$, whereas a slow change with a scale-invariant behaviour of $`_z`$ is observed for $`z+1>0`$.
(ii) due to the $`N`$-dependence of the curves $`\widehat{\tau }_{coex}(N,z)`$ for $`z+1<0`$, a size-controlled crossover occurs at the passage across the curves $`\widehat{\tau }_{coex}(N,c)`$ (if $`c<1`$) or $`\widehat{\tau }_{coex}(N,\alpha +c)`$ (if $`\alpha +c<1`$).
(iii) when either $`c`$, either $`\alpha +c`$ or both are lower than $`1`$ (cases 2 to 5 of Figure 2), an intermediate region arises where an anomalous exponent (neither $`\nu `$ nor $`\nu _\theta `$) is observed. Usual tricritical scaling is observed only for $`c>1`$ and $`\alpha +c>1`$ (case 1 of Figure 2).
(iv) for $`c<1`$ (cases 3 to 5 of Figure 2), on the coexistence line $`\widehat{\tau }_{coex}(N,c)`$, the scaling behaviour of $`R_\alpha `$ exhibits multiplicative logarithmic corrections with the same exponent whatever $`\alpha `$ is:
$$R_\alpha =N^{\nu _\theta }(\mathrm{ln}N)^{\frac{1}{2}(\nu 1/d)}=N^{\nu _\theta }(\mathrm{ln}N)^{\frac{2}{15}}$$
(25)
What we here suggest is that the scaling behaviour around the $`\mathrm{\Theta }`$-point might be more complex than the standard tricritical behaviour. The peculiarity of the scaling behaviour associated with the first-order features of the transition is the occurence of two โtheta regimesโ. A first one is observed when the globule peak is overwhelming; it thus involves the exponent $`\nu _\theta `$ arising from the scale invariance of the globule phase and the scaling function is a power law. A second, anomalous one involves an observable-dependent exponent $`\nu _\theta ^{}(\alpha `$) (if $`c<1`$ and $`\alpha >0`$) or $`\nu _\theta ^{\prime \prime }(\alpha )`$ (if $`\alpha +c<1`$ and $`\alpha <0`$) where $`\alpha `$ refers to the observable $`R_\alpha `$; this anomalous exponent arises from a non trivial balance between the scale invariant globule phase and the size-controlled coil phase, which is required to match their incompatible scaling behaviours. Only the detailed analysis of equilibrium distributions like $`P_N^{(\beta )}(t)`$ gives reliable predictions about the actual coil-globule transition experienced by a chain of finite size.
A related result is the irrelevance of standard finite-size scaling approaches to describe the size dependence of the thermal coil-globule transition. Indeed, the standard finite-size analysis used for first-order transitions here failed as the two states merge in the infinite-size limit. On the other hand, an analysis based on the knowledge of the infinite-size second-order transition cannot account for the first-order features here observed in finite size. In a word, standard finite-size scaling approaches are designed to predict the rounding of the transition features in finite size but cannot introduce back the change of nature of the transition observed here. A novel finite-size scaling analysis should thus be designed to interpret experimental or numerical data.
Captions
Figure 1: Scaling behaviour of $`_z`$ for $`z+1<0`$ (here $`z=1.66`$, corresponding to the integral $`_{\alpha +c}`$ involved in the mean radius of gyration $`<r^2>^{1/2}`$). The two scaling regimes are separated by the curve $`\widehat{\tau }_{coex}(N,z)`$, depending on $`z`$ but of similar shape for any value $`z<1`$. The curve is restricted to the domain $`\widehat{\tau }\widehat{\tau }_g3.2`$ (independent of $`N`$) where a well-identified globule state (a globule peak) exists. For $`c<1`$, $`\widehat{\tau }_{coex}(N,z=c)`$ is the coexistence line of the then first-order coil-globule transition.
Figure 2: Critical exponents of $`R_\alpha (N,\widehat{\tau })`$ for the different scaling regimes determined by the signs of $`1+c`$ and $`1+\alpha +c`$. The bold line (cases 3,4 and 5) is the coexistence curve $`\widehat{\tau }_{coex}(N,c)`$ (defined for $`c+1<0`$, here $`c=1.13`$); the thin line (cases 2, 4 and 5) is the curve $`\widehat{\tau }_{coex}(N,\alpha +c)`$ (defined for $`1+\alpha +c<0`$, here $`\alpha +c=1.66`$ which corresponds to the mean radius of gyration $`<r^2>^{1/2})`$.
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# Time symmetry in Rigged Hilbert Spaces
## I Introduction
Asking for the time evolution of a system means solving the Schrรถdinger equation $`i\mathrm{}\frac{}{t}\psi =H\psi `$. Restricted to the case of closed systems - which is done from now on - this equation can formally be solved by
$$\psi (t)=e^{iHt}\psi ,$$
(1)
with $`\mathrm{}=1`$ for the sake of simplicity and $`\psi =\psi (0)`$. The whole dynamics is thus governed by the unitary propagator $`U(t)=e^{iHt}`$ which maps any given โinitial stateโ $`\psi `$ at time $`t=0`$ to its propagated counterpart $`\psi (t)`$ at time $`t`$. Quantum mechanical determinism is expressed by the fact that $`\psi `$ can also be regarded as a โfinal stateโ into which a state $`\psi (t)`$ evolves by satisfying $`U(t)\psi (t)=\psi `$ (see fig. 1).
Since equation (1) is valid for all times $`t`$, giving $`\psi (t)=U(t)\psi `$, one is led to the operator identity
$$U(t)U(t)=\mathrm{๐}.$$
(2)
In other words: The set of evolution operators $`\{U(t)\}`$ forms a unitary group $`๐ฐ`$ with the group property (2), giving any propagator $`U(t)`$ its inverse element $`U(t)=U^{}(t)`$. The physical meaning of this mathematical statement is profound: in quantum mechanics there is no microscopic arrow of time. Any state $`\psi `$ can be propagated to the future as well as to the past, so that there is a complete time symmetry.
In order to explain an arrow of time despite these facts, there are several ideas. A few of them are:
1. Any measurement on the system interrupts the deterministic evolution of the state by applying random projections (plus normalization) instead of unitary operations (the โcollapse of the wave functionโ). This process cannot be reversed.
2. Real physical systems can never be regarded as closed. There is always an interacting environment with many unobservable degrees of freedom, inducing non-unitary time evolution of the observed system.
3. Time symmetry in classical mechanics is only broken with respect to macroscopic properties of a system (e.g. temperature or entropy). Considering macroscopic observables in quantum mechanics, this effect also shows up (coarse graining argument).
All these ideas have been worked out to elaborate physical models with strong experimental support and all of them break time symmetry (see e.g. ). So one could be content with the situation and regard those models as fairly good explanations of an arrow of time, which is indeed an undeniable part of human experience.
A group of scientists around *Ilja Prigogine* (nobel prize in chemistry 1977), referred to as the โBrussels Schoolโ, is not satisfied with the situation and has worked out a quantum mechanical formalism especially applying to unstable systems (i.e. systems involving decay), where time symmetry is claimed to be broken at the microscopic level (see e.g. ).
This fact is expressed by the splitting of the unitary group $`๐ฐ`$ into two semigroups $`๐ฐ_{}`$ and $`๐ฐ_+`$ of propagators for negative and positive times, such that past and future states can no longer be regarded as equivalent counterparts. This leads to a so-called โintrinsic irreversibilityโ, implying a microscopic arrow of time (see fig. 2).
This paper is separated into two different parts. The first part is dedicated to a review of the Rigged Hilbert Space formalism worked out by the Brussels School. It is shown how it works, how it applies to physical systems and what its advantages over standard quantum mechanics are. The second part is devoted to the role of time symmetry within the Rigged Hilbert Space formalism. Apart from its practical advantages, the Brussels School points out that a microscopic arrow of time can be implemented here. It will be shown that such an implementation is possible but not mathematically necessary.
## II The basic concepts
### A Spectral decompositions
Writing down the formal solution (1) of a closed systemโs Schrรถdinger equation is not yet really satisfying. How can expectation values of observables be numerically calculated? To this aim one has to find a spectral decomposition of the systemโs Hamiltonian $`H`$, thus one has to solve the time-independent Schrรถdinger equation $`H\psi =E\psi `$. Since $`H`$ is a self-adjoint operator, its spectral decomposition is complete and orthogonal and may therefore be used as a spectral decomposition of the unity operator $`\mathrm{๐}`$. Furthermore the spectrum of $`H`$ is real and consits in general of a discrete part $`\sigma _d`$ and a continuous part $`\sigma _c`$. In special cases there are other parts, e.g. singular continuous and dense spectra, but these cases are excluded here. With this restriction, using the Dirac bra-ket notation, the spectral decomposition of the unity operator reads in general
$$\mathrm{๐}=\underset{n}{}|E_nE_n|+_{\sigma _c}๐E\rho (E)|EE|,$$
(3)
where $`H|E_n=E_n|E_n`$, $`H|E=E|E`$ and $`\rho (E)`$ is the energy density. As soon as the energy decomposition of the initial state $`|\psi `$, generally given by
$$|\psi =\underset{n}{}\psi _n|E_n+_{\sigma _c}๐E\rho (E)\psi (E)|E,$$
(4)
is known, the propagated state is also known and reads
$`|\psi (t)`$ $`=`$ $`{\displaystyle \underset{n}{}}\psi _ne^{iE_nt}|E_n`$ (6)
$`+{\displaystyle _{\sigma _c}}๐E\rho (E)\psi (E)e^{iEt}|E.`$
The expectation value of an observable $`A`$ can be calculated for any time $`t`$ using
$$<A>(t)=\psi (t)|A|\psi (t).$$
(7)
So the fundamental key to the dynamics of a system is the spectral decomposition of its Hamiltonian. All dynamical properties are monitored by this decomposition and here is where the new formalism enters. There is a class of systems where the Brussels Schoolsโs spectral decomposition of the Hamiltonian differs from the usual decomposition: systems envolving decay, called *unstable systems*.
### B Resolvent techniques
The basic concept of the Brussels Schoolโs formalism is to find a different spectral decomposition of the Hamiltonian. In contrast to the method of *complex scaling* (see ), the Hamiltonian is not modified and remains self-adjoint; it only acquires a different representation involving generalized eigenvectors contained in an extended distribution space. The modified decomposition can most easily be achieved by using the Hamiltonianโs *resolvent*
$$R(z):=\frac{1}{zH}.$$
(8)
The spectrum of $`H`$ is defined by the irregularities of its resolvent. The discrete eigenvalues are the singularities of $`R(z)`$, i.e. those points in the complex plane where the operator $`R`$ is not everywhere defined on the Hilbert space $``$. The continuous eigenvalues are those points where $`R`$ is an unbound (and thus non-continuous) operator on $``$. Anywhere else on the complex plane the resolvent is an analytic function mapping complex numbers onto the set of linear bounded operators. As $`H`$ is self-adjoint, any irregularities of $`R`$ take place on the real axis, where the discrete eigenvalues form the point spectrum and the continuous eigenvalues define a cut in the complex plane (see fig. 3).
It is possible to find an analytic continuation of $`R`$ from both sides of the cut, each leading to a different operator
$$R_\pm (E):=\frac{1}{EH\pm iฯต}=R_{}^{}(E),$$
(9)
with $`ฯต+0`$ understood here and in the following. Note that this continuation is not possible at singularities embedded in the continuous spectrum. Note further that on the real axis outside the continuous spectrum both operators $`R_+`$ and $`R_{}`$ coincide with $`R`$. The two continued resolvents are known as the *retarded* and the *advanced* *Green operator*, respectively.
The resolvent can be used to construct eigenvectors. The residuum of the resolvent at a discrete eigenvalue yields the dyadic product of the corresponding eigenvector:
$$\underset{zE_n\pm iฯต}{lim}(zE_n)R(z)=|E_n^\pm E_n^\pm |.$$
(10)
Note that both dyadic products project onto the same subspace corresponding to $`E_n`$. If the discrete eigenvalue is isolated, i.e. not embedded in the continuous spectrum, both limits from above and below the real axis coincide. In the following we will assume that all discrete eigenvalues are isolated and we therefore have
$$\underset{zE_n}{lim}(zE_n)R(z)=|E_nE_n|.$$
(11)
A discrete eigenvector $`|E_n`$ of $`H`$ can be obtained by applying the residual resolvent to any Hilbert vector $`|\psi `$ being not orthogonal to the eigenspace of $`E_n`$. :
$$|E_n=\frac{1}{\psi _n}\underset{zE_n}{lim}(zE_n)R(z)|\psi ,$$
(12)
where $`\psi _n=E_n|\psi 0`$ is an adequate normalization constant, such that $`E_n|E_m=\delta _{nm}`$.
In analogy to the discrete case the residuum of the resolvent at one of the continuous eigenvalues yields the dyadic product of the corresponding continuous eigenvector:
$$\underset{zE\pm iฯต}{lim}(zE)R(z)=|E^\pm E^\pm |.$$
(13)
Note that both dyadic products are equivalent representations of the $`\delta `$-Operator $`\delta (EH_c)`$, defined by
$$_{\sigma _c}๐E\rho (E)\delta (EH_c)f(E):=f(H_c),$$
(14)
where $`H_c`$ is the continuous part of the total Hamiltonian $`H`$. Since the continuous eigenvalues are not isolated, there are always pairs of eigenvectors. They correspond to retarded and advanced solutions of the continuous eigenvalue equation and are known as *incoming* and *outgoing waves*. Both solutions can equivalently be used in the spectral decomposition and they can be constructed by
$$|E^\pm =\frac{1}{\psi _\pm (E)}\frac{\pm iฯต}{EH\pm iฯต}|\psi ,$$
(15)
where $`\psi _\pm (E)=E^\pm |\psi 0`$ is a suitable normalization constant, such that $`E^\pm |E_{}^{}{}_{}{}^{\pm }=\delta (EE^{})`$.
Using dyadic products, the retarded and advanced Green operators can be written as
$$E\sigma _c:R_\pm (E)=\frac{๐ซ_E}{EH}i\pi |EE|,$$
(16)
where $`|EE|`$ is any of the two dyadic products $`|E^\pm E^\pm |`$ and the principal value distribution $`๐ซ_E`$ is defined by
$$๐E^{}๐ซ_Ef(E^{}):=_{Eฯต}^{E+ฯต}๐E^{}f(E^{}).$$
(17)
The difference of the two Green operators gives the dyadic product of the continuous eigenvector:
$$\frac{i}{2\pi }\left\{R_+(E)R_{}(E)\right\}=|E^\pm E^\pm |.$$
(18)
Altogether, the spectral decomposition of the unity operator may be expressed by the exclusive use of the resolvent:
$`\mathrm{๐}`$ $`=`$ $`{\displaystyle \underset{n}{}}\underset{zE_n}{lim}(zE_n)R(z)`$ (20)
$`+{\displaystyle \frac{i}{2\pi }}{\displaystyle _{\sigma _c}}๐E\rho (E)\left\{R_+(E)R_{}(E)\right\}.`$
This unity decomposition will be used as the starting point of the Brussels Schoolsโs unity spectral decomposition.
### C The second Riemann sheet
The resolvent can be analytically continued to the complex plane beyond the real axis entering the *second Riemann sheet*:
$$R_\pm (z):=\underset{Ez}{lim}R_\pm (E),E\sigma _c.$$
(21)
On the one half of the complex plane where the continuation started (the first Riemann sheet), both operators coincide with the original resolvent. On the other half (the second Riemann sheet) it might happen that new singularities appear, and that in a pairwise manner: If $`z_k`$ is a singularity of $`R_+(z)`$ in the lower half of the complex plane, then $`z_k^{}`$ is a singularity of $`R_{}(z)`$ in the upper half (see fig. 4). Unstable systems are exactly those with a resolvent having complex poles on the second Riemann sheet.
The second Riemann sheet poles act as complex resonances and cause the decay of certain states.
### D An extended spectral decomposition
Consider the continuous part of the unitary decomposition (20):
$$\mathrm{๐}_c:=\frac{i}{2\pi }_{\sigma _c}๐E\rho (E)\left\{R_+(E)R_{}(E)\right\}.$$
(22)
The following manipulations of $`\mathrm{๐}_c`$ are reducing its domain. As a consequence one is forced to use a smaller *test space*, enlarging at the same time the distribution space and inducing this way a certain *Gelfand tripel* or *Rigged Hilbert Space* where the modified completeness relation remains valid. We delay this discussion to section IV.
If one analytically deforms the integration path along the continuous spectrum $`\sigma _c`$ to a path $`\mathrm{\Gamma }_+`$ leading through the lower half of the complex plane, one enters the second Riemann sheet and crosses the singularities of $`R_+(z)`$. Let the set of these singularities be defined by $`\sigma _+=\{z_k\}`$ and be enclosed by $`\sigma _c\mathrm{\Gamma }_+`$. In the course of the deformation circular integrations around the singularities are split off, each leading to a residuum of $`R_+(z)`$ at these points (see fig. 5).
The modified operator now reads
$`\mathrm{๐}_c^+`$ $`=`$ $`{\displaystyle \underset{k}{}}\underset{Ez_k}{lim}(Ez_k)R_+(E)`$ (24)
$`+{\displaystyle \frac{i}{2\pi }}{\displaystyle _{\mathrm{\Gamma }_+}}๐z\rho (z)\left\{R_+(z)R_{}(z)\right\}.`$
Introducing a *complex distribution* $`D_\pm (E)`$, defined by
$$_{\sigma _c}๐ED_\pm (E)f(E):=_{\mathrm{\Gamma }_\pm }๐zf(z),$$
(25)
where $`\mathrm{\Gamma }_{}=\mathrm{\Gamma }_+^{}`$ and $`f(z)`$ being the analytic continuation of $`f`$ to respectively the upper and lower half of the complex plane, one can rewrite $`\mathrm{๐}_c^+`$ as
$`\mathrm{๐}_c^+`$ $`=`$ $`{\displaystyle \underset{k}{}}\underset{Ez_k}{lim}(Ez_k)R_+(E)`$ (27)
$`+{\displaystyle \frac{i}{2\pi }}{\displaystyle _{\sigma _c}}๐E\rho (E)D_+(E)\left\{R_+(E)R_{}(E)\right\}.`$
In analogy to the standard case one defines dyadic products of *generalized eigenvectors* of $`H`$ by
$$|f_k^+f_k^{}|:=\underset{Ez_k}{lim}(Ez_k)R_+(E)$$
(28)
for the discrete case and, using (18),
$$|f^+(E)f^{}(E)|=D_+(E)|EE|$$
(29)
for the continuous case, where $`|EE|`$ is any of the two dyadic products $`|E^\pm E^\pm |`$. Hence we have
$`{\displaystyle _{\sigma _c}}๐E\rho (E)|EE|`$ $`=`$ $`{\displaystyle \underset{k}{}}|f_k^+f_k^{}|`$ (30)
$`+`$ $`{\displaystyle _{\sigma _c}}๐E\rho (E)|f^+(E)f^{}(E)|.`$ (31)
The generalized eigenvectors appear as pairs of right and left eigenvectors. The discrete eigenvectors are called *Gamov vectors* and can be constructed by:
$`|f_k^+`$ $`=`$ $`{\displaystyle \frac{1}{\psi _k}}\underset{Ez_k}{lim}(Ez_k)R_+(E)|\psi ,`$ (32)
$`f_k^{}|`$ $`=`$ $`{\displaystyle \frac{1}{\psi _k}}\underset{Ez_k}{lim}(Ez_k)\psi |R_+(E),`$ (33)
where $`|\psi `$ is a suitable vector and $`\psi _k=f_k^+|\psi 0`$ is a normalization constant, such that $`f_k^{}|f_l^+=\delta _{kl}`$. By adjunction one sees that $`|f_k^+|f_k^{}`$. Furthermore we have $`f_k^\pm |f_k^\pm =f_k^\pm =0`$. These objects are left and right eigenvectors of the Hamiltonian $`H`$ corresponding to complex eigenvalues:
$`H|f_k^+`$ $`=`$ $`z_k|f_k^+,`$ (34)
$`f_k^{}|H`$ $`=`$ $`f_k^{}|z_k.`$ (35)
Since the Gamov vectors have zero norm, they all coincide in Hilbert space with the null vector, which is not allowed to be an eigenvector. Also, since $`H`$ is self-adjoint, it cannot have any eigenvectors corresponding to complex eigenvalues. Concluding, the Gamov vectors cannot be elements of the Hilbert space. In fact they are distributions in a suitably chosen distribution space (see section IV B). The continuous eigenvectors can be constructed in two ways:
$`|f^\pm (E)=D_\pm (E)|E^+`$ $``$ $`f^\pm (E)|=D_{}E^+|,`$ (36)
$`|g^\pm (E)=D_\pm (E)|E^{}`$ $``$ $`g^\pm (E)|=D_{}(E)E^{}|.`$ (37)
By application to test vectors with analytic extensions to the complex plane one sees that
$`D_+(E)|E^{}`$ $`=`$ $`|E^{}`$ (38)
$`D_{}(E)|E^+`$ $`=`$ $`|E^+.`$ (39)
Hence there are two equivalent bases given by
$`B_1^+`$ $`=`$ $`\{|E_nE_n|,|f_k^+f_k^{}|,|f^+(E)E^+|\},`$ (40)
$`B_2^+`$ $`=`$ $`\{|E_nE_n|,|f_k^+f_k^{}|,|E^{}g^{}(E)|\}.`$ (41)
Both bases can equally be used to span the Hilbert space $``$. Arbitrarily we select $`B_1^+`$ and obtain a *retarded* unity decomposition
$`\mathrm{๐}_+`$ $`=`$ $`{\displaystyle \underset{n}{}}|E_nE_n|+{\displaystyle \underset{k}{}}|f_k^+f_k^{}|,`$ (43)
$`+{\displaystyle _{\sigma _c}}๐E\rho (E)|f^+(E)E^+|.`$
In complete analogy one can modify the spectral decomposition by deformation of the integration path to $`\mathrm{\Gamma }_{}=\mathrm{\Gamma }_+^{}`$ in the *upper* half of the complex plane, entering the domain of $`R_{}(z)`$ with singularities at $`\sigma _{}=\{z_k^{}\}`$. One then obtains the two bases
$`B_1^{}`$ $`=`$ $`(B_1^+)^{},`$ (44)
$`B_2^{}`$ $`=`$ $`(B_2^+)^{},`$ (45)
such that
$`|f_k^{}`$ $`=`$ $`{\displaystyle \frac{1}{\psi _k}}\underset{Ez_k^{}}{lim}(Ez_k^{})R_{}(E)|\psi `$ (46)
$`f_k^+|`$ $`=`$ $`{\displaystyle \frac{1}{\psi _k}}\underset{Ez_k^{}}{lim}(Ez_k^{})\psi |R_{}(E)`$ (47)
$`|f^{}(E)`$ $`=`$ $`|E^+,f^+(E)|=D_{}(E)E^+|,`$ (48)
$`|g^{}(E)`$ $`=`$ $`D_{}(E)|E^{},g^+(E)|=E^{}|.`$ (49)
Arbitrarily choosing $`B_1^{}`$ one obtains an *advanced* unity decomposition:
$`\mathrm{๐}_{}`$ $`=`$ $`{\displaystyle \underset{n}{}}|E_nE_n|+{\displaystyle \underset{k}{}}|f_k^{}f_k^+|`$ (51)
$`+{\displaystyle _{\sigma _c}}๐E\rho (E)|E^+f^+(E)|,`$
which is just the adjoint of the retarded unity decomposition, so $`(\mathrm{๐}_+)^{}\mathrm{๐}_{}`$. Note, however, that both decompositions are equivalent representations of the unity operator, thus in the sense of an operator identity we have $`\mathrm{๐}_+=\mathrm{๐}_{}=\mathrm{๐}`$.
### E Using the extended decomposition
The unity decompositions obtained in the preceding section can be used for different representations of vectors and operators. Using the retarded decomposition the Hamiltonian reads
$`H`$ $`=`$ $`{\displaystyle \underset{n}{}}E_n|E_nE_n|+{\displaystyle \underset{k}{}}z_k|f_k^+f_k^{}|`$ (53)
$`+{\displaystyle _{\sigma _c}}๐E\rho (E)E|f^+(E)E^+|.`$
The adjoint of the above expression yields the Hamiltonian in the *advanced* decomposition. Since both decompositions are equivalent representations, one sees that $`H^{}=H`$, i.e. the Hamiltonian remains self-adjoint.
The propagator in the retarded decomposition reads
$`U(t)`$ $`=`$ $`{\displaystyle \underset{n}{}}e^{iE_nt}|E_nE_n|+{\displaystyle \underset{k}{}}e^{iz_kt}|f_k^+f_k^{}|`$ (55)
$`+{\displaystyle _{\sigma _c}}๐E\rho (E)e^{iEt}|f^+(E)E^+|.`$
Since the complex numbers $`z_k`$ are located in the lower half of the complex plane, the contribution from the Gamov vectors decays in the future. If an initial Hilbert state $`|\psi `$ has Gamov components, they will disappear for $`t\mathrm{}`$. So if the state has no rotating components, i.e. $`E_n|\psi =0`$, the whole state will disappear in the future with the decay rates of the Gamov contributions given by the imaginary part of the contributing $`z_k`$:
$$\psi _k^+e^{iz_kt}=\psi _k^+e^{\gamma _kt}e^{i\nu _kt}\stackrel{t}{}0,$$
(56)
with $`\psi _k^+=f_k^+|\psi `$, $`z_k=\nu _ki\gamma _k`$ and $`\nu _k,\gamma _k>0`$. Each contribution from a Gamov state also induces a certain energy shift, given by the real part of the $`z_k`$. The advantage of the extended spectral decomposition is getting clear now: Not only the stable frequencies can immediately be read from the spectrum, but the decay rates and the energy shifts, too. By additionally checking the amplitudes $`\psi _k^+`$ between the initial state and the Gamov vectors one has immediate access to the dissipative dynamics of the system.
### F A closer look
Apart from the contribution of the Gamov vectors one has to consider the contribution from the generalized continuous eigenvectors. Let us call it the *background contribution*. Since the decomposition involves the complex distribution $`|f^+(E)=D_+(E)|E`$, a test vector $`\phi `$ has to be introduced. The background contribution then reads:
$`\phi |\psi _{bg}(t)`$ $`:=`$ $`{\displaystyle _{\sigma _c}}๐E\rho (E)e^{iEt}\phi |f^+(E)E^+|\psi `$ (57)
$`=`$ $`{\displaystyle _{\sigma _c}}๐E\rho (E)D_+(E)\psi _+(E)e^{iEt}\phi _+^{}(E)`$ (58)
where $`\psi _+(E)=E^+|\psi `$ and $`\phi _+^{}(E)=\phi |E^+`$. Due to $`D_+(E)`$ the integration path is transformed into the curve $`\mathrm{\Gamma }_+`$ below the complex eigenvalues $`z_k`$.
$$\phi |\psi _{bg}(t)=_{\mathrm{\Gamma }_+}๐z\rho (z)e^{izt}\psi (z)\phi ^{}(z),$$
(59)
where $`\psi _+(z)`$ and $`\phi _+^{}(z)`$ are the analytic continuations of respectively $`\psi _+(E)`$ and $`\phi _+^{}(E)`$ to the lower half of the complex plane (provided their existence). The curve $`\mathrm{\Gamma }_+`$ is arbitrary as long as $`\mathrm{\Gamma }_+\sigma _c`$ surrounds $`\sigma _+`$, i.e., $`\mathrm{\Gamma }_+`$ is placed *below* the poles $`z_k`$, which all have negative imaginary parts. For $`t>0`$, due to the exponential factor $`e^{izt}`$, the background contribution thus becomes more and more neglectable in comparison with the pole contributions. For initial states with high decay rates, i.e. pole contributions with small imaginary parts, the background contribution may be neglected already for small times $`t>0`$. This approximation is called the *pole approximation* and it is essentially equivalent to the *Wigner-Weisskopf approximation* (see ). The time evolution of the initial state may be approximated by an exponential decay and its energy spectrum by a Breit-Wigner distribution. This approximation is very popular, e.g. in quantum optics, and indeed in most cases it is a very good one.
### G Do divergences break time symmetry?
Obviously the above argument does not hold for negative times. Here, the contributions from the Gamov vectors diverge. In several publications of the Brussels School (see e.g. or *A. Bohm* in ) this is an argument in favour of a broken time symmetry, since for negative times the *advanced* decomposition (51) is the better choice. However, as will be shown below, this is a rather intuitive reasoning and does not enforce a break of time symmetry.
Let us see what happens with the retarded decomposition of a Hilbert state for negative times. Of course the Gamov contributions now diverge, but they are compensated by the background contribution. For $`t\mathrm{}`$ the background integral (59) gains influence, since the integration path $`\mathrm{\Gamma }_+`$ is placed *below* the complex eigenvalues $`z_k`$. The same argument, that for increasing $`t>0`$ the background contribution can be neglected in comparison with the pole contributions, is turned upside down for $`t<0`$. The smaller imaginary part of $`\mathrm{\Gamma }_+`$ induces a higher divergence than the contribution of the poles. By construction their divergences cancel each other, since the sum of both yields the contribution of the standard continuous eigenvectors (see eq. (31)), so a divergence of an initial Hilbert state backwards in time is prevented. The divergence of the Gamov contributions is even more than compensated: If a Hilbert state decays in the future, it will also decay in the past. To see this one can use the advanced decomposition (51) to investigate the systemโs evolution to the past and one would find a completely symmetric situation with majorizing pole contributions decaying backwards in time.
The situation is different if a Gamov vector $`|f_k^+`$ is taken as the initial state. Now there is only one pole contribution $`z_k`$ decaying in the future and diverging in the past. But Gamov vectors are only โexistingโ in the spectral representation of Hilbert states and they themselves do not refer to elements of physical reality, similiar to all other โgeneralized eigenstatesโ appearing in spectral decompositions. As constituted by the axiomatic foundations of quantum mechanics a physically realizable state is represented by a Hilbert vector (actually a family of Hilbert vectors normalized to unity and differing only in phase). Gamov states are no Hilbert vectors, so their properties are not expected to be physically meaningful.
Concluding, both decompositions are equivalent, but the use of the retarded decomposition is more *intuitive* for $`t>0`$. However, this should not be taken as the reason for a broken time symmetry. A more convincing argument in favour of a broken time symmetry is the splitting of *test spaces*. As we have seen in (59), analytic continuations of the test functions are needed to have the decompositions well-defined. In section IV we will get back to this point.
### H Perturbation theory
Besides the advantages pointed out in section II E we now face another important advantage of the Brussels Schoolโs formalism. Perturbation theory can be applied even in cases where standard quantum mechanics fails. Those cases are featured by the phenomenon of instability, which is strongly connected to the sensitivity of a system towards small perturbations. Let the systemโs Hamiltonian be of the form
$$H_\lambda =H_0+\lambda W,$$
(60)
where $`W`$ is the perturbation operator, switched on by the perturbation parameter $`\lambda [0,1]`$, and $`H_0`$ is the unperturbed part of the Hamiltonian, whose spectral decomposition is known and reads
$$H_0=\underset{n}{}\omega _n|\omega _n\omega _n|+_{\sigma _c}๐\omega \rho _0(\omega )\omega |\omega \omega |.$$
(61)
To simplify the discussion let the spectrum of $`H_0`$ be non-degenerate, so $`\omega _n\omega _m`$ for $`nm`$. As long as the perturbation strength $`\lambda `$ is not explicetely needed, it is set to 1 and $`H_{\lambda =1}`$ is written as $`H`$. The spectral decomposition of $`H`$ is yet unknown and reads in general
$$H=\underset{n}{}E_n|u_nu_n|+_{\sigma _c}๐\omega \rho (\omega )\omega |u^\pm (\omega )u^\pm (\omega )|,$$
(62)
where the sum is taken over the discrete spectrum $`\sigma _d`$ and the integration is performed over the continuous spectrum $`\sigma _c`$ of $`H`$.
The continuous spectrum $`\sigma _c`$ is not affected by the perturbation (provided the perturbation operator $`W`$ is $`H_0`$-compact, which is assumed here) and there is an identity mapping of perturbed and unperturbed continuous eigenvalues, so in both cases the integration is performed over $`\sigma _c`$. The discrete case is different. Switching to *Brioullin-Wigner* perturbation theory let us define pairs of orthogonal projection operators
$`P_n`$ $`:=`$ $`|\omega _n\omega _n|\text{and}`$ (63)
$`Q_n`$ $`:=`$ $`\mathrm{๐}P_n.`$ (64)
Furthermore let the perturbed eigenvector $`|u_n`$ of $`H`$ be not normalized to 1 but rather let $`u_n|\omega _n=1`$, such that $`P_n|u_n=|\omega _n`$. Applying these projections to the perturbed eigenvalue equation $`H|u_n`$ one obtains the *Brioullin-Wigner formulas* for the shifted eigenstate $`|u_n`$ and the shifted energy $`E_n`$:
$`|u_n`$ $`=`$ $`|\omega _n+Q_n{\displaystyle \frac{\lambda }{E_nH_0}}Q_nW|u_n,`$ (65)
$`E_n`$ $`=`$ $`\omega _n+\lambda \omega _n|W|u_n`$ (66)
Iteration leads to the corresponding perturbation series
$`|u_n`$ $`=`$ $`{\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}\left\{Q_n{\displaystyle \frac{\lambda }{E_nH_0}}Q_nW\right\}^p|\omega _n,`$ (67)
$`E_n`$ $`=`$ $`\omega _n+\lambda {\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}\omega _n|W\left\{Q_n{\displaystyle \frac{\lambda }{E_nH_0}}Q_nW\right\}^p|\omega _n.`$ (68)
Take a closer look at the above formulas (67) and (68). They are only valid if the expression
$`Q_n{\displaystyle \frac{1}{E_nH_0}}Q_n`$ $`=`$ $`{\displaystyle \underset{mn}{}}{\displaystyle \frac{|\omega _m\omega _m|}{E_n\omega _m}}`$ (70)
$`+{\displaystyle ๐\omega \rho (\omega )\frac{|\omega \omega |}{E_n\omega }}`$
is well-defined. Here *resonance* appears as a major obstacle. There are two different types of resonance making a perturbation expansion of eigenstates and eigenvalues impossible:
1. discrete resonance: The eigenvalue $`\omega _n`$ may be shifted in the course of the perturbation to a neighboured eigenvalue of $`H_0`$ (see fig. 6a). Expression (70) is no longer defined and the corresponding perturbation series diverges on account of the summation term.
2. continuous resonance: If the eigenvalue $`\omega _n`$ lies next to the continuous spectrum of $`H_0`$, it may be shifted there (see fig. 6b). Now the integral term in (70) causes a divergence of the perturbation series.
In both cases the perturbation expansion is limited to a finite convergence radius of the perturbation parameter $`\lambda `$ . Outside its range the shifted eigenvalue disappears from the spectrum of $`H_\lambda `$ and the mapping from the unperturbed to the perturbed spectrum is no longer one-to-one (see fig. 7).
Now regard the case of discrete unperturbed eigenvalues *embedded* in the continuous spectrum. Here the convergence radius of $`\lambda `$ reduces to zero, i.e. the perturbation series are no longer analytic in $`\lambda `$. An arbitralily small perturbation induces continuous resonance. As a consequence embedded states get unstable and start to decay. Decaying states cannot be part of the spectral decomposition of $`H_\lambda `$, they disappear from there and cannot be reconstructed for $`\lambda 0`$, i.e. $`H_\lambda \to ฬธH_0`$. Hence perturbation theory is not aplicable.
There is a way out: One has to include unstable states, i.e. Gamov vectors, into the spectral decomposition of $`H_\lambda `$. In terms of perturbation theory the Gamov vectors can be constructed by modified Brioullin-Wigner equations:
$`|f_k^\pm `$ $`=`$ $`|\omega _k+\underset{Ez_k^\pm }{lim}Q_k{\displaystyle \frac{\lambda }{EH_0\pm iฯต}}Q_kW|f_k^\pm ,`$ (71)
$`z_k^\pm `$ $`=`$ $`\omega _k+\lambda \omega _k|W|f_k^\pm ,`$ (72)
where $`z_k^+=z_k`$ are the singularities of $`R_+`$ on the lower half of the complex plane and $`z_k^{}=z_k^{}`$ are those of $`R_{}`$ on the upper half. There is a one-to-one mapping between the disappearing discrete eigenvalues $`\omega _k`$ and the complex eigenvalues $`z_k`$. Iteration leads to the corresponding perurbation series
$`|f_k^\pm `$ $`=`$ $`\underset{Ez_k^\pm }{lim}{\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}\left\{Q_k{\displaystyle \frac{\lambda }{EH_0\pm iฯต}}Q_kW\right\}^p|\omega _k,`$ (73)
$`z_k^\pm `$ $`=`$ $`\omega _k+\lambda \omega _k|W|f_k^\pm .`$ (74)
The generalized continuous eigenstates are
$`|f^+(\omega )`$ $`=`$ $`D_+(\omega )|u^+(\omega )`$ (75)
$`f^{}(\omega )|`$ $`=`$ $`u^+(\omega )|,`$ (76)
with the complex distribution $`D_+`$ defined in (25). The standard continous eigenvectors can be constructed using the *Lippmann-Schwinger equation*
$$|u^+(\omega )=|\omega +\frac{\lambda }{\omega H_0+iฯต}W|u^+(\omega )$$
(77)
or its iteration, the *Born series*:
$$|u^+(\omega )=\underset{p=0}{\overset{\mathrm{}}{}}\left\{\frac{\lambda }{\omega H_0+iฯต}W\right\}^p|\omega .$$
(78)
Take a closer look at the construction of Gamov states in eq. (71). The divergences due to continuous resonance are prevented by introducing a small complex shift in the denominator and analytic continuation to the second Riemann sheet. As a consequence, the energy eigenvalues are shifted to respectively the lower and upper complex plane, inducing exponential decay and Breit-Wigner energy distribution of the corresponding Gamov states. The expressions are still analytic in $`\lambda `$, so perturbation theory can be applied. This is a major advantage of the Brussels Schoolโs formalism.
## III Example: the Friedrichs model
The Friedrichs model (K. Friedrichs 1948, see ) is the simplest model involving instability. Here one single discrete state is coupled to a continuum of states by an interaction Hamiltonian $`W`$. A physical analogon to this model is e.g. the *Wigner-Weisskopf model* (Wigner and Weisskopf 1930, see ) of a twolevel atom coupled to the vacuum field. Another example is the *Auger effect* of autoionisation (see e.g. ).
The total Hamiltonian is of the form
$`H`$ $`=`$ $`H_0+W,`$ (79)
$`\text{where }H_0`$ $`=`$ $`\omega _1|11|+{\displaystyle _0^{\mathrm{}}}๐\omega \omega |\omega \omega |,`$ (80)
$`W`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}๐\omega \left\{W(\omega )|\omega 1|+W^{}(\omega )|1\omega |\right\}.`$ (81)
Using the *second resolvent identity*
$$R=R_0+RWR_0,$$
(82)
with the resolvent $`R_0(z)=\frac{1}{zH_0}`$ of the free Hamiltonian $`H_0`$ and the resolvent $`R(z)=\frac{1}{zH}`$ of the total Hamiltonian. The resolvent $`R(z)`$ can be obtained by applying it to the free states $`|1`$ and $`|\omega `$ and doing some rearrangement:
$`R(z)`$ $`=`$ $`{\displaystyle \frac{1}{\eta (z)}}\left[|1+{\displaystyle _0^{\mathrm{}}}๐\omega {\displaystyle \frac{W(\omega )}{z\omega }}|\omega \right]`$ (85)
$`\left[1|+{\displaystyle _0^{\mathrm{}}}๐\omega {\displaystyle \frac{W^{}(\omega )}{z\omega }}\omega |\right]`$
$`+{\displaystyle _0^{\mathrm{}}}๐\omega {\displaystyle \frac{1}{z\omega }}|\omega \omega |.`$
The complex function
$$\eta (z)=z\omega _1_0^{\mathrm{}}๐\omega \frac{|W(\omega )|^2}{z\omega }$$
(86)
has analytic extensions $`\eta _\pm (\omega )=\eta (\omega \pm iฯต)`$ from above and below to the real axis. The zeros of $`\eta `$ on the second Riemann sheet define the complex poles of the resolvent:
$$\eta _+(z_1)=0,\eta _{}(z_1^{})=0.$$
(87)
The pole $`z_1`$ can be obtained by numerically calculating the complex zero of $`\eta _+(z)`$. In first order one gets
$$z_1\omega _1+๐ซ_0^{\mathrm{}}๐\omega \frac{|W(\omega )|^2}{\omega _1\omega }i\pi |W(\omega _1)|^2.$$
(88)
The Gamov vectors can be obtained using eqns. (32) and (33):
$`|f_1^+`$ $`=`$ $`{\displaystyle \frac{1}{c_1}}\underset{Ez_1}{lim}(Ez_1)R_+(E)|1`$ (89)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{\eta _+{}_{}{}^{}(z_1)}}}\left[|1+{\displaystyle _0^{\mathrm{}}}๐\omega {\displaystyle \frac{W(\omega )}{[z_1\omega ]_+}}|\omega \right]`$ (90)
$`f_1^{}|`$ $`=`$ $`{\displaystyle \frac{1}{c_1}}\underset{Ez_1}{lim}(Ez_1)1|R_+(E)`$ (91)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{\eta _+{}_{}{}^{}(z_1)}}}\left[1|+{\displaystyle _0^{\mathrm{}}}๐\omega {\displaystyle \frac{W^{}(\omega )}{[z_1\omega ]_+}}\omega |\right],`$ (92)
where the complex distribution is defined by
$$_0^{\mathrm{}}๐\omega \frac{\phi (\omega )}{[z_1\omega ]_+}:=\underset{Ez_1}{lim}_0^{\mathrm{}}๐\omega \frac{\phi (\omega )}{E\omega +iฯต}.$$
(93)
The standard outgoing states can be constructed using eq. (15):
$`|u^+(\omega )`$ $`=`$ $`{\displaystyle \frac{1}{c(\omega )}}{\displaystyle \frac{iฯต}{EH+iฯต}}|\omega `$ (94)
$`=`$ $`|\omega +{\displaystyle \frac{W^{}(\omega )}{\eta _+(\omega )}}\times `$ (96)
$`\times \left[|1+{\displaystyle _0^{\mathrm{}}}๐\omega ^{}{\displaystyle \frac{W(\omega ^{})}{\omega \omega ^{}+iฯต}}|\omega \right].`$
Using eq. (36) the generalized continuous eigenvectors read
$`|f_+(\omega )`$ $`=`$ $`D_+(\omega )|u^+(\omega )`$ (97)
$`=`$ $`|\omega +D_+(\omega ){\displaystyle \frac{W^{}(\omega )}{\eta _+(\omega )}}\times `$ (99)
$`\times \left[|1+{\displaystyle _0^{\mathrm{}}}๐\omega ^{}{\displaystyle \frac{W(\omega ^{})}{\omega \omega ^{}+iฯต}}|\omega \right],`$
$`f^{}(\omega )|`$ $`=`$ $`u^+(\omega )|`$ (100)
$`=`$ $`\omega |+{\displaystyle \frac{W(\omega )}{\eta _{}(\omega )}}\times `$ (102)
$`\times \left[1|+{\displaystyle _0^{\mathrm{}}}๐\omega ^{}{\displaystyle \frac{W^{}(\omega ^{})}{\omega \omega ^{}iฯต}}\omega |\right],`$
where the complex distribution $`D_+`$ is defined by (25) with the curve $`\mathrm{\Gamma }_+`$ leading from $`0`$ to $`\mathrm{}`$ below $`z_1`$.
In the pole approximation, i.e. neglecting the background contribution from the continuous eigenvectors, the initial state $`|1`$ decays with a decay rate $`\mathrm{\Gamma }`$ given by two times the imaginary part of $`z_1`$,
$`P_1(t)`$ $`=`$ $`|1|U(t)|1|^2e^{\mathrm{\Gamma }t},`$ (103)
$`\text{where}\mathrm{\Gamma }`$ $`=`$ $`2\mathrm{I}mz_12\pi |W(\omega _1)|^2,`$ (104)
which is in agreement with Fermiโs Golden Rule.
## IV Time symmetry
### A The Rigged Hilbert Space
When dealing with infinite dimensions, the Hilbert space appears to be too small in the following sense: Position and momentum operator have no Hilbert eigenvectors and Diracโs bra-ket formalism is mathematically not justified. To correct this problem, the Hilbert space must be enlarged to include *improper states* like $`\delta `$-functions and plane waves. The framework to enable these features is given by the theory of the *Gelfand triplet*, also called the *Rigged Hilbert Space* (RHS).
The concept is the following (comp. e.g. ). Take the Hilbert space $`=L^2`$ of squared integrable functions and find a suitable test space $`\mathrm{\Phi }`$ with the properties
1. $`\mathrm{\Phi }`$ is a locally convex topological vector space,
2. $`\mathrm{\Phi }`$ is complete with respect to its own topology,
3. $`\mathrm{\Phi }`$ is dense in $``$,
such that the topological dual $`\mathrm{\Phi }^{}`$ of $`\mathrm{\Phi }`$, consisting of all linear continuous functionals on $`\mathrm{\Phi }`$, called *distributions*, includes the desired improper states. The Hilbert space itself is his own topological dual, i.e. $`^{}=`$. For any test space $`\mathrm{\Phi }`$ being smaller than $``$, the Gelfand triplet
$$\mathrm{\Phi }\mathrm{\Phi }^{}$$
(105)
makes it possible to extend the bracket $`\chi |\phi `$ to all pairs $`\phi \mathrm{\Phi }`$, $`\chi \mathrm{\Phi }^{}`$ by
$`\chi |\phi `$ $`:=`$ $`\chi (\phi )`$ (106)
$`\phi |\chi `$ $`:=`$ $`\chi ^{}(\phi ).`$ (107)
On pairs of Hilbert vectors the bracket conicides with the scalar product in $``$. The โketsโ and โbrasโ are now identified by linear functionals on the opposite dual space, i.e.
$`|\xi `$ $`:=`$ $`|\xi `$ (108)
$`\text{and}\xi |`$ $`:=`$ $`\xi |,`$ (109)
where $`\xi ,\mathrm{\Phi }`$ or $`\mathrm{\Phi }^{}`$. An operator $`A`$ can be extended to the distribution space $`\mathrm{\Phi }^{}`$ using
$$A\chi |\phi =\chi |A^{}\phi .$$
(110)
Hence $`A`$ can only be extended to $`\mathrm{\Phi }^{}`$ if its adjoint $`A^{}`$ is defined on and does not lead out of the test space $`\mathrm{\Phi }`$:
$$A^{}\mathrm{\Phi }\mathrm{\Phi }.$$
(111)
The test space of standard quantum mechanics is the Schwartz space $`๐ฎ`$ of rapidly decreasing functions, defined by
$$๐ฎ:=\{fC^{\mathrm{}}|\underset{x}{sup}|x^n_x^mf(x)|<\mathrm{}n,m\}.$$
(112)
The Schwartz space is a physically convenient test space, since here position and momentum operator are everywhere defined, they fulfill (111), they are self-adjoint and have their eigenvectors among the *tempered distributions* in $`๐ฎ^{}`$. Furthermore, they keep their real spectrum, which makes the RHS $`(๐ฎ,L^2,๐ฎ^{})`$ be called a *tight rigged* Hilbert space. Another popular test space is the space $`๐`$ of $`C^{\mathrm{}}`$-functions with compact support. But $`๐`$ is not tight rigged, since the momentum operator obtains a complex spectrum with complex plane waves $`e^{izx}๐^{}`$ as eigenvectors.
### B Breaking time symmetry
Now we turn to the RHS of the Brussels School. As already stated in section II D the definition of generalized eigenvectors is only possible if the test functions can be analytically continued to whether the lower or the upper half of the complex plane. The retarded right eigenvectors $`|f_k^+`$ and $`|f^+(E)`$ involve complex distributions continuating the test functions $`\phi ^{}(E)=\phi |E`$ to the lower half of the complex plane, denoted by $`_{}`$. There is a class of functions fulfilling this condition: the Hardy class $`H_{}^2`$ from below (for detailed information on Hardy classes see ). If $`\phi ^{}(E)`$ is in $`H_{}^2`$ then its complex conjugate $`\phi (E)=E|\phi `$ is in the Hardy class $`H_+^2`$ from above, whose members can be analytically continuated to the upper half of the complex plane, denoted by $`_+`$. In detail: The function $`\phi (E)`$ is in $`H_\pm ^2`$ if and only if it has an analytic continuation to $`_\pm `$, such that there is a $`C<\mathrm{}`$ and
$$\underset{y>0}{sup}_{\mathrm{}}^{\mathrm{}}๐E|\phi (E\pm iy)|^2<C.$$
(113)
Thus, restricted to the real line, Hardy functions are $`L^2`$-functions. An equivalent definition can be given using the *Paley-Wiener theorem*: The Hardy spaces $`H_\pm ^2`$ are formed by the inverse Fourier transforms of $`L^2`$-functions with support on respectively the positive and negative semiaxis $`_\pm `$:
$$H_\pm ^2=\mathrm{F}^1\{L^2(_\pm )\}.$$
(114)
For physical convenience it is necessary, that the test functions are also in the Schwartz class $`๐ฎ`$. Furthermore, their energy representation can be restricted to the positive semiaxis, since energy is bounded from below by zero. Altogether, the test spaces of the Brussels School are defined by
$$\mathrm{\Phi }_\pm :=(H_\pm ^2๐ฎ)|__+.$$
(115)
The topological duals $`\mathrm{\Phi }_\pm ^{}`$ then include the complex distributions used in the generalized eigenvectors. The formalism of the Brussels School can be applied in the RHS
$$\mathrm{\Phi }_\pm \mathrm{\Phi }_\pm ^{}.$$
(116)
Since there are *two* RHSโs, one has to investigate how to deal with them. The retarded unity decomposition (43) contains the dyadic products $`|f_k^+f_k^{}|`$. Applying test vectors $`\phi _+`$ and $`\phi _{}`$ to both sides one obtains
$$\phi _+|f_k^+f_k^{}|\phi _{}.$$
(117)
The above expression only makes sense, if $`\phi _+^{}(E)`$ and $`\phi _{}(E)`$ are both in $`\mathrm{\Phi }_{}`$, thus $`\phi _+(E)`$ must be in $`\mathrm{\Phi }_+`$. Hence the instructions for computing brackets while using the retarded unity decomposition are: On the left hand you have to use $`\mathrm{\Phi }_+`$-vectors and on the right hand $`\mathrm{\Phi }_{}`$-vectors. However, since both test spaces are dense in $``$, you would not have to care too much, if there was not a significant subtlety: time symmetry is broken. To see this we investigate the time evolution of a test vector $`\phi _{}\mathrm{\Phi }_{}`$:
$$\phi _{}^t(E)=e^{iEt}\phi _{}(E).$$
(118)
Since $`\phi (E)`$ is in $`H_{}^2`$, there is an analytic continuation $`\phi _{}(Eiy)`$ to $`_{}`$ and so there is one for its propagated counterpart, given by
$$\phi _{}^t(Eiy)=e^{yt}e^{iEt}\phi _{}(Eiy).$$
(119)
Obviously, condition (113) is violated for $`\phi _{}^t`$ and $`t<0`$, thus time evolution leads out of the test space for $`t<0`$. As a consequence test functions in $`\mathrm{\Phi }_{}`$ can only be propagated to the future. Analogically, test functions in $`\mathrm{\Phi }_+`$ can only be propagated to the past. Using condition (111) one finds in addition that the extension of the unitary propagator $`U(t)`$ to the distribution spaces $`\mathrm{\Phi }_+^{}`$ and $`\mathrm{\Phi }_{}^{}`$ is only possible for respectively $`t>0`$ and $`t<0`$. So retarded Gamov vectors $`|f_k^+`$ in $`\mathrm{\Phi }_+^{}`$ โpropagateโ (as part of the spectral decomposition of Hilbert vectors) only to the future and advanced Gamov vectors $`|f_k^{}`$ in $`\mathrm{\Phi }_{}^{}`$ propagate only to the past. The unitary propagator group $`๐ฐ`$ hence splits into a semigroup $`๐ฐ_+`$ propagating to the future (and applying to vectors in $`\mathrm{\Phi }_{}`$ and $`\mathrm{\Phi }_+^{}`$) and a second one $`๐ฐ_{}`$ propagating to the past (and applying to vectors in $`\mathrm{\Phi }_+`$ and $`\mathrm{\Phi }_{}^{}`$). This feature is called by the Brussels School the *intrinsic irreversibility* of quantum mechanics and is identified with a microscopic arrow of time.
### C Time symmetry and Gamov vectors
Is the special choice of the two RHSโs *necessary* for the correct mathematical implementation of Gamov states? If so, one could regard the splitting of the unitary time evolution group into two semigroups as enforced by the mathematical concept. This point of view is taken by *A. Bohm* and *N.L. Harshman* in :
> โPhilosophizing alone would not be enough to take a semigroup instead. To arrive at the semigroup, we start from the empirically desirable properties of Gamov resonance states and let mathematics determine the path.โ
The above mentioned properties are given by
$`H|f_k^+`$ $`=`$ $`(E_Ri\mathrm{\Gamma }/2)|f_k^+,`$ (120)
$`|f_k^+(t)`$ $`=`$ $`e^{iHt}|f_k^+=e^{iE_Rt}e^{i\frac{\mathrm{\Gamma }}{2}t}|f_k^+,`$ (121)
with $`E_R,\mathrm{\Gamma }>0`$. Starting from these conditions, following the reasoning of section IV B, one is led to the two RHSโs with broken time symmetry. However, the argumentation is not valid if a Rigged Hilbert Space can be found, where Gamov vectors are included and time symmetry is preserved. Let us start with the space $`๐`$ of $`C^{\mathrm{}}`$-functions with compact support. As already mentioned, this space is a perfect test space fulfilling the requirements given in section IV A. Furthermore we have $`๐๐ฎ`$. Now regard the space $`๐ต`$ of inversely Fourier transformed $`๐`$-functions:
$$๐ต:=\mathrm{F}^1\{๐\}.$$
(122)
Since the inverse Fourier transform is a one-to-one mapping on the Schwartz space $`๐ฎ`$, we have
$$๐ต๐ฎ.$$
(123)
Let $`\stackrel{~}{\phi }(s)`$ be a $`๐`$-function. Its inverse Fourier transform $`\phi (E)๐ต`$ reads
$$\phi (E)=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}๐se^{iEs}\stackrel{~}{\phi }(s)=\frac{1}{2\pi }_K๐se^{iEs}\stackrel{~}{\phi }(s),$$
(124)
where $`K`$ is the compact support of $`\stackrel{~}{\phi }(s)`$. The analytic continuation of $`\phi (E)`$ to the complex plane exists and is an entire function, given by
$$\phi (E+iy)=\frac{1}{2\pi }_K๐se^{ys}e^{iEs}\stackrel{~}{\phi }(s),$$
(125)
for all $`E,y`$ (see , theorem 7.22). Since the inverse Fourier transform is a linear continuous one-to-one mapping, the transformed topology $`\tau _๐`$ of $`๐`$,
$$\tau _๐ต:=\mathrm{F}^1\{\tau _๐\},$$
(126)
transfers all topological properties from $`๐`$ to $`๐ต`$. Thus, $`\{๐ต,\tau _๐ต\}`$ is a topological vector space fulfilling all requirements of a test space in $`L^2()`$. Since $`๐ต`$ contains functions being analytic in the whole complex plane, the topological dual $`๐ต^{}`$ contains the retarded Gamov vectors defined by (120) and (121). Furthermore, it contains also the *advanced* Gamov vectors $`|f_k^{}`$ with eigenvalues $`z_k^{}`$ in the upper half of the complex plane, as well as any other generalized eigenvector used in the extended spectral decompositions (43) and (51). Because of relation (123) all $`๐ต`$-functions are $`๐ฎ`$-functions, so $`๐ต`$ is already a physical convenient test space and standard quantum mechanics is applicable, too. Since $`๐๐ฎ(_\pm )๐`$ the space $`๐ต`$ does neither include any of the test spaces $`\mathrm{\Phi }_\pm `$ nor is it part of it, i.e.
$$๐ต\mathrm{\Phi }_\pm ๐ต,$$
(127)
but of course $`๐ต\mathrm{\Phi }_\pm \mathrm{}`$. So $`๐ต`$ is a different test space, whose topological duals include any of the complex distributions needed for a comfortable treatment of decaying systems.
However, there is no intrinsic irreversibility. Let the $`๐`$-function $`\stackrel{~}{\phi }(s)`$ have a compact support bounded by the finite interval $`[a,b]`$. Its inverse Fourier transform $`\phi (E)`$ can be propagated in time yielding $`\phi _t(E)=e^{iEt}\phi (E)`$, whose Fourier transform,
$$\stackrel{~}{\phi }_t(s)=_{\mathrm{}}^{\mathrm{}}๐Ee^{iEs}e^{iEt}\phi (E)=\stackrel{~}{\phi }(s+t),$$
(128)
has a compact support bounded by the finite interval $`[a+t,b+t]`$, thus it remains compact for all $`t`$ and $`\stackrel{~}{\phi }_t(s)`$ stays in $`๐`$, so $`\phi _t(E)`$ stays in $`๐ต`$. Hence the propagator $`U(t)`$ does not lead out of the test space for any $`t`$. Consequently, the unitary propagator group $`๐ฐ=\{U(t)\}`$ can be extended to the distribution space $`๐ต^{}`$ without any splitting. The support of $`๐ต`$-functions can be restricted to the positive semiaxis, due to the positivity of energy, and the RHS
$$\left(๐ตL^2๐ต^{}\right)|__+$$
(129)
is thus suitable for the presented concepts. Any desired unity decomposition including Gamov vectors, retarded or advanced, can equivalently be used. As far as Hilbert states are concerned, no unphysical behaviour is predicted and any calculated physical quantity is numerically identical to the one calculated by standard methods.
### D Physical justification
Apart from pure mathematics there is a physical reasoning presented by members of the Brussels School who try to derive the special choice of the two Hardy classes from the assumption of *causality* (see e.g. Schulte and Twarock in or Antoine, Bohm and Harshman in , from where the citations below are taken). This is done the following way.
#### 1 Argument
Starting point is a special assumption of *causality* which is called the โpreparation$``$registration arrow of timeโ:
> โTime translation of the registration apparatus relative to the preparation apparatus makes sense only by an amount of $`t0`$.โ
This is translated into the language of quantum mechanics to a so-called โquantum mechanical arrow of timeโ:
> โA state $`\varphi ^+(t)`$ must be prepared before an observable $`|\psi ^{}\psi ^{}|=|\psi ^{}(0)\psi ^{}(0)|`$ can be measured in that state, i.e. $`\varphi ^+=\varphi ^+(0)`$ must be prepared during $`t0`$.โ
Interpreting measurement formally as a scattering experiment, the prepared state $`\varphi ^+`$ is identified with an incoming state and the measured state $`\psi ^{}`$ with an outgoing state (see fig. 9). Survival probabilities of the initial (i.e. prepared) state are thus calculated by transition probabilities of the incoming and outgoing state:
$`P(t)`$ $`=`$ $`|\varphi ^+|e^{iHt}|\psi ^{}|^2\text{for }t0\text{,}`$ (130)
$`P(t)`$ $`=`$ $`0\text{for }t<0\text{.}`$ (131)
Using the energy decomposition this implies
$`{\displaystyle ๐E\rho (E)e^{iEt}\varphi ^+(E)}`$ $`=`$ $`0\text{for }t>0`$ (132)
$`\text{and }{\displaystyle ๐E\rho (E)e^{iEt}\psi ^+(E)}`$ $`=`$ $`0\text{for }t<0\text{.}`$ (133)
Using the Paley-Wiener theorem (114) one is led to the desired conclusion
$$\varphi ^+H_+^2,\psi ^{}H_{}^2.$$
(134)
Hence, founding on the given formulation of causality the only correct choice of a Rigged Hilbert Space is the choice of the two RHSโs of the Brussels School. In other words: Standard quantum mechanics has not yet found a microscopic arrow of time because it has been incomplete in its axiomatic foundation.
#### 2 Critics
The above argumentation is correct, but it cannot be taken as a *derivation* of the arrow of time. As is obvious even by the choice of words (โbeforeโ, โafterโ, etc.), the given โcausalityโ assumption is already based on the arrow of time, hence it is logically impossible to derive the latter from the former. Causality itself does not require an arrow of time, since neither Newtonโs law nor Schrรถdingerโs law contradict causality, although both theories are perfectly time-symmetric. Yet, the โpreparation$``$registration arrow of timeโ refers explicitely to *open subsystems*, namely the preparation apparatus and the measuring apparatus, which are coupled together to form a composite system. Though one cannot generally infer from statements about the time evolution of open subsystems to statements about the evoution of closed composite systems, especially not when the subsystems are involved in a measurement process. Hence postulating a โpreparation$``$registration arrow of timeโ for *all* systems equals introducing a microscopic arrow of time *by hand* into quantum mechanics. Once the arrow of time has been put into the mathematics, it is not surprising that it will reappear eventually. The situation can be visualized by fig. 10.
Thus, the special choice of the two RHSโs must be regarded as being *equivalent to a postulation of a global microscopic arrow of time*. The split RHSโs are a mathematical mirror image of the empirically motivated splitting between past and future.
### E The arrow of time as a matter of choice
As already stated in the introduction, it is possible to explain an arrow of time due to time-asymmetric evolution of *macroscopic* observables and due to *local* observations on subsystems. But respecting a *global microscopic arrow of time*, there is a major difference between standard and Rigged Hilbert Space quantum mechanics. In standard quantum mechanics there is no way to implement a global microscopic arrow of time except perhaps by applying brute force and splitting the unitary evolution group axiomatically into two semigroups. In Rigged Hilbert Space quantum mechanics one has the choice of either time-symmetric or time-asymmetric microscopic evolution, due to the choice of the test spaces. Hence, in the RHS formulation it is possible to implement a global microscopic arrow of time in contrast to standard quantum mechanics. So one could argue that nature *obviously* has taken the time-asymmetric choice of two test spaces instead of one. But what does โobviouslyโ mean? One might mention our everyday experiences or global time-asymmetric boundary conditions like the big bang or the second law of thermodynamics. Though these experiences and boundary conditions are all related to subsystems or far from being microscopic. The โuniverseโ is a construction from observations on subsystems on a *very* macroscopic scale. Altogether, these time-asymmetric boundary conditions all base on *local macroscopic observations* and hence, as already stated, they do not give rise to a real problem with quantum mechanics, which is a set of microscopic laws. So the choice of global microscopic time-symmetry or time-asymmetry remains merely axiomatic and cannot be derived from other physical postulates.
## V Summary
The quantum mechanical formalism, developped by the Brussels School, offers a comfortable treatment of decaying systems, where standard quantum mechanics requires a lot more effort. The generalized spectral decompositions involve Gamov states and can be interpreted in a physically intuitive manner. They directly yield important physical quantities like decay rates and energy level shifts. Perturbation expansions become possible in cases, where standard quantum mechanics fails.
In contrast to standard quantum mechanics, where the implementation of a microscopic arrow of time is not possible, the Rigged Hilbert Space formalism offers the possibility to incorporate such an arrow by the special choice of two different Gelfand triplets. Due to this choice, the unitary time evolution group splits into two semigroups, each one for respectively the future and the past. However, it is shown that the use of Gamov states does not require this splitting. As a counter example, one single Gelfand Triplet is constructed, where the formalism of the Brussels School can be applied and where time symmetry is preserved. Furthermore, it is pointed out that the physical arguments in favour of the special choice of two Gelfand Triplets are not independent from the arrow of time and thus cannot be used for a derivation of the latter. Hence even within the framework of the Rigged Hilbert Space formalism, the global microscopic arrow of time remains a postulate.
## VI Acknowledgements
I would like to thank F. Petruccione, H.-P. Breuer, M. Bordemann, and M. Wilkens for stimulating discussions and editorial advice and the referee for valuable comments. I also thank all other people talking with me about the topic of this paper and contributing in many ways to its completion.
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# 1 Introduction
## 1 Introduction
The topological particle, whose canonical BRST quantization is developed and applied in this paper, is the simplest example of a topological quantum theory. There are many reasons, both physical and mathematical, for studying such theories; however much of the work to date has been carried out in the Lagrangian approach, using the functional integral as the starting point. The underlying motivation for this paper is the belief that a serious canonical analysis of such theories should be fruitful.
The topological particle and its quantization is described by Beaulieu and Singer in , but these authors concentrate on the case based on a constant function on the manifold, while in this paper the model is based on a Morse function for the manifold, a function which is anything but constant, having only isolated critical points, and encodes information about the topology of the manifold. It is first shown, in section 2, that the the supersymmetric quantum mechanical system which arises on BRST quantization of the model is the system used by Witten in his work on supersymmetry and Morse theory. The topological origin of this model gives a natural explanation for the form of the matrix elements for the theory which in Wittenโs paper are calculated by instanton methods.
The main new result of the present paper is a path integral formula for the calculation of these matrix elements, which is derived in section 4 using the methods of stochastic calculus on manifolds. Paths are defined by a stochastic differential equation which is essentially the Nicolai map for the model ; the paths encode fluctuations about classical trajectories and thus lead to a fully rigorous path integral WKB method (as derived by Blau, Keski-Vakkuri and Niemi using physicistsโ methods ). The result makes contact again with the classical action of the topological particle which was the starting point.
Having established the precise form of the matrix elements of the model, in section 5 these are used to construct explicitly the cochains and cohomology for the model of the manifold cohomology introduced by Witten based on the critical points of a Morse function.
## 2 Classical Dynamics
The topological particle model introduced by Beaulieu and Singer is defined by the action
$$S\left(x(.)\right)=_0^tv_\mu (x(t^{}))\dot{x}^\mu (t^{})dt^{}.$$
(1)
The fields $`x`$ are smooth maps $`x:IM`$ from $`I`$ (the interval $`[0,t]`$ of the real line) into a compact $`n`$-dimensional Riemannian manifold $`M`$, while $`v=dh`$ is an exact one-form on $`M`$. The components $`v_\mu `$ in local coordinates $`x^\mu ,\mu =1,\mathrm{},n`$ are as usual defined by $`v=v_\mu dx^\mu `$, so that $`v_\mu =\frac{h}{x^\mu }`$.
Clearly, since $`v`$ is the differential of $`h`$, this action can be expressed more simply as
$$S\left(x(.)\right)=h(x(t))h(x(0)).$$
(2)
This form of the action shows that the model is indeed topological in nature, a related point being that the equation of motion for $`x`$ is trivially satisfied. However the form of the action (1) involving positions and velocities is required for the passage from the Lagrangian to the Hamiltonian form. While Beaulieu and Singer considered the case $`v=0`$, in this paper a more general situation is considered; in particular in section 5 the function $`h`$ is taken to be a Morse function, that is, a function on $`M`$ with isolated critical points.
It is evident that the action (1) is highly symmetric, depending only on the endpoints of the path. It might thus be naively supposed that the path integral
$$_{\text{paths}/\text{symmetries}}๐x(.)\mathrm{exp}\left(S(x(.))\right)$$
(3)
would be trivial. This is not in fact the case because the โmeasureโ $`๐x`$ is not simply some limit of a product measure, but must be derived by careful canonical quantization of the theory, which is carried out below.
The first step in this process is to investigate the classical Hamiltonian dynamics. From the action (1) the Lagrangian of the theorem is seen to be
$$(x,\dot{x})=v_\mu (x)\dot{x}^\mu ,$$
(4)
so that the Euclidean time Legendre transformation to the phase space $`T^{}M`$ (the cotangent bundle of $`M`$) gives as momentum conjugate to $`x^\mu `$
$$p_\mu =i\frac{\delta (x,\dot{x})}{\delta \dot{x}^\mu }=iv_\mu .$$
(5)
The symmetries of the system now manifest themselves as $`n`$ constraints on the phase space $`T^{}M`$:
$$T_\mu p_\mu +iv_\mu (x)=0,\mu =1,\mathrm{},n.$$
(6)
The Poisson brackets on the phase space $`T^{}M`$ are obtained from the standard symplectic form $`\omega =dp_\mu dx^\mu `$, so that as usual $`\{x^\mu ,p_\nu \}=\delta _\nu ^\mu `$. Direct calculation shows that (since $`v`$ is closed)
$$\{T_\mu ,T_\nu \}=0.$$
(7)
The Hamiltonian of the system is, by the Euclidean time prescription,
$$H=ip_\mu \dot{x}^\mu +(x,\dot{x})=0,$$
(8)
so that the constraints are first class and abelian.
As is standard in a topological theory, the constraints are of a number that seems to preclude any interesting dynamics - in this case the system has a $`2n`$-dimensional phase space with $`n`$ first class constraints, so that by naive counting one would expect the corresponding reduced phase space to be trivial. In fact the theory does capture some topological information as will emerge below.
The first indication of this comes from considering gauge-fixing, which shows that the reduced phase space, while as expected zero dimensional, corresponds to the critical points of $`h`$. The reduced phase space is defined to be the quotient of the subspace of the phase space $`T^{}M`$ on which the constraints hold by the action of the group generated by the constraints . Classically gauge-fixing conditions are sought which pick out one point in each orbit of this group; in this case a natural choice is the set of $`n`$ conditions $`X^\mu g^{\mu \nu }(p_\nu iv_\nu )=0`$. (Justification for this choice can only be fully made on quantization.) Taken together the constraints and the gauge-fixing condition are satisfied when $`p_\mu =0`$ and $`v_\mu =0`$, that is, at the critical points of the manifold. To see how this finite reduced phase space can provide topological information we turn to quantization, using the BRST approach.
## 3 BRST quantization
To implement the constraints and gauge-fixing at the quantum level we use the BRST quantization in canonical form , introducing ghosts and their conjugate momenta. For this process two supermanifolds are required, a super configuration space $`SM`$ with even local coordinates $`x^\mu `$ and odd local coordinates $`\eta ^\mu `$ and a super phase space $`SPM`$ with even local coordinates $`x^\mu `$ and $`p_\mu `$ and odd local coordinates $`\eta ^\mu `$ and $`\pi _\mu `$. (In each case the index $`\mu `$ runs from $`1`$ to $`n`$.) The $`(n,n)`$ dimensional supermanifold $`SM`$ is built from the tangent bundle of $`M`$, with coordinate patches corresponding to those on $`M`$ and changes of the coordinates $`x^\mu `$ between patches being those on $`M`$ while those of the coordinates $`\eta ^\mu `$ are defined by
$$\stackrel{~}{\eta ^\mu }=\frac{\stackrel{~}{x}^\mu }{x^\nu }\eta ^\nu .$$
(9)
For future reference we note that there is a well-defined projection $`ฯต:SMM`$ defined by
$$ฯต(x,\eta )=x.$$
(10)
The super phase space $`SPM`$ is the cotangent bundle to $`SM`$, so that $`p_\mu `$ and $`\pi _\mu `$ transform according to the rule
$$\stackrel{~}{p_\mu }=\frac{x^\nu }{\stackrel{~}{x}^\mu }p_\nu ,\stackrel{~}{\pi _\mu }=\frac{x^\nu }{\stackrel{~}{x}^\mu }\pi _\nu .$$
(11)
The simplest, and natural, choice of symplectic form on this manifold, which makes $`\pi _\mu `$ the conjugate momentum to $`\eta _\mu `$, is
$`w_s`$ $`=`$ $`d\left(p_\mu dx^\mu +\pi _\mu D\eta ^\mu \right)`$ (12)
$`=`$ $`dp_\mu dx^\mu +D\pi _\mu D\eta ^\mu {\displaystyle \frac{1}{2}}R_{\mu \nu \lambda }{}_{}{}^{\kappa }\pi _{\kappa }^{}\eta ^\lambda dx^\mu dx^\nu ,`$
where the Levi-Cevita connection corresponding to the Riemannian metric $`g`$ has been used, with Christoffel symbols $`\mathrm{\Gamma }_{\mu \nu }^\kappa `$ and curvature tensor components $`R_{\mu \nu \lambda }^\rho `$, so that
$$D\eta ^\mu =d\eta ^\mu +\mathrm{\Gamma }_{\nu \lambda }{}_{}{}^{\mu }\eta _{}^{\lambda }dx^\nu ,D\pi _\mu =d\pi _\mu \mathrm{\Gamma }_{\nu \mu }{}_{}{}^{\lambda }\pi _{\lambda }^{}dx^\nu .$$
(13)
The corresponding Poisson brackets (which are calculated in appendix A) are:
$`\{p_\nu ,x^\mu \}`$ $`=`$ $`\delta _\nu ^\mu ,\{p_\mu ,p_\nu \}=R_{\mu \nu \lambda }{}_{}{}^{\kappa }\pi _{\kappa }^{}\eta ^\lambda `$
$`\{p_\mu ,\eta ^\nu \}`$ $`=`$ $`\mathrm{\Gamma }_{\mu \lambda }{}_{}{}^{\nu }\eta _{}^{\lambda },\{p_\mu ,\pi _\lambda \}=\mathrm{\Gamma }_{\mu \lambda }{}_{}{}^{\nu }\pi _{\nu }^{},`$
$`\text{and}\{\pi _\nu ,\eta ^\mu \}`$ $`=`$ $`\delta _\nu ^\mu .`$ (14)
the others being zero. To quantize, we take wave functions to be functions $`\psi (x,\eta )`$ on the super configuration space $`SM`$. The observables $`x^\mu `$ and $`\eta ^\mu `$ are simply represented by multiplication by these variables, while the momenta $`p_\mu `$ and $`\pi _\mu `$ are represented as
$$p_\mu =iD_\mu i\left(\frac{}{x_\mu }+\eta ^\nu \mathrm{\Gamma }_{\mu \nu }{}_{}{}^{\lambda }\frac{}{\eta ^\lambda }\right)\text{and}\pi _\mu =i\frac{}{\eta ^\mu }.$$
(15)
The BRST operator $`\mathrm{\Omega }`$ is constructed from the constraints in the standard way, giving
$$\mathrm{\Omega }=\eta ^\mu T_\mu =i\eta ^\mu \left(\frac{}{x^\mu }+v_\mu \right).$$
(16)
(The symmetry of the Christoffel symmetry removes the covariant part of $`p_\mu `$ in this case, as in exterior differentiation of forms.) The gauge-fixing fermion $`\chi `$ is then constructed from the gauge-fixing functions $`X^\mu `$ in the standard way:
$$\chi =\pi _\mu x^\mu =ig^{\mu \nu }\pi _\mu (D_\nu v_\nu ).$$
(17)
States of the system can of course naturally be identified with forms on $`M`$ via the identification
$$a_{\mu _1\mathrm{}\mu _p}(x)\eta ^{\mu _1}\mathrm{}\eta ^{\mu _p}a_{\mu _1\mathrm{}\mu _p}(x)dx^{\mu _1}\mathrm{}dx^{\mu _p}$$
(18)
Under this identification we see that
$`\mathrm{\Omega }`$ $`=`$ $`i\eta ^\mu \left({\displaystyle \frac{}{x_\mu }}v_\mu \right)=i(d+\eta ^\mu v_\mu )=ie^hde^h,`$
$`\chi `$ $`=`$ $`ig^{\mu \nu }\pi _\nu \left({\displaystyle \frac{}{x_\mu }}+\eta ^\nu \mathrm{\Gamma }_{\mu \nu }{}_{}{}^{\lambda }{\displaystyle \frac{}{\eta ^\lambda }}+v_\mu \right)=\delta ig^{\mu \nu }v_\mu \pi _\mu =e^h\delta e^h`$ (19)
where $`d`$ is exterior differentiation and $`\delta `$ is the adjoint operator to $`d`$, that is $`\delta ={}_{}{}^{}d_{}^{}`$ with the Hodge star operator. Thus we see that $`\mathrm{\Omega }`$ and $`\chi `$ are the supersymmetry operators used by Witten in his study of supersymmetry and Morse theory . (The identification of states with forms also leads to a natural inner product on states; conventions for this may be found in appendix B.)
These expressions for $`\mathrm{\Omega }`$ and $`\chi `$ simplify the calculation of the explicit expression for the canonical BRST Hamiltonian $`H=\frac{i}{2}[\mathrm{\Omega },\chi ]`$, leading to
$$H=\frac{1}{2}(d+\delta )^2+\frac{1}{2}g^{\mu \nu }\frac{h}{x^\mu }\frac{h}{x^\nu }+\frac{i}{2}g^{\mu \lambda }(\eta ^\nu \pi _\mu \pi _\mu \eta ^\nu )\frac{D^2h}{Dx^\lambda Dx^\nu }$$
(20)
which is (up to a factor $`\frac{1}{2}`$) the Hamiltonian used by Witten . Witten also shows that the mapping $`\psi e^h\psi `$ induces an isomorphism of de Rham cohomology classes of $`d`$ and $`\mathrm{\Omega }=e^hde^h`$, and that forms with zero $`H`$ eigenvalue give exactly one representative of each $`\mathrm{\Omega }`$ cohomology class. Since additionally $`H`$ has the same eigenvalues as the Laplacian $`d+\delta `$, we see that the gauge-fixing fermion $`\chi `$ is a good one .
## 4 Path Integrals
In this section stochastic calculus is used to derive a rigorous path integral expression for the action of the evolution operator $`\mathrm{exp}Ht`$ on the states of the system.
In the special case of flat space (that is, where the manifold $`M`$ is simply $`^n`$ with the Euclidean metric $`g^{\mu \nu }=\delta ^{\mu \nu }`$) this calculation was done by Salmonson and van Holten using WKB methods (also known as instanton methods), which are a standard approximation technique in quantum mechanics, . (An an accessible account of the method applied to instantons may be found in the lectures of Coleman .) The basic idea is to consider fluctuations about the classical trajectories. In the conventional WKB approach only second order fluctuations are considered (first order ones vanishing because the expansion is about the classical trajectory) so that the method used is an approximate one; in this paper we give an exact path integral formula in which the usual WKB factor appears along with further factors. The approach is valid on a manifold with a general Riemannian metric as well as in flat space.
The stochastic calculus calculations which will now be given show plainly how this arises. We will begin by working in flat space, where the Hamiltonian is
$$H_f=\frac{1}{2}(d+\delta )^2+\frac{1}{2}\frac{h}{x^\mu }\frac{h}{x^\mu }+\frac{i}{2}(\eta ^\nu \pi _\mu \pi _\mu \eta ^\nu )\frac{^2h}{x^\mu x^\nu }.$$
(21)
One simple step will then adapt the method to a general Riemannian manifold.
The starting point is the stochastic differential equation
$$dx_t^\mu =db_tv_\mu (x_t)dt,x_0=x$$
(22)
which implements the WKB approach of taking fluctuations about the classical trajectories, and corresponds to the Nicolai map . Here $`x_t^\mu `$ is a stochastic process on the Wiener space of paths in $`^n`$ starting from the point $`x`$ and $`b_t^\mu `$ is a standard Brownian path in $`^n`$. Next, for positive $`t`$, we consider the operator $`U_t`$ defined on functions on the superspace $`^{n,n}`$ by
$`U_t\psi (x,\eta )=`$
$`{\displaystyle ๐\mu \left[\mathrm{exp}\left(_0^t\left(_\mu h(x_s)dx_s^\mu +_\mu _\nu h(x_s)i\theta _s^\mu \rho _{s\nu }ds\right)\right)\psi (x_t,\theta _t)\right]}`$
where $`d\mu `$ denotes Wiener measure for paths $`(b_t^\mu ,\theta _s^\mu ,\rho _{s\mu })`$ in superspace $`^{n,2n}`$ (for the fermionic paths $`\theta _s`$ and $`\rho _s`$ see ).
Now by Itรด calculus,
$`d\left[\mathrm{exp}\left({\displaystyle _0^t}\left(_\mu h(x_s)dx_s^\mu +_\mu _\nu h(x_s)i\theta _s^\mu \rho _{s\nu }ds\right)\right)\psi (x_t,\theta _t)\right]`$
$`=[\mathrm{exp}`$ $`\left({\displaystyle _0^t}(_\mu h(x_s)dx_s^\mu +_\mu _\nu h(x_s)i\theta _s^\mu \rho _{s\nu }ds)\right)(H_f)\psi (x_t,\theta _t)]dt`$ (24)
$`+\text{terms of zero measure},`$
so that
$$\frac{U_tf(x)}{t}=U_tH_ff(x)$$
(25)
and we conclude that
$$U_t=\mathrm{exp}tH_f.$$
(26)
Now, again by Itรด calculus,
$$_0^t\left(_\mu h(x_s)dx_s^\mu +\frac{1}{2}_\mu _\mu h(x_s)ds\right)=h(x_t)h(x)$$
(27)
so that we can simplify (4) to obtain the Feynman-Kac-Itรด formula
$`\mathrm{exp}tH\psi (x,\eta )={\displaystyle ๐\mu \mathrm{exp}\left((h(x)h(x_t))\right)}`$ (28)
$`\mathrm{exp}\left({\displaystyle _0^t}_\mu _\nu h(x_s)i\theta _s^\mu \rho _{s\nu }ds\right)\psi (x_t,\theta _t).`$
This expression shows how the WKB factor $`\mathrm{exp}\left(\mathrm{\Delta }h\right)`$ (which clearly corresponds to the classical action (2) of the original topological theory) appears in the path integral.
The Feynman-Kac-Itรด formula is easily adapted to curved space with a general Riemannian metric $`g^{\mu \nu }`$ by replacing the Euclidean Brownian paths $`b_t`$ with the standard Brownian paths $`\stackrel{~}{b}_t`$ on a Riemannian manifold, and adjusting the fermion paths by using the (stochastic) vielbein $`e_{a,s}^\mu `$ as specified below. The bosonic Brownian paths on a Riemannian manifold, which were introduced by Elworthy and by Ikeda and Watanabe , are defined by the stochastic differential equations
$`d\stackrel{~}{b}_t`$ $`=`$ $`e_{a,t}^\mu db_t^a+{\displaystyle \frac{1}{2}}\mathrm{\Gamma }_{\nu \rho }{}_{}{}^{\mu }(\stackrel{~}{b}_t)dt`$
$`e_{a,t}^\mu `$ $`=`$ $`\mathrm{\Gamma }_{\nu \lambda }{}_{}{}^{\mu }(\stackrel{~}{b}_t)db_t^b+{\displaystyle \frac{1}{2}}e_{a,t}^\nu R_\nu {}_{}{}^{\mu }(\stackrel{~}{b}_t)dt`$ (29)
$`\stackrel{~}{b}_0^\mu =x^\mu ,e_{a,0}^\mu =e_a^\mu (x)`$
where $`x`$ is the point on the manifold from which the Brownian motion is chosen to start and $`\{e_a=e_a^\mu (x)\frac{}{x^\mu },a=1\mathrm{}n\}`$ is a choice of orthonormal basis at that point. The fermionic paths $`\stackrel{~}{\theta }_t^\mu ,\stackrel{~}{\rho }_{t,\nu }`$ are obtained from the flat space fermionic paths by rotating with the stochastic vielbein:
$`\stackrel{~}{\theta }_t^\mu `$ $`=`$ $`\theta _t^ae_{a,t}^\mu `$
$`\stackrel{~}{\rho }_{t,\nu }`$ $`=`$ $`\stackrel{~}{\rho }_{at}e_{a,t}^\mu g_{\nu \mu }(\stackrel{~}{b}_t).`$ (30)
Using paths $`\stackrel{~}{x}_t`$ on $`M`$ satisfying
$`d\stackrel{~}{x}_t^\mu `$ $`=`$ $`d\stackrel{~}{b}_tg^{\nu \mu }(\stackrel{~}{x}_t)v_\mu (\stackrel{~}{x}_t)dt`$
$`\stackrel{~}{x}_0`$ $`=`$ $`x,`$ (31)
similar steps to those above lead to the Feynman-Kac-Itรด formula
$`\mathrm{exp}tH\psi (x,\eta )={\displaystyle ๐\mu \mathrm{exp}((h(x)h(\stackrel{~}{x}_t)))}`$ (32)
$`\mathrm{exp}({\displaystyle _0^t}(D_\mu D_\nu h(\stackrel{~}{x}_s)ig^{\lambda \nu }(\stackrel{~}{x}_s)\stackrel{~}{\theta }_s^\mu \stackrel{~}{\rho }_{s\lambda }`$
$`+`$ $`R_\mu {}_{}{}^{\nu }(\stackrel{~}{x}_s)\stackrel{~}{\theta }_s^\mu \stackrel{~}{\rho }_{\nu s}+\frac{1}{2}R_{\mu \kappa }^{\lambda \nu }(\stackrel{~}{x}_s)\stackrel{~}{\theta }_s^\mu \stackrel{~}{\theta }_s^\kappa \stackrel{~}{\rho }_{\lambda s}\stackrel{~}{\rho }_{\nu s})ds)\psi (\stackrel{~}{x}_t,\stackrel{~}{\theta }_t).`$
In both cases care must be taken when $`x`$ is a critical point, since there will not in general be a unique solution to the stochastic differential equation concerned.
## 5 Morse theory and cohomology
In Witten rescales the function $`h`$ by a constant factor (here called $`u`$) to obtain the scaled Hamiltonian
$$H_u=\frac{1}{2}(d+\delta )^2+u^2g^{\mu \nu }\frac{h}{x^\mu }\frac{h}{x^\nu }+u\frac{i}{2}g^{\mu \lambda }(\eta ^\nu \pi _\mu \pi _\mu \eta ^\nu )\frac{D^2h}{Dx^\lambda Dx^\nu },$$
(33)
and, taking the large $`u`$ limit, builds an explicit model for the cohomology of the manifold in terms of the critical points of the manifold with differential derived from the exterior derivative. This leads directly to the weak and strong Morse identities for $`M`$, but also gives considerably further insight into the mechanism relating the critical points to the manifold topology. Some parts of Wittenโs analysis have been proved rigorously, (for instance by Bismut and by Simon et al ); however the explicit modelling of the manifoldโs cohomology via critical points and instanton calculations does not appear to have received a full mathematical treatment of the nature given below.
For the rest of this paper it will be assumed that $`h`$ is a Morse function, that is, it has only isolated critical points. (For terminology and notation see appendix C.) If $`a`$ is a critical point of $`h`$ with index $`p`$ then Witten shows that for large $`u`$ there is exactly one $`p`$-form $`\psi _{(u),a}(x,\eta )`$ on $`M`$ which is concentrated near $`a`$ and is an eigenstate of $`H_u`$ with eigenvalue $`\lambda _a(u)`$ which is low, that is, which does not grow like $`u`$ but is $`o(u)`$. (This result is derived analytically by Simon et al in .) Additionally it is shown that there are no other forms which have low $`H_u`$ eigenvalues. Witten also shows (as was remarked in section 3 for the $`u=1`$ case) that the mapping $`\psi e^{hu}\psi `$ induces an isomorphism of de Rham cohomology classes of $`d`$ and $`d_u=e^{uh}de^{uh}`$, and that forms with zero $`H_u`$ eigenvalue give exactly one representative of each $`d_u`$ cohomology class.
Observing that $`d_u`$ and $`H_u`$ commute, we see that if $`\psi `$ is an eigenstate of $`H_u`$, then $`d_u\psi `$ is either zero or an eigenstate of $`H_u`$ with the same eigenvalue. Thus if $`a`$ is a critical point of $`h`$ with index $`p`$,
$$d_u\psi (a)=\underset{bC_h,\mathrm{index}\mathrm{of}b=p+1}{}c_{ab}\psi _b$$
(34)
for some real numbers $`c_{ab}`$. As a result the cohomology of $`M`$ can be modelled by $`p`$-cochains (with $`p=1,\mathrm{},n=\mathrm{dim}M`$) of the form
$$c=\underset{aC_h,\mathrm{index}\mathrm{of}a=p}{}c_a\psi _{(u),a}$$
(35)
where the coefficients $`c_a`$ are real numbers, with the modified exterior derivative $`d_u`$ as coboundary operator. The calculation of $`c_{ab}`$ in Wittenโs paper is done by instanton methods, which may be made both more rigorous and more transparent by using the path integral expression for $`\mathrm{exp}H_ut`$ developed in the preceding section, as will now be described.
The constants $`c_{ab}`$ in equation(34) may be evaluated by considering the matrix elements $`d_{u(2)}\mathrm{exp}H_ut(A,B)`$ in the large $`u`$ limit. (Here the notation $`d_{u(2)}`$ means that the operator $`d_u`$ acts with respect to the second argument.) To see that these matrix elements are relevant, we choose an orthonormal basis of eigenstates of $`H_u`$ consisting of the low eigenvalue states $`\psi _c,cC_h`$ (where we have simplified the notation by dropping explicit reference to $`u`$) together with further eigenstates $`\{\psi _n|n=0,\mathrm{},\mathrm{}\}`$ with eigenvalues $`\lambda _n(u)`$ which will be at least of order $`u`$. We can then express the matrix elements of the evolution operator as
$$\mathrm{exp}H_ut(Y,X)=\underset{cC_h}{\overset{\mathrm{}}{}}e^{\lambda _c(u)t}{}_{}{}^{}\psi _{c}^{}(Y)\psi _c(X)+\underset{n=0}{\overset{\mathrm{}}{}}e^{\lambda _nt}{}_{}{}^{}\psi _{n}^{}(Y)\psi _n(X).$$
(36)
For large $`u`$ we have an approximate expression
$$\mathrm{exp}H_ut(Y,X)=\underset{cC_h}{\overset{\mathrm{}}{}}{}_{}{}^{}\psi _{c}^{}(Y)\psi _c(X),$$
(37)
so that we have at leading order for large $`u`$
$$d_{u(2)}\mathrm{exp}H_ut(Y,X)=\underset{cC_h}{\overset{\mathrm{}}{}}{}_{}{}^{}\psi _{c}^{}(Y)d\psi _c(X).$$
(38)
Now as was remarked above, each $`\psi _c`$ is concentrated around $`c`$. We thus expect that at leading order for large $`u`$
$$d_{u(2)}\mathrm{exp}H_ut(A,B)=c_{ab}{}_{}{}^{}\psi _{a}^{}(A)\psi _b(B),$$
(39)
if $`ฯต(A)=a`$ and $`ฯต(B)=b`$.
In order to evaluate this expression we make use of the Feynman-Kac-Itรด formula (32), together with the explicit form of the kernel in the neighbourhood of a critical point. We choose a metric which globally satisfies the Smayle transversality condition for $`h`$ (see appendix C), and additionally one which is Euclidean within the Morse coordinate neighbourhood $`N_a`$ of each critical point $`a`$ and on a neighbourhood of each steepest descent curve $`\mathrm{\Gamma }_{ab}`$ joining critical points.
Before proceeding further it is useful to introduce some specific coordinate systems. For each critical point $`c`$ in $`M`$ we will choose on $`N_c`$ a fiducial set of Morse coordinates (appendix C) $`x_{[c]}^\mu `$ and fermionic partners $`\eta _{[c]}^\mu `$. Additionally for each steepest descent curve $`\mathrm{\Gamma }_{ab}`$ (satisfying (55)) joining the pair of critical points $`a`$ and $`b`$ with indices $`p`$ and $`p+1`$ respectively we will choose a coordinate neighbourhood $`U_{_{\mathrm{\Gamma }_{ab}}}`$ which contains $`N_aN_bU_{_{\mathrm{\Gamma }_{ab}}}`$ with coordinates $`x{}_{_{\mathrm{\Gamma }_{ab}}}{}^{},\eta _{_{\mathrm{\Gamma }_{ab}}}`$ such that $`\mathrm{\Gamma }_{ab}`$ lies along $`x_{_{\mathrm{\Gamma }_{ab}}}^{np}`$ while $`x{}_{_{\mathrm{\Gamma }_{ab}}}{}^{},\eta _{_{\mathrm{\Gamma }_{ab}}}`$ match $`x{}_{[b]}{}^{},\eta _{[b]}`$ on $`N_b`$ apart from possible rotations, and also match $`x_{[a]}`$ on $`N_a`$ apart from possible rotations and (necessarily) a translation in the $`x^{np}`$ coordinate with $`x^{np}{}_{_{\mathrm{\Gamma }_{ab}}}{}^{}=x{}_{[a]}{}^{}{}_{}{}^{np}+k_a`$ for some positive constant $`k_a`$. Reconciliation with the fiducial coordinates will ultimately bring in sign factors.
Within $`N_a`$ the Hamiltonian then has the form
$`H_u`$ $`=`$ $`{\displaystyle \underset{\mu =1}{\overset{n}{}}}[\frac{1}{2}(\frac{^2}{x{}_{_{\mathrm{\Gamma }_{ab}}}{}^{\mu }x_{_{\mathrm{\Gamma }_{ab}}}^\mu }+u^2(x{}_{_{\mathrm{\Gamma }_{ab}}}{}^{\mu }a{}_{_{\mathrm{\Gamma }_{ab}}}{}^{\mu })(x{}_{_{\mathrm{\Gamma }_{ab}}}{}^{\mu }a{}_{_{\mathrm{\Gamma }_{ab}}}{}^{\mu }))`$ (40)
$`+\frac{i}{2}u\sigma _\mu (\eta {}_{_{\mathrm{\Gamma }_{ab}}}{}^{\mu }\pi {}_{_{\mathrm{\Gamma }_{ab}}}{}^{}{}_{\mu }{}^{}\pi {}_{_{\mathrm{\Gamma }_{ab}}}{}^{}{}_{\mu }{}^{}\eta {}_{_{\mathrm{\Gamma }_{ab}}}{}^{\mu })],`$
where $`\sigma _\mu =1,\mu =1,\mathrm{}np`$ while $`\sigma _\mu =1,\mu =np+1,\mathrm{}n`$. The bosonic and fermionic parts commute so that their heat kernels may be considered separately; the bosonic part is the Harmonic oscillator Hamiltonian whose heat kernel is given by Mehlerโs formula , while the fermionic part is (apart from sign factors $`\sigma _\mu `$) the fermionic oscillator whose heat kernel is given in . If $`x`$ is near $`a`$ and in $`N_a`$ then at leading order for large $`u`$
$`\mathrm{exp}H_ut(A,X)=_{\mathrm{def}}M(A,X)`$
$`=`$ $`\left({\displaystyle \frac{u}{\pi }}\right)^{n/2}\mathrm{exp}(\frac{1}{2}u(x{}_{_{\mathrm{\Gamma }_{ab}}}{}^{}k_a)^2){\displaystyle \underset{\mu =1}{\overset{np}{}}}(\alpha {}_{_{\mathrm{\Gamma }_{ab}}}{}^{\mu }){\displaystyle \underset{\nu =np+1}{\overset{n}{}}}\eta _{_{\mathrm{\Gamma }_{ab}}}^\nu `$
where $`X`$ is a point over $`x`$ in $`N_a`$, with coordinates $`(x{}_{_{\mathrm{\Gamma }_{ab}}}{}^{},\eta {}_{_{\mathrm{\Gamma }_{ab}}}{}^{})`$.
Next we calculate $`\mathrm{exp}H_ut(A,X)`$ for $`x`$ near $`\mathrm{\Gamma }_{ab}`$ using the Feynman-Kac-Itรด formula (32). In this case the steepest descent curve (satisfying (55)) from $`x`$ approaches $`a`$ very fast. Thus after very small time $`\delta t`$ the path $`\stackrel{~}{x}_{\delta t}`$ is almost certainly near $`a`$, so that to leading order in $`u`$ we have a contribution from $`\mathrm{\Gamma }_{ab}`$ of
$`\mathrm{exp}H_ut(A,X){}_{_{\mathrm{\Gamma }_{ab}}}{}^{}=\mathrm{exp}H_u\delta t\mathrm{exp}H_u(t\delta t)(A,X)`$
$`=`$ $`{\displaystyle }d\mu [\mathrm{exp}(u(h(x)h(\stackrel{~}{x}_{\delta t})))`$
$`\times `$ $`\mathrm{exp}({\displaystyle _0^{\delta t}}uD_\mu D_\nu h(\stackrel{~}{x}_s)ig^{\lambda \nu }(\stackrel{~}{x}_s)\stackrel{~}{\theta }_s^\mu \stackrel{~}{\rho }_{s\lambda }`$
$`+R_\mu {}_{}{}^{\nu }(\stackrel{~}{x}_s)\stackrel{~}{\theta }_s^\mu \stackrel{~}{\rho }_{\nu s}+{\displaystyle \frac{1}{2}}R_{\mu \kappa }^{\lambda \nu }(\stackrel{~}{x}_s)\stackrel{~}{\theta }_s^\mu \stackrel{~}{\theta }_s^\kappa \stackrel{~}{\rho }_{\lambda s}\stackrel{~}{\rho }_{\nu s}ds)M(a{}_{_{\mathrm{\Gamma }_{ab}}}{}^{},\alpha {}_{_{\mathrm{\Gamma }_{ab}}}{}^{},\stackrel{~}{x}_{\delta t},\stackrel{~}{\theta }_{\delta t})]`$
$`=`$ $`\left({\displaystyle \frac{u}{\pi }}\right)^{n/2}\mathrm{exp}(u(h(x)h(a))){\displaystyle \underset{\mu =1}{\overset{np}{}}}(\alpha {}_{_{\mathrm{\Gamma }_{ab}}}{}^{\mu }){\displaystyle \underset{\nu =np+1}{\overset{n}{}}}\eta {}_{_{\mathrm{\Gamma }_{ab}}}{}^{\nu }.`$
Here we have used the fact that the operator $`\eta ^\mu \pi _\mu `$ which corresponds to the term $`g^{\lambda \nu }(\stackrel{~}{x}_s)\stackrel{~}{\theta }_s^\mu \stackrel{~}{\rho }_{s\lambda }`$ in the path integral has zero eigenvalue on $`\mathrm{exp}H_ut(A,\stackrel{~}{x}_{\delta t},\stackrel{~}{\theta }_{\delta t})`$ when $`x`$ lies on $`\mathrm{\Gamma }_{ab}`$.
To calculate $`d_{u(2)}\mathrm{exp}H_ut(A,B)`$ we cannot take the derivative of the separate contributions from each $`\mathrm{\Gamma }_{ab}`$ using (5) because as we vary $`x`$ around $`b`$ we will jump from one $`\mathrm{\Gamma }_{ab}`$ to another. To avoid this difficulty we note that
$`d_{u(2)}\mathrm{exp}H_ut(A,B)`$ (43)
$`=`$ $`d_{u(2)}\mathrm{exp}H_us\mathrm{exp}H_u(ts)(A,B)`$
$`=`$ $`{\displaystyle _M}d^nxd^n\eta \mathrm{exp}H_u(ts)(A,X)d_{u(2)}\mathrm{exp}H_us(X,B).`$
Because of the concentration of $`d_{u(2)}\mathrm{exp}Hs(X,B)`$ near $`b`$ we can integrate over $`^n`$ rather than $`M`$ using the form of $`\mathrm{exp}Hs(X,B)`$ which is approximately true for large $`u`$ on $`N_b`$; although ultimately we will obtain a result independent of $`s`$ and $`t`$, at this stage we must use Mehlerโs formula in full (including terms of order $`e^{us}`$ whose equivalent we could neglect near $`a`$ for our purposes) because it is not the zero mode of $`H_u`$ which will contribute to $`d\psi _a(b)`$ at leading order. Thus for $`x`$ and $`y`$ near $`b`$ we use
$`\mathrm{exp}H_us(X,Y)=\left({\displaystyle \frac{u}{\pi }}\right)^{n/2}\mathrm{exp}\left(\frac{1}{2}u\left(x{}_{_{\mathrm{\Gamma }_{ab}}}{}^{2}{\displaystyle \frac{\mathrm{cosh}us}{\mathrm{sinh}us}}\right)+u{\displaystyle \frac{x{}_{_{\mathrm{\Gamma }_{ab}}}{}^{}y_{_{\mathrm{\Gamma }_{ab}}}}{\mathrm{sinh}us}}\right)`$
$`\times {\displaystyle \underset{\mu =1}{\overset{np1}{}}}(\varphi {}_{_{\mathrm{\Gamma }_{ab}}}{}^{\mu }e_{}^{us}\eta {}_{_{\mathrm{\Gamma }_{ab}}}{}^{\mu }){\displaystyle \underset{\nu =np}{\overset{n}{}}}\left(\varphi {}_{_{\mathrm{\Gamma }_{ab}}}{}^{\nu }\eta {}_{_{\mathrm{\Gamma }_{ab}}}{}^{\nu }e_{}^{us}\right).`$
where and $`(x{}_{_{\mathrm{\Gamma }_{ab}}}{}^{},\eta {}_{_{\mathrm{\Gamma }_{ab}}}{}^{})`$, $`(y{}_{_{\mathrm{\Gamma }_{ab}}}{}^{},\varphi {}_{_{\mathrm{\Gamma }_{ab}}}{}^{})`$ are the coordinates of $`X`$ and $`Y`$ respectively, so that at leading order in $`u`$ the relevant term of $`d_{u(2)}\mathrm{exp}H_us(X,B)`$ (that is, the term which contains $`d\psi _a(b)`$)is
$$\left(\frac{u}{\pi }\right)^{n/2}u\frac{x^{np}}{\mathrm{sinh}us}\mathrm{exp}(\frac{1}{2}u\left(x{}_{_{\mathrm{\Gamma }_{ab}}}{}^{2}\frac{\mathrm{cosh}us}{\mathrm{sinh}us}\right))\underset{\mu =1}{\overset{np}{}}\eta {}_{_{\mathrm{\Gamma }_{ab}}}{}^{\mu }e_{}^{us}\underset{\nu =np}{\overset{n}{}}\beta {}_{_{\mathrm{\Gamma }_{ab}}}{}^{\nu }.$$
(45)
Using (43) we see that
$`d_{u(2)}\mathrm{exp}H_ut(A,B)`$ (46)
$`=`$ $`{\displaystyle _^n}d^nx{}_{_{\mathrm{\Gamma }_{ab}}}{}^{}\left({\displaystyle \frac{u}{\pi }}\right)_{}^{n}\theta (x{}_{_{\mathrm{\Gamma }_{ab}}}{}^{np})u{\displaystyle \frac{e^{us}}{\mathrm{sinh}us}}x{\displaystyle {}_{_{\mathrm{\Gamma }_{ab}}}{}^{np}\underset{\mu =1}{\overset{np}{}}}(\alpha {}_{_{\mathrm{\Gamma }_{ab}}}{}^{\mu }){\displaystyle \underset{\nu =np}{\overset{n}{}}}\beta _{_{\mathrm{\Gamma }_{ab}}}^\nu `$
$`\times \mathrm{exp}\left(\frac{1}{2}ux{}_{_{\mathrm{\Gamma }_{ab}}}{}^{2}{\displaystyle \frac{\mathrm{cosh}us}{\mathrm{sinh}us}}\right)\mathrm{exp}u(h(x)h(a))`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}dx{}_{_{\mathrm{\Gamma }_{ab}}}{}^{np}\left({\displaystyle \frac{u}{\pi }}\right)_{}^{(n+1)/2}u{\displaystyle \frac{e^{us}}{\mathrm{sinh}us}}x{\displaystyle {}_{_{\mathrm{\Gamma }_{ab}}}{}^{np}\underset{\mu =1}{\overset{np}{}}}(\alpha {}_{_{\mathrm{\Gamma }_{ab}}}{}^{\mu }){\displaystyle \underset{\nu =np}{\overset{n}{}}}\beta _{_{\mathrm{\Gamma }_{ab}}}^\nu `$
$`\times \mathrm{exp}(\frac{1}{2}ux{}_{_{\mathrm{\Gamma }_{ab}}}{}^{2}({\displaystyle \frac{\mathrm{cosh}us}{\mathrm{sinh}us}}1))\mathrm{exp}u(h(b)h(a))`$
$`=`$ $`\left({\displaystyle \frac{u}{\pi }}\right)^{(n+1)/2}{\displaystyle \underset{\mu =1}{\overset{np}{}}}(\alpha {}_{_{\mathrm{\Gamma }_{ab}}}{}^{\mu }){\displaystyle \underset{\nu =np}{\overset{n}{}}}\beta {}_{_{\mathrm{\Gamma }_{ab}}}{}^{\nu }\mathrm{exp}u(h(b)h(a))`$
$`\text{ at leading order in }u.`$
Here the $`\theta `$-function occurs because the contribution from $`\mathrm{\Gamma }_{ab}`$ to $`\mathrm{exp}H_u(ts)`$ is zero on the side of $`b`$ away from $`a`$. Now using equation (39) and the fact that
$${}_{}{}^{}\psi _{a}^{}(A)=\left(\frac{u}{\pi }\right)^{n/4}\underset{\mu =1}{\overset{np}{}}\alpha {}_{[a]}{}^{}{}_{}{}^{\nu },\psi _b(B)=\left(\frac{u}{\pi }\right)^{n/4}\underset{\nu =np}{\overset{n}{}}\beta _{[b]}^\nu $$
(47)
we see that
$$c_{ab}=\left(\frac{u}{\pi }\right)^{1/2}\mathrm{exp}u(h(b)h(a))\underset{\mathrm{\Gamma }_{ab}}{}(1)^{\sigma _{_{\mathrm{\Gamma }_{ab}}}}$$
(48)
where $`(1)^{\sigma _{_{\mathrm{\Gamma }_{ab}}}}`$ is a sign factor which comes from the change between the $`[a]`$ and $`[b]`$ coordinates and the $`[\mathrm{\Gamma }_{ab}]`$ coordinates.
If (once again following Witten ) we rescale each $`\psi _c`$ to $`\stackrel{~}{\psi }_c=e^{uh(c)}\psi `$, and additionally use $`\stackrel{~}{d_u}=\sqrt{\frac{\pi }{u}}d_u`$, we obtain
$$\stackrel{~}{d}_u\stackrel{~}{\psi }_a=\underset{\mathrm{\Gamma }_{ab}}{}(1)^{\sigma _{\mathrm{\Gamma }_{ab}}}\stackrel{~}{\psi }_b$$
(49)
which coincides with the geometrical approach using ascending and descending spheres.
## 6 Conclusions and further possibilities
In this paper we have carried out a full canonical quantization of the simplest topological quantum theory, the topological particle, and demonstrated precisely the way in which the quantization captures topological information. The path integral formula developed in Section 4, which implements the WKB approximation in a mathematically rigorous way, even in curved space, should be useful in other situations involving quantum tunnelling.
Recent work by Hrabak on the BRST operator for the two-dimensional topological sigma model leads (in an elegant and original way using the multi-symplectic formalism) to the supersymmetric theory considered and exploited by Witten . A novel approach to quantization in the multi-symplectic formalism has been developed by Kanatchikov which might make possible a new approach to quantization of the two dimensional model.
Appendix
## Appendix A Poisson brackets
To calculate the Poisson brackets of the canonical variables $`x^\mu `$, $`p_\mu `$, $`\eta ^\mu `$ and $`\pi _\mu `$ determined by the symplectic form (12) on the super phase space $`SPM`$ introduced in section 3 we first need the Hamiltonian vector fields of these variables. The Hamiltonian vector field $`X_f`$ of a function $`f`$ on phase space is defined by
$$X_f\iota \omega _s=df,$$
(50)
where $`\iota `$ denotes interior product. With $`\omega _s`$ defined as in (12), by inspection we see that
$`X_{x^\mu }`$ $`=`$ $`{\displaystyle \frac{}{p_\mu }}`$
$`X_{p_\mu }`$ $`=`$ $`{\displaystyle \frac{}{x^\mu }}\mathrm{\Gamma }_{\mu \nu }{}_{}{}^{\rho }\pi _{\rho }^{}{\displaystyle \frac{}{\pi _\nu }}+\mathrm{\Gamma }_{\mu \nu }{}_{}{}^{\rho }\eta _{}^{\nu }{\displaystyle \frac{}{\eta ^\rho }}{\displaystyle \frac{1}{2}}R_{\mu \nu \rho }{}_{}{}^{\lambda }\pi _{\lambda }^{}\eta ^\rho {\displaystyle \frac{}{p_\nu }}`$
$`X_{\eta ^\mu }`$ $`=`$ $`{\displaystyle \frac{}{\pi _\mu }}`$
$`X_{\pi _\mu }`$ $`=`$ $`{\displaystyle \frac{}{\eta ^\mu }}.`$ (51)
Poisson brackets are then defined by the rule
$$\{f,g\}=\frac{1}{2}\left(X_f\iota dgX_g\iota df\right)$$
(52)
which leads to (3).
## Appendix B Sign conventions for integrals
For the supermanifold $`SM`$ Berezin integration corresponds to integration of top forms on $`M`$. If a function $`f`$ on $`SM`$ takes the form $`f(x,\eta )=f_{\mu _1\mathrm{}\mu _p}(x)\eta ^{\mu _1}\mathrm{}\eta ^{\mu _p}`$ then the conventional integral
$$_Mf_{\mu _1\mathrm{}\mu _p}(x)๐x^{\mu _1}\mathrm{}๐x^{\mu _p}$$
is equal to the Berezin integral
$$_{SM}d^nxd^n\eta f_{\mu _1\mathrm{}\mu _p}(x)\eta ^{\mu _1}\mathrm{}\eta ^{\mu _p}.$$
If $`K`$ is a linear operator on functions on the supermanifold $`SM`$, then the integral kernel of $`K`$ (if it exists) is defined by
$$Kf(Y)=_{SM}d^nxd^n\eta f(X)K(X,Y).$$
(53)
## Appendix C Morse Theory terminology
Collected together here are some basic definitions and notation for Morse theory; more details, and many more results, may be found in the classic book of Milnor .
We start with a function $`h:M`$. The critical points of $`h`$ are the points where all the partial derivatives are zero. In this paper we will assume that $`h`$ is a Morse function, that is, its critical points are all isolated; a critical point $`a`$ of $`h`$ is said to be of index $`p`$ if the Hessian matrix $`\left(D_{x^\mu }D_{x^\nu }h(a)\right)`$ has exactly $`p`$ negative eigenvalues. The set of critical points of $`h`$ will be denoted $`C_h`$. (Although we give here a coordinate-based definition of critical point and index, the definitions are of course intrinsic, independent of any choice of local coordinates.)
Each critical point $`a`$ of $`h`$ has a neighbourhood $`N_a`$ on which special coordinate systems, known as Morse coordinates, can be chosen in which the Morse function $`h`$ takes the standard form
$$h(x)=h(a)+\frac{1}{2}\underset{\mu =1}{\overset{n}{}}\sigma _\mu (x^\mu a^\mu )^2$$
(54)
with $`\sigma _\mu =+1`$ for $`\mu =1,\mathrm{},np`$ and $`\sigma _\mu =1`$ for $`\mu =np+1,\mathrm{},n`$. (Here, abusing notation for simplicity, points and their coordinates are identified.) On the corresponding neighbourhood of the supermanifold $`SM`$ we use odd coordinates $`\kappa _\mu `$ corresponding to $`dx^\mu `$.
A metric $`g`$ on $`M`$ satisfies the Smayle transversality condition for $`h`$ if the solution curves $`\mathrm{\Gamma }_{ab}`$ to the โsteepest descentโ differential equation
$$\frac{dx^\mu (t)}{dt}=g^{\mu \nu }\frac{h}{x^\nu }$$
(55)
which start from a critical point $`b`$ and end at a critical point $`a`$ (with $`h(a)`$ necessarily less than $`h(b)`$) are discrete (and finite in number).
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# Acknowledgments
## Acknowledgments
The author wishes to thank A.I. Studenikin and the Organizing Committee of the 9th Lomonosov Conference on Elementary Particle Physics for invitation to participate and give a talk.
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# Null surfaces formulation in 3D
## I Introduction
The Null Surfaces Formulation (NSF) developed by Frittelli, Kozameh, and Newman gives an alternative treatment of Einsteinโs General Relativity (GR), which emphasizes the dynamical role played in the theory by the null hypersurfaces of the space-time metric. The construction is set up ab initio in terms of generic space-time foliations. These are represented by the level surfaces of a two-parametric function $`Z(x^a;\zeta ,\overline{\zeta })`$ of the $`x^a`$ coordinates, on which it is imposed the condition that the foliations become null for some space-time metric $`g_{ab}(x^a)`$ to be obtained as a functional of $`Z(x^a;\zeta ,\overline{\zeta })`$. Then, an auxiliary function $`\mathrm{\Omega }(x^a;\zeta ,\overline{\zeta })`$ fixes the resulting conformal freedom in such a way that the Einsteinโs field equations be satisfied for $`g_{ab}(x^a)`$. Such equations are expressed entirely in terms of the non-local variables $`Z(x^a;\zeta ,\overline{\zeta })`$ and $`\mathrm{\Omega }(x^a;\zeta ,\overline{\zeta })`$, with no mention at all of the metric and its associated tensors, which are merely functions of them.
One important issue in the NSF is how to restrict the freedom of (infinitely many) different families of null hypersurfaces corresponding to the same metric. For asymptotically flat space-times a canonical family of null foliations can be selected by considering past null cones from points at future null infinity, $`๐ฅ^+`$. In this case the function $`Z`$ has a second meaning. For fixed values of $`x^a`$ the function $`u=Z(x^a;\zeta ,\overline{\zeta })`$ parametrically describes the intersection of the future light cone from $`x^a`$ with $`๐ฅ^+`$. This intersection is called a light cone cut of null infinity and it plays a central role in reformulating General Relativity in terms of the free data (representing outgoing gravitational radiation) at $`๐ฅ^+`$.
It is worth mentioning that the final equations of NSF become a system of partial differential equations for the main variables that, in spite of the clear physical and geometric contents of the equations, are technically very involved and difficult to analyze. This is best exemplified in an alternative formulation of NSF that decouples the original set of equations and produces a conformally invariant equation solely for $`Z`$, a desirable feature since $`Z`$ only contains information of the conformal structure . Although the resulting equation is equivalent to the vanishing of the Bach tensor, it cannot be explicitly written in terms of $`Z`$ and its derivatives on a reasonable number of lines.
Thus, from a purely technical point of view it is desirable to derive an analogous NSF formalism in a lower dimensional space. The obvious choice is 3-D gravity since its almost trivial character should simplify the formalism and make it easier to investigate all the hoped features that the formulation in four dimensions is required to posses in order to be consistent with standard GR. Examples of these are caustics and other singularities of null congruences predicted by general theorems, which are at present searched for only in a perturbative approach to the NSF field equations .
Moreover, three-dimensional GR can be used to study many examples of 4-D vacuum space-times with a space translation symmetry. By going to the โmanifolds of orbitsโ of the associated Killing field, the problem is reduced to 3-D gravity coupled to a scalar field acting as matter source of the 3-D Einsteinโs equation. Assuming suitable fall off conditions on the scalar field, the concept of asymptotic flatness can be stated as close as possible to the analogous four dimensional treatment . Consequently, for asymptotically flat 3-D space-times, one can introduce our canonical family of null surfaces, i.e., the light cone cuts of null infinity. In this way, we extend our formalism to 4-D space-times which are not asymptotically flat in the usual sense since they include cylindrical or axisymmetric waves.
Finally, from the quantum theory point of view, there is some evidence, at least at a linear level in 4-D and at a non perturbative for some midi-superspaces in 3-D, that in contrast with the quantization in terms of the local fields like the metric, the quantization of the non-local variable $`Z`$ and $`\mathrm{\Omega }`$ of NSF, results in operators with non-diverging expectation values for semi-classical states. Concrete results have been obtained in 3-D since some midi-superspaces can be exactly solved and large quantum gravity effects appear for local fields such as the metric . Recently, it has been shown that the quantum light cone cut admits a semiclassical approximation which is stable against small perturbations at transplanckian frequencies . Thus, it appears that this non local variable is best suited for semiclassical approximations. To develop a more thorough discussion on these issues one should have available the full classical theory.
This work is organized as follows. In section two we introduce one-parameter families of hypersurfaces in 3-D and discuss which conditions should be imposed on these surfaces so that they become null. The main results of this section, the metric components and metricity conditions, were originally obtained by E. Cartan and S. Chern many years ago and recently rederived by M. Tanimoto but we follow an approach which is closely related to our work done in four dimensions. In section three we derive the field equations for these surfaces including matter sources. As in previous results one can decouple those equations and obtain a conformally invariant equation for the null hypersurfaces. In section four we review the notion of asymptotic flatness in 3-D and define the light cone cuts. Several kinematical issues which differ from the four dimensional case are pointed out and discussed. A particular example of a light cone cut is constructed and analyzed. In the conclusions we outline an approach to handle caustics and singularities that come together with null surfaces.
## II NSF formalism in three dimensions.
In this section the NSF formalism, as developed in , is adapted to dimension $`2+1`$ and the corresponding kinematical structures are constructed. Although the basic framework was outlined in detail in 4-D , it is worth repeating the construction in 3-D since it is conceptually equivalent to the higher dimensional case though technically less involved and therefore easier to follow.
We begin with a three-dimensional manifold $`M`$, and assume we are given one-parameter functions $`Z(x^a,\zeta )`$ of the space-time coordinates $`x^a`$; and parameter $`\zeta `$ on $`S^1`$ running between $`0`$ and $`2\pi `$. We also assume that for a fixed value of the parameter $`\zeta `$, the level surfaces of $`Z`$,
$$Z(x^a,\zeta )=\mathrm{const}.,$$
(1)
locally foliate the manifold $`M`$.
The statement that the level surfaces of $`Z`$ be indeed the characteristic or null surfaces of some metric $`g_{ab}(x^a)`$ for $`M`$ is that for fixed values of $`x^a`$ and arbitrary values of $`\zeta `$ they satisfy
$$g^{ab}(x^a)_aZ(x^a,\zeta )_bZ(x^a,\zeta )=0.$$
(2)
An equivalent geometrical statement is that as the families of foliations intersect at a single but arbitrary point $`x^a`$ we demand that the enveloping surface forms the light cone from the point $`x^a`$. Thus, the parameter $`\zeta `$ spans the circle of null directions.
The idea now is to solve eq. (2) for the five components of the conformal metric in terms of $`_aZ(x^a,\zeta )`$. Given an arbitrary function $`Z`$ the problem has no solution since we have an infinite number of algebraic equations (one for each value of $`\zeta `$) for five unknowns. Therefore we must impose conditions on $`Z(x^a,\zeta )`$ so that a solution exists. The solution and conditions are best expressed when written in a canonical coordinate system constructed from knowledge of $`Z`$.
To see this, we introduce three functions defined as
$$(\theta ^0,\theta ^1,\theta ^2)(u,\omega ,r)(Z(x^a,\zeta ),Z(x^a,\zeta ),^2Z(x^a,\zeta )),$$
(3)
where $`\frac{}{\zeta }`$ denotes the derivative of $`\zeta `$ holding $`x^a`$ fixed. For each value of $`\zeta `$ these functions form a coordinate system intrinsically adapted to the surfaces. Its associated gradient $`\theta ^i_a`$ and vector dual basis $`\theta _i^a`$ are given by
$$\theta ^i{}_{a}{}^{}(d\theta ^i)_a,\theta _i{}_{}{}^{a}\left(\frac{}{\theta ^i}\right)^a.$$
(4)
We now demand that $`u=const.`$ is a null surface. Then on that surface $`\omega =const.`$ singles out a null geodesic whereas $`r=const.`$ identifies a point on that geodesic. The next available scalar, $`^3Z(x^a,\zeta )`$, determines the conformal metric and is constrained to satisfy the metricity conditions. Using eq. (3) and assuming the $`\theta ^i`$ coordinate system is well behaved, one can obtain $`x^a=x^a(\theta ^i,\zeta )`$ and define $`\mathrm{\Lambda }(u,\omega ,r,\zeta )^3Z(x^a(\theta ^i,\zeta ),\zeta )`$.
A technical point worth mentioning is that the action of the operator $``$, the parameter derivative holding $`x^a`$ fixed, acting on a function $`F(u,\omega ,r,\zeta )`$ is given by
$$=_\zeta +\omega _u+r_\omega +\mathrm{\Lambda }_r.$$
(5)
Thus, the action of the $``$ operator on $`\mathrm{\Lambda }`$ is given by (5).
Likewise, $``$ does not commute with the directional derivatives $`_i\frac{}{\theta ^i}`$. The explicit form of the commutation relations, needed later, are
$$[_u,]=_u\mathrm{\Lambda }_r,$$
(6)
$$[_\omega ,]=_u+_\omega \mathrm{\Lambda }_r,$$
(7)
$$[_r,]=_\omega +_r\mathrm{\Lambda }_r.$$
(8)
The metric components and metricity conditions are obtained by repeatedly operating $``$ on ( 2). Taking $``$ on (2) yields
$$(g^{ab}_aZ_bZ)=2g^{ab}_aZ_bZ=2g^{01}=0,$$
(9)
where we have used that $`g^{ab}=0`$. If we take again $``$ on the above equation we obtain
$$g^{ab}(_aZ_b^2Z+_aZ_bZ)=g^{02}+g^{11}=0.$$
(10)
Repeating this procedure gives
$$g^{ab}(_aZ_b^3Z+3_aZ_b^2Z)=g^{02}_r\mathrm{\Lambda }+3g^{12}=0,$$
(11)
$`g^{ab}(_aZ_b^4Z+4_aZ_b^3Z+3_a^2Z_b^2Z)`$ $`=`$ $`0,`$ (12)
$`g^{02}_r\mathrm{\Lambda }+4(g^{11}_\omega \mathrm{\Lambda }+g^{12}_r\mathrm{\Lambda })+3g^{22}`$ $`=`$ $`0.`$ (13)
Using the commutator (8), together with the already obtained metric components, we can rewrite this last expression as,
$$g^{02}\left(_r\mathrm{\Lambda }3_\omega \mathrm{\Lambda }\frac{1}{3}(_r\mathrm{\Lambda })^2\right)+3g^{22}=0.$$
(14)
Note that the four nontrivial metric components are proportional to $`g^{02}`$. It is therefore natural to take $`g^{02}`$ as a $`\zeta `$-dependent conformal factor $`\mathrm{\Omega }^2`$ and write the metric tensor as $`g^{ab}=\mathrm{\Omega }^2h^{ab}=\mathrm{\Omega }^2h^{ij}\theta _i^a\theta _j^b`$ with
$$h^{ij}=\left(\begin{array}{ccc}0& 0& 1\\ 0& 1& \frac{1}{3}_r\mathrm{\Lambda }\\ 1& \frac{1}{3}_r\mathrm{\Lambda }& \frac{1}{3}(_r\mathrm{\Lambda })+\frac{1}{9}(_r\mathrm{\Lambda })^2+_\omega \mathrm{\Lambda }\end{array}\right).$$
(15)
The conformal factor $`\mathrm{\Omega }`$ cannot be an arbitrary function of $`\zeta `$ since it is defined as
$`\mathrm{\Omega }^2=g^{ab}_aZ_b^2Z.`$
Therefore,
$`\mathrm{\Omega }^2=g^{ab}(_aZ_b^2Z+_aZ_b^3Z).`$
Thus, $`\mathrm{\Omega }`$ and $`\mathrm{\Lambda }`$ must satisfy
$$3\mathrm{\Omega }=\mathrm{\Omega }_r\mathrm{\Lambda }.$$
(16)
Equation (16) is invariant under $`\mathrm{\Omega }(x,\zeta )\mathrm{\Omega }^{}(x,\zeta )=f(x)\mathrm{\Omega }(x,\zeta )`$ for an arbitrary $`f(x)`$. This freedom is a consequence of the conformal invariance of the formulation.
Note also that nothing prevents us from taking further $``$ derivatives on (12). Since the r.h.s of this equation vanishes and all the components of the conformal metric have been explicitly constructed from $`_i\mathrm{\Lambda }`$, the next $``$ on (12) must impose a condition on $`\mathrm{\Lambda }`$ so that a conformal metric exists. This equation can be written as
$`g^{ab}(5_aZ_b^4Z+_aZ_b^5Z+10_a^2Z_b^3Z)`$ $`=`$ $`0`$ (17)
$`5g^{11}_\omega \mathrm{\Lambda }+5g^{12}_r\mathrm{\Lambda }+g^{02}_r^2\mathrm{\Lambda }+10(g^{02}_u\mathrm{\Lambda }+g^{12}_\omega \mathrm{\Lambda }+g^{22}_r\mathrm{\Lambda })`$ $`=`$ $`0.`$ (18)
Using the explicit form of the metric components and the commutation relations and dividing out by the conformal factor gives
$$M[\mathrm{\Lambda }]2\left((_r\mathrm{\Lambda })_\omega \mathrm{\Lambda }\frac{2}{9}(_r\mathrm{\Lambda })^2\right)_r\mathrm{\Lambda }^2(_r\mathrm{\Lambda })+3(_\omega \mathrm{\Lambda })6_u\mathrm{\Lambda }=0.$$
(19)
Conversely, if (19) is satisfied, then further $``$ derivatives on (14) are identically satisfied.
We summarize this section as follows: given a function $`Z(x^a,\zeta )`$ such that $`\mathrm{\Lambda }(\theta ^i,\zeta )^3Z(x^a,\zeta )`$ satisfies $`M[\mathrm{\Lambda }]=0`$, then $`Z=const.`$ is a null hypersurface of $`g^{ab}`$.
Remark 1: Equation (19) is the main metricity condition, (16) is only used to fix the $`\zeta `$-dependence of the conformal factor. In fact, given a function $`Z(x^a,\zeta )`$ such that $`^3Z(x^a,\zeta )`$ satisfies $`M[\mathrm{\Lambda }]=0`$, then $`Z(x^a,\zeta )=const.`$ is a null hypersurface of $`h^{ab}(x^a,\zeta ^{})`$ for an arbitrary value of $`\zeta ^{}`$.
To see this we first show that $`h^{ab}(x^a,\zeta ^{})`$ satisfies
$$h^{ab}=\frac{2}{3}_r\mathrm{\Lambda }h^{ab},$$
(20)
which follows from the definition of $`h^{ab}`$. The above equation can be integrated immediately to give
$`h^{ab}(x^a,\zeta ^{})=\text{ exp}\left({\displaystyle \frac{2}{3}}{\displaystyle _\zeta ^\zeta ^{}}_r\mathrm{\Lambda }(x^a,\eta )d\eta \right)h^{ab}(x^a,\zeta ).`$
Thus, if $`Z(x^a,\zeta )`$ yields null surfaces for $`h^{ab}(x^a,\zeta )`$ it also gives null surfaces for any $`h^{ab}(x^a,\zeta ^{})`$ that satisfies (20). Since $`h^{ab}(x^a,\zeta )`$ is an explicit functional of $`\mathrm{\Lambda }`$ we conclude that (19) is the only condition to be imposed on $`Z`$ so that a conformal geometry can be constructed on the three dimensional manifold.
Remark 2: In fact, it is not necessary to assume that $`g^{ab}=g^{ab}(x)`$ to obtain the explicit form of the conformal metric and metricity condition (19). If the metric $`g^{ab}`$ is such that
$`g^{ab}=\mu (x,\zeta )g^{ab}`$
for an arbitrary function $`\mu (x,\zeta )`$, we can repeat the steps shown above to obtain the same conformal metric $`h^{ab}(x^a,\zeta )`$ and metricity condition (19).
Remark 3: Starting from a completely different perspective, E. Cartan and S. Chern studied the geometry of the ordinary differential equation
$`^3Z=\mathrm{\Lambda }(\zeta ,Z,Z,^2Z).`$
They found that the condition (19) has a simple geometrical meaning, namely, that the Cartan normal conformal connection defined in the solution space is torsion free. It is therefore natural in our case to obtain (19) as a metricity condition.
Remark 4: A very interesting (and open) problem in 4-D is to understand the geometrical meaning of the metricity conditions. It is natural to conjecture that the complex PDE derived in that case is also the torsion free condition for the conformal connection defined in the solution space of the starting second order PDE for $`Z`$ on the sphere of null directions.
## III The field equations.
In this section we derive the field equations for the basic variables of this formalism. In contrast with the 4-D case, the vanishing of the trace free part of the Ricci tensor does not contain interesting solutions. Since the Weyl tensor vanishes in three dimensions, the solutions of
$$R_{ab}\frac{1}{3}Rg_{ab}=0$$
yield spaces with constant curvature. Thus, we must add a stress energy tensor as a source term on the r.h.s. of the above equation. For simplicity we will assume that $`T_{ab}`$ is a given trace free tensor. If the matter source contains a trace part then the trace equation is automatically satisfied by virtue of the vanishing of the divergence of $`T_{ab}`$. This can be easily seen by writing
$`G_{ab}=\underset{ยฏ}{G}_{ab}+{\displaystyle \frac{1}{3}}g_{ab}G,`$
where $`G`$ and $`\underset{ยฏ}{G}_{ab}`$ are the trace and the trace free part of the Einstein tensor respectively. If the trace free equations
$`\underset{ยฏ}{G}_{ab}=\underset{ยฏ}{T}_{ab},`$
are satisfied, then using the vanishing of the divergence of the stress energy tensor and the Einstein tensor one obtains
$`_a(GT)=4^b(\underset{ยฏ}{G}_{ab}\underset{ยฏ}{T}_{ab})=0.`$
We thus start with the Einstein equations and contract the indices $`a,b`$ with the null vectors $`\theta _2{}_{}{}^{a}\theta _{2}^{}^b`$. Note that this operation automatically removes the trace part of the equation. Rewriting the resulting equation in terms of our variables yields
$$_r^2\mathrm{\Omega }=T_{rr}\mathrm{\Omega }.$$
(21)
At a first glance it appears that this single equation cannot be equivalent to the five components of the Einstein equations. Note, however, that (21) is true for any value of $`\zeta `$. Thus, if we add to (21) the metricity conditions which gives the change of $`\mathrm{\Omega }`$ and $`\mathrm{\Lambda }`$ under variations of $`\zeta `$ we obtain a consistent set of equations equivalent to the standard Einstein equations. The final equations read,
$$\begin{array}{cc}(\text{E})_r^2\mathrm{\Omega }\hfill & =T_{rr}\mathrm{\Omega },\hfill \\ (m_I)\mathrm{\Omega }\hfill & =\frac{1}{3}_r\mathrm{\Lambda }\mathrm{\Omega },\hfill \\ (m_{II})6_u\mathrm{\Lambda }\hfill & =\frac{4}{9}(_r\mathrm{\Lambda })^3+2_r\mathrm{\Lambda }(_r\mathrm{\Lambda })2_\omega \mathrm{\Lambda }_r\mathrm{\Lambda }^2(_r\mathrm{\Lambda })+3(_\omega \mathrm{\Lambda }).\hfill \end{array}\}$$
(22)
The above constitute a coupled set of equations for $`\mathrm{\Omega }`$ and $`\mathrm{\Lambda }`$. These two functions, however, have completely different meaning; whereas $`\mathrm{\Lambda }`$ determines the entire conformal structure, $`\mathrm{\Omega }`$ only fixes the scale. It is thus desirable to decouple the above set of equations and obtain equations that only involve $`\mathrm{\Lambda }`$. By construction, the resulting equations would be conformally invariant.
The approach followed to decouple the field equations is to study the integrability conditions for ($`m_I`$) and (E) since both equations must be satisfied by a single function $`\mathrm{\Omega }`$. Therefore, conditions are imposed on $`\mathrm{\Lambda }`$ so that a solution exists.
To illustrate the procedure and to discuss the freedom left in the solutions we first consider the case when $`T_{rr}=0`$. Using (6), (7) and (8), we calculate
$$[_r^2,]\mathrm{\Omega }=2_\omega _r\mathrm{\Omega }+2_r\mathrm{\Lambda }_r^2\mathrm{\Omega }+_r^2\mathrm{\Lambda }_r\mathrm{\Omega }.$$
(23)
On the other hand, using ($`m_I`$) and (E), we have
$$[_r^2,]\mathrm{\Omega }=_r^2(\frac{1}{3}_r\mathrm{\Lambda }\mathrm{\Omega }),$$
(24)
the r.h.s. of equations (23) and (24), gives
$$2_\omega _r\mathrm{\Omega }=\frac{1}{3}_r^3\mathrm{\Lambda }\mathrm{\Omega }\frac{1}{3}_r^2\mathrm{\Lambda }_r\mathrm{\Omega }.$$
Applying $`_r`$ to the above equation, commuting the partial derivatives, and using eq. (E) yields
$$_r^4\mathrm{\Lambda }=0.$$
(25)
Thus, for spaces with constant curvature the field equations for $`\mathrm{\Lambda }`$ are,
$`_r^4\mathrm{\Lambda }`$ $`=`$ $`0,`$ (26)
$`6_u\mathrm{\Lambda }`$ $`=`$ $`{\displaystyle \frac{4}{9}}(_r\mathrm{\Lambda })^3+2_r\mathrm{\Lambda }(_r\mathrm{\Lambda })2_\omega \mathrm{\Lambda }_r\mathrm{\Lambda }^2(_r\mathrm{\Lambda })+3(_\omega \mathrm{\Lambda }).`$ (27)
By construction, equation (26) is conformally invariant. It is thus interesting to ask what tensorial quantity is represented by $`_r^4\mathrm{\Lambda }`$. It can be easily shown that
$$\frac{1}{6}_r^4\mathrm{\Lambda }=B_{ab}Z^aZ^b,$$
with
$$B_{ab}=2ฯต_a{}_{}{}^{mn}_{[m}^{}S_{n]b},$$
$$S_{ab}=R_{ab}\frac{1}{4}Rg_{ab}.$$
$`2_{[m}S_{n]b}`$ known as the Bach tensor in the literature, was first defined and used by E. Cotton in 1899 to show that its vanishing determines the conformal structure of spaces with constant curvature (see Appendix A).
The solutions to equations (26) and (27) should contain functional degrees of freedom since any null surface compatible with a given Einstein metric should be a solution of those equations. This feature is shown below where we discuss the simplest set of equations corresponding to flat space and we obtain the collection of all null hypersurfaces in Minkowski space. For this we consider the special class of solutions to the above equations given by $`\mathrm{\Lambda }`$โs that satisfy
$$_r\mathrm{\Lambda }=0,_\omega \mathrm{\Lambda }=1,$$
since this restriction yields a Minkowski metric (see eq. (15)). We immediately see that equation (26) is identically satisfied whereas the other equation gives
$$_u\mathrm{\Lambda }=0.$$
Thus, the solution can be written as
$$\mathrm{\Lambda }+\omega =g(\zeta ),$$
where $`g`$ is an arbitrary function on $`S^1`$. Using the relationship between $`\mathrm{\Lambda }`$, $`\omega `$, and $`Z`$ we rewrite the above equation as
$$^3Z+Z=g(\zeta ),$$
whose solution contains three arbitrary constants. Denoting those constants as $`(t,x,y)`$ we write the solution as
$$Z=t+x\mathrm{cos}(\zeta )+y\mathrm{sin}(\zeta )+\alpha (\zeta ),$$
with $`^3\alpha =g(\zeta )`$. The above solution is the generating function for an arbitrary null hypersurface in flat 3-D space. Its freedom is given by an arbitrary function on the circle of null directions. In the next section we will cut down this freedom by finding a canonical family of surfaces.
Following a similar approach we obtain the field equations for $`\mathrm{\Lambda }`$ when matter source is included. Then eq. (26) is replaced by (see Appendix)
$$_r^4\mathrm{\Lambda }=6_\omega T_{rr}6_rT_{r\omega }+_r^2\mathrm{\Lambda }+6\mathrm{\Omega }^1(T_{rr}_\omega \mathrm{\Omega }T_{r\omega }_r\mathrm{\Omega }).$$
(28)
Note that $`\mathrm{\Omega }`$ appears explicitly in the above equation. Thus, it would appear that we are unable to get rid of the conformal factor. This is not so. Applying several times the $``$ operator on (28) gives five independent equations involving the three components of $`_a\mathrm{\Omega }`$. We can algebraically solve for these three components and leave two extra equations for $`\mathrm{\Lambda }`$ and the energy-momentum tensor. An explicit illustration of this procedure written in tensorial language is given in Appendix B.
## IV Light Cone Cuts of $`๐ฅ^+`$
The main results obtained in previous sections, the metricity conditions and the field equations, are valid for any family of null hypersurfaces. In this section we introduce a canonical family with a well defined geometrical meaning. They represent past null cones from future null infinity. As a necessary step before introducing this family of null hypersurfaces we first review the notion of asymptotic structure available in 3-D.
A space-time is considered to be asymptotically flat at null infinity if one can attach a boundary $`๐ฅ^+`$ with topology $`S^1\times R`$ to the original manifold, a smooth metric field $`g_{ab}^{}`$ and scalar field $`\omega `$ on $`M๐ฅ^+`$, such that:
1. on $`M`$ $`\omega 0`$, $`g_{ab}^{}=\omega ^2g_{ab}`$,
2. on $`๐ฅ^+`$ $`\omega =0`$, $`\omega _a=_a\omega 0`$, and $`g^{ab}\omega _a\omega _b=0`$,
3. if $`\omega `$ is chosen such that $`^a\omega _a=0`$ on $`๐ฅ^+`$, then the vector field $`n^a=g^{ab}\omega _b`$ is complete on $`๐ฅ^+`$.
As in 4-D, the smooth unphysical metric on $`M๐ฅ^+`$ is used to discuss the fall off behavior of matter fields near $`๐ฅ^+`$. By bringing null infinity at a finite distance, one can use the techniques of differential geometry near $`๐ฅ^+`$ instead of taking limits along null directions. Moreover, the space of null geodesics is the same for both the physical and unphysical metric, the only difference being in the corresponding affine parametrization of the geodesics.
There are, however, several differences with 4-D asymptotia. In 3-D the Weyl tensor vanishes and the metric is flat outside sources. Thus, the dynamical degrees of freedom are coded in the behavior of matter fields near $`๐ฅ^+`$.
We now turn our attention to the description of light cone cuts of $`๐ฅ^+`$ for 3-D asymptotically flat space-times. A straightforward method to obtain these cuts is to start with the future light cone from an arbitrary point $`x^a`$ of the space time and define a light cone cut of null infinity as the intersection of this cone with the null boundary. Using Bondi coordinates $`(u,\zeta )`$ at $`๐ฅ^+`$ one can locally describe this intersection as
$$u=Z(x^a,\zeta ).$$
(29)
The function $`Z`$ is called the light cone cut function. It is easy to show that $`Z`$ has a second meaning. For fixed values of $`u`$ and $`\zeta `$, the points $`x^a`$ that satisfy (29) lie on the past null cone from the point $`(u,\zeta )`$ at $`๐ฅ^+`$. Thus, the level surfaces of $`Z(x^a,\zeta )`$ form a canonical one-parameter family of null surfaces on the space time.
In the above construction one makes explicit use of the background metric to first obtain the null cone from $`x^a`$ and then intersect this cone with $`๐ฅ^+`$. Since our goal is to regard the cuts as the main variables, we search for a construction that does not involve explicit use of the metric. The idea is to define the cuts as the solutions to a third order ODE where the source term, $`\mathrm{\Lambda }`$, is obtained from field equations that satisfy the desired asymptotic conditions.
As was mentioned in the previous section one has available two sets of field equations for $`\mathrm{\Lambda }`$ and it is worth reviewing both approaches since they involve different variables. The first set of equations, namely eq. (22), are PDEโs for $`\mathrm{\Omega }`$ and $`\mathrm{\Lambda }`$ with the matter field $`T_{rr}`$ as a source term. Imposing regularity and asympotic conditions on the stress-energy tensor as discussed in together with apropriate initial conditions on $`\mathrm{\Omega }`$ and $`\mathrm{\Lambda }`$ one should be able to find solutions that preserve the asymptotic structure. For example, $`\mathrm{\Omega }`$ and $`_\omega \mathrm{\Lambda }`$ should tend to 1 and -1 respectively as we approach to $`๐ฅ^+`$. Assuming one has solutions to (22) and denoting those solutions as
$$\mathrm{\Omega }(u,\omega ,r,\zeta )=\mathrm{\Omega }[T_{rr}],\mathrm{\Lambda }(u,\omega ,r,\zeta )=\mathrm{\Lambda }[T_{rr}],$$
where the term in bracket denotes functional dependence, one should then use the obtained $`\mathrm{\Lambda }`$ to define the light cone cuts as the solutions to
$$^3Z=\mathrm{\Lambda }(Z,Z,^2Z,\zeta )=\mathrm{\Lambda }[T_{rr}],$$
(30)
where the coordinates $`(u,\omega ,r,\zeta )`$ have been replaced by $`(Z,Z,^2Z,\zeta )`$. The resulting equation is a third order ODE with a 3-dim solution space. Note that the metricity condition is automatically satisfied since it is part of the field equations for $`\mathrm{\Lambda }`$. Thus, one can algebraically obtain a metric in the solution space such that the level surfaces $`Z=const.`$ are null surfaces of that metric.
If we follow the second approach we should start with eqs.(27) and (28) and follow a similar construction to end up with (30). The implicit assumption being done here is that one is able to eliminate $`\mathrm{\Omega }`$ in (28) and obtain a set of field equations for $`\mathrm{\Lambda }`$. Although this second approach offers the advantage of having just one variable the first set of equations for $`\mathrm{\Omega }`$ and $`\mathrm{\Lambda }`$ appears to be technically more tractable and easier to handle.
At this point we mention a possible third approach that yields a different equation for the light cone cuts. We first derive an analog of the Sachsโ theorem in four dimensions (which gives a relationship between the cuts and the shears associated with the Bondi and light cone cuts ). Since the shear vanishes in three dimensions we use the only available structure that is left, the divergence of null congruences. To be more precise, we start with the divergence $`\rho `$ of a null cone with apex at $`x^a`$. This optical scalar satisfies the following ODE
$$\frac{\rho }{s}=\rho ^2+\mathrm{\Phi }_{oo},$$
(31)
where s is an affine length and $`\mathrm{\Phi }_{oo}=R_{ab}ล^aล^b`$. The null cone congruence condition is imposed by assuming that near the apex $`s=s_o`$, $`\rho ={\displaystyle \frac{1}{ss_o}}`$. Following a similar calculation outlined in one can show that at points of intersection between the light cone from $`x^a`$ and $`๐ฅ^+`$, the Bondi and light cone divergences are related through
$$^2Z=\rho _B\rho _Z,$$
(32)
where $`\rho _B`$ is the divergence of the Bondi cuts (usually taken to be zero) and $`\rho _Z`$ the solution of (31) evaluated at $`๐ฅ^+`$. Note that $`\rho _Z`$ is a functional of the Ricci tensor $`\mathrm{\Phi }_{oo}`$. Thus, we can regard (31) and (32) as the field equations for $`Z`$. Note also that the resulting equation is a second order ODE which in principle is equivalent to the third order ODE previously obtained. However, the solution space of this last equation is 2-dim and it is not clear at this point how to generate the required extra dimension. A thorough discussion of these issues will be presented in the future.
A very interesting and difficult problem is to analyze the behaviour of the solutions to the field equations. Since these solutions represent characteristic surfaces, they will have caustics and self-intersections. It can be shown that both $`r=^2Z`$ and $`\mathrm{\Lambda }=^3Z`$ diverge at conjugate points. Thus, the field equations as written above are not suited to analyze the behaviour of the solutions in the neighborhood of conjugate points. Possible alternative formulations are discussed in the last section.
### A Light cone cuts: an example
In this subsection we find the light cone cuts associated with particular simple models of space-times. These cuts of null infinity exhibit the typical behaviour that one has to deal with in NSF, namely, the presence of caustics and self-intersections.
As we show below, the simplest axisymmetric and static line element,
$$\mathrm{๐๐ }^2=dt^2+\rho ^{8M}\left(d\rho ^2+\rho ^2d\varphi ^2\right)$$
(33)
with a constant $`M`$ and polar coordinates $`(t,\rho ,\varphi )`$ does not give interesting cuts of $`๐ฅ^+`$. Nevertheless it is useful to follow the construction since it gives us a hint of how to modify the background space time to obtain non trivial cuts.
The metric (33) is interpreted as the space-time of a particle at the origin since its Einstein tensor $`G_{ab}`$ is given by $`M\delta (\rho )_at_bt`$.
After the coordinate change $`r=\frac{\rho ^\alpha }{\alpha }`$ with $`\alpha =14M`$, (no relationship with the canonical coordinate $`r`$ defined in previous sections) it can be put in the form
$$\mathrm{๐๐ }^2=dt^2+dr^2+\alpha ^2r^2d\varphi ^2,$$
(34)
which is a locally flat metric with a global deficit angle of $`2\pi \left(1\alpha \right)`$. Since the solution is not regular at the world-line of the particle $`\rho =0`$, this line must be removed from the otherwise flat space-time. Thus, the integration of the geodesic equation is trivial except for the following. For any apex $`x^a`$ there will be one null geodesic that will point to the origin and will not reach $`๐ฅ^+`$. Therefore, the light-cone cuts at null infinity will be open curves without caustics.
In order to obtain closed cuts of null infinity, we must then search for a matter source with a non-singular metric. The simplest choice is a ring of mass whose metric can be thought of a flat interior matched to a vacuum external solution of the form ( 34):
$$ds^2=\{\begin{array}{cc}dt^2+dr^2+r^2d\varphi ^2,\hfill & 0rR,\hfill \\ dt^2+dr^2+\left(\frac{R}{R+a}\right)^2(r+a)^2d\varphi ^2,\hfill & Rr,\hfill \end{array}$$
(35)
where the parameter $`R`$ is the radius of the ring and $`a`$ is related to the mass. A calculation of its energy momentum gives,
$$T_{ab}=\frac{a}{8\pi R(R+a)}\delta (rR)_at_bt.$$
(36)
Integration on a spacelike surface gives $`M={\displaystyle \frac{a}{4(R+a)}}`$, so that $`M0`$ implies that we must take $`a0`$. The Minkowski space corresponds to $`a=0`$. There is another range of negative $`a`$โs that also gives a positive mass but these paramenters are not continuously connected to $`a=0`$.
To solve the null geodesics equation and obtain the null cones we introduce the retarded time coordinate $`u=tr`$, and rewrite (35) as
$$ds^2=du^22dudr+x(r)^2d\varphi ^2,$$
(37)
with
$$x(r)=\{\begin{array}{cc}r,\hfill & 0rR,\hfill \\ R\frac{(r+a)}{(R+a)}.\hfill & Rr\hfill \end{array}$$
(38)
A first integral of the geodesic equations can be obtained by observing that $`k^ak_a=0`$ and that the scalar product between the null vector $`k^a`$ and the Killing fields $`\chi _u^a`$ and $`\xi _\varphi ^a`$ are constant along the null geodesic.
$`k_a\chi _u^a`$ $`=`$ $`\dot{u}\dot{r}=E,`$ (39)
$`k_a\xi _\varphi ^a`$ $`=`$ $`x(r)^2\dot{\varphi }=L,`$ (40)
$`k_ak^a`$ $`=`$ $`\dot{u}^22\dot{u}\dot{r}+x(r)^2\dot{\varphi }^2=0.`$ (41)
Then, inserting equations ( 39) and ( 40) in ( 41) gives
$$\dot{r}=\pm E\sqrt{1\left(\frac{b}{x(r)}\right)^2};$$
(42)
with $`bL/E`$, and from ( 39), ( 40) and ( 42), we get
$`{\displaystyle \frac{du}{dr}}`$ $`=`$ $`\pm {\displaystyle \frac{1\sqrt{1\left(\frac{b}{x(r)}\right)^2}}{\sqrt{1\left(\frac{b}{x(r)}\right)^2}}},`$ (43)
$`{\displaystyle \frac{d\varphi }{dr}}`$ $`=`$ $`\pm {\displaystyle \frac{b}{x(r)^2\sqrt{1\left(\frac{b}{x(r)}\right)^2}}}.`$ (44)
Integrating ( 43) and ( 44) between $`r=r_0`$ and $`r=\mathrm{}`$, gives the desired light cone cut. The more interesting case occurs when $`r_0>R`$ so we will restrict ourselves to this case. We now consider the circle of null directions from the point $`(u_0,r_0,\varphi _0=0)`$ and divide this circle in four quadrants. Since the solution has axial symmetry with respect to the line that joins $`r_0`$ with the origin we will only consider the upper quadrants which are distinguished by the initial value of $`\dot{r}(r_0)`$. The integration on the quadrant where $`\dot{r}(r_0)0`$ is straightforward and does not give any caustics. (Along those lines one takes the positive sign in ( 43) and ( 44).) The value of $`b`$ determines the starting null direction. This parameter ranges between $`0`$, when it points away from the origin and $`b_m=x(r_0)`$, when $`\dot{r}(r_0)=0`$. It is does appropriate to introduce a null angle $`0\theta {\displaystyle \frac{\pi }{2}}`$ defined by $`\mathrm{sin}\theta ={\displaystyle \frac{b}{b_m}}`$.
Care must be taken in the quadrant where $`\dot{r}(r_0)0`$. Here $`b`$ also ranges between $`b_m`$, and $`0`$, when it points to the origin. Then, the range of $`\theta `$ in this quadrant is given by $`{\displaystyle \frac{\pi }{2}}\theta \pi `$. However, there exists a critical value of b, called $`b_c`$ (and a corresponding $`\theta _c`$) such that for $`bb_c`$ the geodesics enter the ring. This value of $`b_c=R`$ is obtained from the condition that $`\dot{r}(R)=0`$.
We start the integration with $`\dot{r}<0`$ until the turning point $`\dot{r}(r_t)=0`$ is reached. At this point $`x(r_t)=b`$ from which the value of $`r_t`$ is obtained. Then, for $`r_tr\mathrm{}`$ we choose the positive sign in ( 43) and ( 44) to finish the integration.
The light cone cut of $`๐ฅ^+`$ from a point with coordinates $`(u_0,r_0,\varphi _0=0)`$ will be given by the following $`\theta `$-parametrized relations between $`u`$ and $`\varphi `$:
$`uu_0`$ $`=`$ $`(r_0+a)(1\mathrm{cos}\theta ),`$ (45)
$`\varphi `$ $`=`$ $`\theta \left(1+{\displaystyle \frac{a}{R}}\right),0\theta \theta _c`$ (46)
$`uu_0`$ $`=`$ $`(r_0+a)(1\mathrm{cos}\theta )2a\sqrt{1\left({\displaystyle \frac{r_0+a}{R+a}}\mathrm{sin}\theta \right)^2},`$ (47)
$`\varphi `$ $`=`$ $`\theta \left(1+{\displaystyle \frac{a}{R}}\right)2{\displaystyle \frac{a}{R}}\mathrm{arccos}({\displaystyle \frac{r_0+a}{R+a}}\mathrm{sin}\theta ),\theta _c\theta \pi ;`$ (48)
We now give a brief analysis of the singularity structure of the light cone cuts for this particular case. Without involving heavy use of singularity theorems we can see that caustics will occur when $`\varphi `$ is not an injective function of $`\theta `$. (If otherwise, we could invert the relation to obtain $`\theta =\theta (\varphi )`$ insert this back into $`u=U(\theta (\varphi ))`$ and trivially show that the cut is the graph of a function without self intersections.)
The graphs $`\varphi `$ vs $`\theta `$ for fixed values of $`R=3`$ and $`a=1`$ and two values of $`r_0=5,10`$ given below serve as illustrations of the most general situation.
In Fig. 1.a we observe that around $`\theta =\theta _c`$ the function fails to be injective. This case, however, is not analyzed in the theory of caustics since at $`\theta _c`$ the function is not differentiable. We thus adopt the criteria that whenever the function $`\varphi (\theta )`$ is continuous but not differentiable, we define a caustic point to be such that there is a change of sign of the left and right derivatives at this point. It is worth mentioning that the non differentiability around this caustic point follows from the fact that the Christoffel symbols, the coefficients of the geodesic equation, have a jump at $`r=R`$. By smoothing out the shell with a mass density function we can recover differentiability in the solution of the geodesic equation.
For the other value $`r_0=5`$ (see Fig. 1.b), we note the presence of a second caustic point $`\theta _0`$ where $`{\displaystyle \frac{d\varphi }{d\theta }}(\theta _0)=0`$.
In general, we will always have a caustic point $`\theta _c`$ for any value of $`R`$, $`a`$, and $`r_0`$ and a second caustic point whenever $`r_0`$ ($`>R`$) satisfies the condition
$`2a(r_0+a)<(R+a)^2.`$
We can also plot the corresponding light cone cuts.
As we can see both graphically and analytically, $`{\displaystyle \frac{du}{d\varphi }}`$ is well defined and finite at the caustic points. These caustics are called โcuspsโ. It is rather surprising that, to first order, the singularity structure around $`\theta _c`$ and $`\theta _0`$ is similar. The behaviours of $`{\displaystyle \frac{d^2u}{d\varphi ^2}}`$ and $`{\displaystyle \frac{d^3u}{d\varphi ^3}}`$ however, are completely different. Whereas at $`\theta _0`$ the second and third derivatives blow up, at $`\theta _c`$ they are distributional.
## V Conclusions
We close this work with several comments:
* The cut function $`Z`$ only gives a local description of a light cone cut. Using the available formalism to describe characteristic surfaces one can show that a light cone cut will have self intersections (thus $`Z`$ will not be unique) and singular points where $`^2Z`$ and $`^3Z`$ diverge. It is thus desirable to find alternative frameworks to describe the cut near singular points and/or globally on $`๐ฅ^+`$. One of these alternative approaches is given by the method of generating families of Legendre submanifolds which might lead to a global description of the light cone cuts.
* Likewise, the null hypersurface $`Z=const.`$ starts as a smooth surface near $`๐ฅ^+`$ but later develops conjugate points and caustics after the surface enters a region with $`T_{ab}0`$. A similar approach using generating families might give a global description of these null surfaces. In particular, there is a need to find a version of the metricity condition that will remain smooth around caustic points.
* We mention again that an interesting and open problem is to study the geometrical meaning of the metricity conditions in four dimensions. We conjecture that they are a torsion free condition on the Cartan connection.
* The most general singularity of a light cone cut in 3-D can be found on the explicit example presented above. In addition, we find another singular behaviour that is not consider in the literature since the Legendre submanifold is assumed to be at least $`C^2`$. It appears that by relaxing differentiability conditions on the Legendre submanifold one should be able to prove generic results for this new type of singularity.
### Acknowledgments
This research has been partially supported by AIT, CONICET, CONICOR, and UNC. We thank Paul Tod for enlightening discussions.
## A Conformal Einstein metrics and the Bach equation.
In this Appendix we derive the field equations for metrics conformally related to Einstein space times. This derivation will be restricted to three and four dimensions. Let $`\widehat{M}`$ be a n-dimensional manifold (n=3 or 4) and $`\widehat{g}_{ab}`$ a metric with arbitrary signature. Note that we are not following the same conventions as in the main body of the paper where the physical metric is โunhattedโ. This is done to simplify the notation for the conformal metrics and their related tensors.
Assume that the metric $`\widehat{g}_{ab}`$ satisfies the Einsteinโs equation, i. e.
$$\widehat{R}_{ab}\frac{\widehat{R}}{2}\widehat{g}_{ab}=T_{ab}$$
where $`T_{ab}`$ is the stress-energy tensor, and $`\widehat{R}_{ab}`$, $`\widehat{R}`$, the Ricci tensor and the Ricci scalar of $`\widehat{g}_{ab}`$ respectively. Introducing
$$\widehat{S}_{ab}=\frac{1}{n2}\left(\widehat{R}_{ab}\frac{\widehat{R}}{2(n1)}\widehat{g}_{ab}\right),$$
we can rewrite the equations as
$$\widehat{S}_{ab}\widehat{S}\widehat{g}_{ab}=\frac{1}{(n2)}T_{ab}.$$
We shall express this equation in terms of a conformal metric $`g_{ab}`$ which is obtained from $`\widehat{g}_{ab}`$ by rescaling with the conformal factor $`\mathrm{\Omega }`$. (Note that this conformal factor is not the conformal factor for an asymptotically flat spacetime as the one given in section IV.)
$$g_{ab}=\mathrm{\Omega }^2\widehat{g}_{ab}\text{or equivalently}\widehat{g}^{ab}=\mathrm{\Omega }^2g^{ab}.$$
(A1)
We now write the transformation law of various tensor fields under this conformal rescaling. Let us first decompose the Riemann curvature tensor of $`g_{ab}`$ according to
$$R_{abcd}=C_{abcd}+2(g_{a[c}S_{d]b}g_{b[c}S_{d]a});$$
here we assume that $`n3`$. In terms of $`S_{ab}`$ the Bianchi identities can be put as
$$^aS_{ab}=_bS$$
(A2)
and
$$^dC_{abcd}=2(n3)_{[a}S_{b]c}.$$
(A3)
Considering the conformal transformation given by (A1), the tensor $`\widehat{S}_{ab}`$ transforms
$$\widehat{S}_{ab}\mathrm{\Omega }=S_{ab}\mathrm{\Omega }+_a_b\mathrm{\Omega }\frac{\mathrm{\Omega }^1}{2}g_{ab}_c\mathrm{\Omega }^c\mathrm{\Omega }.$$
(A4)
Then the conformal Einsteins equation becomes
$$\widehat{S}_{ab}\mathrm{\Omega }=S_{ab}\mathrm{\Omega }+_a_b\mathrm{\Omega }\frac{\mathrm{\Omega }^1}{2}g_{ab}_c\mathrm{\Omega }^c\mathrm{\Omega },$$
(A5)
where
$$\widehat{S}_{ab}=\frac{1}{n2}\left(T_{ab}\frac{T}{n1}g_{ab}\right)$$
with $`T=T_{ab}g^{ab}`$. Since the stress-energy tensor satisfies $`\widehat{}^aT_{ab}=0`$, for later use, we calculate
$$\mathrm{\Omega }^aT_{ab}=(n2)^a\mathrm{\Omega }T_{ab}_b\mathrm{\Omega }T.$$
(A6)
We now ask what conditions should be imposed on $`g_{ab}`$ so that a solution of (A5) exists. If we consider eq. (A5) as a second order differential equation for $`\mathrm{\Omega }`$ it is then clear that, for a non trivial solution to exist, integrability conditions must be imposed on the metric $`g_{ab}`$. Applying $`_c`$ to (A5), and antisymmetrizing, we get
$$_{[c}\widehat{S}_{a]b}\mathrm{\Omega }+_{[c}\mathrm{\Omega }\widehat{S}_{a]b}=_{[c}S_{a]b}\mathrm{\Omega }+_{[c}\mathrm{\Omega }S_{a]b}+_{[c}_{a]}_b\mathrm{\Omega }\frac{g_{b[a}_{c]}}{2}(\mathrm{\Omega }^1_d\mathrm{\Omega }^d\mathrm{\Omega }).$$
(A7)
Contracting equation (A5) with $`^a\mathrm{\Omega }`$ we find
$$^a\mathrm{\Omega }\widehat{S}_{ab}^a\mathrm{\Omega }S_{ab}=\frac{_b}{2}(\mathrm{\Omega }^1_d\mathrm{\Omega }^d\mathrm{\Omega }),$$
and replacing the last term of (A7) by this result, yields
$$_{[c}\widehat{S}_{a]b}\mathrm{\Omega }+_{[c}\mathrm{\Omega }\widehat{S}_{a]b}g_{b[c}^d\mathrm{\Omega }\widehat{S}_{a]d}=_{[c}S_{a]b}\mathrm{\Omega }+_{[c}\mathrm{\Omega }S_{a]b}+_{[c}_{a]}_b\mathrm{\Omega }g_{b[c}^d\mathrm{\Omega }S_{a]d}.$$
(A8)
Finally, using the definition of the Riemann tensor, we calculate
$$_{[c}_{a]}_b\mathrm{\Omega }=\frac{1}{2}C_{cabd}^d\mathrm{\Omega }+g_{b[c}S_{a]d}^d\mathrm{\Omega }_{[c}\mathrm{\Omega }S_{a]b},$$
and equation (A8) becomes
$$2(_{[a}\widehat{S}_{b]c}\mathrm{\Omega }+_{[a}\mathrm{\Omega }\widehat{S}_{b]c}g_{c[a}^d\mathrm{\Omega }\widehat{S}_{b]d})=2\mathrm{\Omega }_{[a}S_{b]c}+C_{abcd}^d\mathrm{\Omega }$$
(A9)
Now we consider the case where the stress-energy tensor is pure trace, i.e.
$$T_{ab}=\mathrm{\Lambda }\widehat{g}_{ab},$$
where $`\mathrm{\Lambda }`$ is the cosmological constant. This gives in terms of $`g_{ab}`$
$$\widehat{S}_{ab}=\frac{\mathrm{\Lambda }\mathrm{\Omega }^2}{(n2)(n1)}g_{ab}.$$
Inserting this in the l.h.s. of (A9), we have
$$_{[a}\widehat{S}_{b]c}\mathrm{\Omega }+_{[a}\mathrm{\Omega }\widehat{S}_{b]c}g_{c[a}^d\mathrm{\Omega }\widehat{S}_{b]d}=\frac{\mathrm{\Lambda }\mathrm{\Omega }^2}{(n2)(n1)}\left(2_{[a}\mathrm{\Omega }g_{b]c}+_{[a}\mathrm{\Omega }g_{b]c}g_{c[a}_{b]}\mathrm{\Omega }\right)=0.$$
In dimension $`n=3`$, using the fact that $`C_{abcd}=0`$, equation (A9) becomes
$$2_{[a}S_{b]c}=0.$$
(A10)
It is easy to prove that the tensor $`_{[a}S_{b]c}`$ is conformally invariant since
$`\mathrm{\Omega }_{[a}S_{b]c}`$ $`=`$ $`_{[a}\widehat{S}_{b]c}\mathrm{\Omega }+_{[a}\mathrm{\Omega }\widehat{S}_{b]c}g_{c[a}^d\mathrm{\Omega }\widehat{S}_{b]d}`$ (A11)
$`=`$ $`\mathrm{\Omega }(\widehat{}_{[a}\widehat{S}_{b]c}C^l{}_{c[a}{}^{}S_{b]l}^{})+_{[a}\mathrm{\Omega }\widehat{S}_{b]c}g_{c[a}^d\mathrm{\Omega }\widehat{S}_{b]d}`$ (A12)
$`=`$ $`\mathrm{\Omega }\widehat{}_{[a}\widehat{S}_{b]c},`$ (A13)
where
$$C^l{}_{ac}{}^{}=\mathrm{\Omega }^1(2\delta _{(a}^l_{c)}\mathrm{\Omega }g_{ac}^l\mathrm{\Omega }).$$
In order to obtain the integrability condition for $`n>3`$, we use Bianchi identity (A3) and equation (A9) becomes
$$C_{abcd}^d\omega \frac{1}{n3}^dC_{abcd}=0,$$
(A14)
where for simplicity in the calculations we have defined $`w=\mathrm{log}\mathrm{\Omega }`$. Applying $`^a`$ to (A14) we find
$$^a^d\omega C_{abcd}+^aC_{abcd}^d\omega \frac{1}{n3}^a^dC_{abcd}=0$$
and using (A5) and (A14), we get
$$\frac{1}{n2}R^{ad}C_{abcd}+\frac{1}{n3}^a^dC_{abcd}\frac{n4}{n3}^d\omega ^aC_{abcd}=0.$$
(A15)
Clearly if $`n=4`$ we have the corresponding Bach equation
$$B_{bc}:=\frac{R^{ad}}{2}C_{abcd}+^a^dC_{abcd}=0.$$
(A16)
## B The Bach equation with stress-energy tensor
In order to handle a more general stress-energy tensor than the pure trace, we may suppose that the stress-energy tensor is given and that it does not depend on $`\mathrm{\Omega }`$.
We rewrite equation (A9) in terms of the stress-energy tensor and after using (A6) we get
$$2_{[a}S_{b]c}+C_{abcd}^d\omega =t_{abc},$$
(B1)
where
$$t_{abc}=\frac{2}{n2}\left(_{[a}T_{b]c}\frac{g_{c[a}^dT_{b]d}}{n2}\frac{_{[a}Tg_{b]c}}{n1}+_{[a}\omega T_{b]c}\frac{(n3)T}{(n2)(n1)}_{[a}\omega g_{b]c}\right).$$
We shall concentrate our attention to two cases, when the dimensions are $`n=3`$ and $`n=4`$. In the first case, equation (B1) becomes
$$B_{cab}=2\left(_{[a}T_{b]c}g_{c[a}^dT_{b]d}\frac{_{[a}Tg_{b]c}}{2}+_{[a}\omega T_{b]c}\right).$$
(B2)
Contracting this equation with $`T_{bc}`$ and using twice (A6) , we get
$$(T_{bc}T^{bc}T^2)_a\omega =T^{bc}B_{cab}2T^{bc}(_{[a}T_{b]c}g_{c[a}^dT_{b]d}\frac{_{[a}Tg_{b]c}}{2})+^bT_{bc}T_a{}_{}{}^{c}+T^cT_{ca}.$$
(B3)
Demanding that $`T_{bc}T^{bc}T^20`$ and inserting (B3) in (B2), we obtain
$`(T_{de}T^{de}T^2)B_{cab}`$ $`=`$ $`2\left(T_{de}T^{de}T^2\right)\left(_{[a}T_{b]c}g_{c[a}^dT_{b]d}{\displaystyle \frac{_{[a}Tg_{b]c}}{2}}\right)+T_{bc}T^{de}B_{ead}`$ (B7)
$`T_{ac}T^{de}B_{ebd}2T^{de}T_{bc}\left(_{[a}T_{d]e}g_{e[a}^lT_{d]l}{\displaystyle \frac{_{[a}Tg_{d]e}}{2}}\right)+`$
$`+2T^{de}T_{ac}\left(_{[b}T_{d]e}g_{e[b}^lT_{d]l}{\displaystyle \frac{_{[b}Tg_{d]e}}{2}}\right)+`$
$`+T_{bc}\left(^dT_{de}T_a{}_{}{}^{e}+T^eT_{ea}\right)T_{ac}\left(^dT_{de}T_b{}_{}{}^{e}T^eT_{eb}\right).`$
When $`n=4`$ equation (B1) becomes
$$C_{abcd}^d\omega ^dC_{abcd}=t_{abc},$$
(B8)
where
$$t_{abc}=_{[a}T_{b]c}\frac{g_{c[a}^dT_{b]d}}{2}\frac{_{[a}Tg_{b]c}}{3}+_{[a}\omega T_{b]c}\frac{T}{6}_{[a}\omega g_{b]c}.$$
(B9)
As before, applying $`^a`$ and using Einsteins equations we can write
$$2B_{ab}=T^{ad}C_{abcd}2\left(t_{acb}^a\omega +^at_{abc}\right).$$
(B10)
A simple calculation gives
$`2\left(t_{acb}^a\omega +^at_{abc}\right)`$ $`=`$ $`2^a(_{[a}T_{b]c}{\displaystyle \frac{g_{c[a}^dT_{b]d}}{2}}{\displaystyle \frac{_{[a}Tg_{b]c}}{3}}){\displaystyle \frac{1}{2}}T_{ac}(T^a{}_{b}{}^{}{\displaystyle \frac{T}{6}}g^a{}_{b}{}^{})+`$
$`+{\displaystyle \frac{1}{2}}T_{ac}R^a{}_{b}{}^{}T_c\omega _b\omega +{\displaystyle \frac{3}{2}}^a\omega _a\omega (T_{bc}{\displaystyle \frac{T}{6}}g_{bc})+`$
$`+2^a\omega _aT_{bc}{\displaystyle \frac{1}{2}}^a\omega _aTg_{bc}+{\displaystyle \frac{1}{2}}_b\omega _cT^a\omega _cT_{ab}+`$
$`+{\displaystyle \frac{1}{2}}g_{bc}^a\omega ^dT_{ad}{\displaystyle \frac{1}{2}}_c\omega ^aT_{ab}^aT_{ac}_b\omega .`$
Contracting (B8) with $`C^{abc}_l`$ and using twice (A6), we get
$`(3C_{abcd}C^{abcd}6T_{bc}T^{bc}+2T^2)_a\omega `$ $`=`$ $`12_d(C^{den}{}_{a}{}^{}T_{en}^{})+3_a(T^{de}T_{de})6_d(T_{ae}T^{de})`$ (B13)
$`3^dT_{de}T^e{}_{a}{}^{}+5T^dT_{da}_aT^2+2T^e{}_{a}{}^{}_{e}^{}T+`$
$`+12^lC_{denl}C^{den}{}_{a}{}^{}.`$
Demanding that $`3C_{abcd}C^{abcd}6T_{bc}T^{bc}+2T^20`$ and inserting (B13) in (B10), we obtain a similar result than for dimension 3.
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# Fluctuation-Facilitated Charge Migration along DNA
\[
## Abstract
We propose a model Hamiltonian for charge transfer along the DNA double helix with temperature driven fluctuations in the base pair positions acting as the rate limiting factor for charge transfer between neighboring base pairs. We compare the predictions of the model with the recent work of J.K. Barton and A.H. Zewail (Proc. Natl. Acad. Sci. USA, 96, 6014 (1999)) on the unusual two-stage charge transfer of DNA.
\]
Charge transport and electrical conduction are known to occur in a wide range of organic linear chain crystals of stacked planar molecules . Transfer rates from molecule to molecule are determined by the single-particle integral $`\tau `$, with typical rates of the order of 10<sup>15</sup> sec<sup>-1</sup>. Strong interaction between the electronic degrees of freedom and molecular vibrations may reduce this to 10<sup>12</sup> sec<sup>-1</sup>, a typical lattice or intramolecular vibration frequency. By comparison, biochemical charge transfer processes, such as those encountered in the metabolic redox (oxydation-reduction) chains , usually are much slower (down to 1 sec). Key steps often involve some form of large scale motion of the molecule.
DNA can be considered as a one-dimensional, a-periodic, linear chain of stacked base-pairs. More than 30 years ago it was suggested that duplex DNA might support electron transport in a manner similar to that of linear chain compounds, namely by tunneling along overlapping $`\pi `$ orbitals located on the base-pairs. Barton et al. first presented evidence that photo-induced, radical cations can travel along DNA molecules in aqueous solution over quite considerable distances (more than 40 ร
, or about ten base-pairs). If so, DNA might present us with flexible, molecular-size wires able to transport charge in aqueous environments. Possible applications range from micro-electronics to long range detection of DNA damage. Subsequent studies by a number of groups reported a wide range of values for the effective inverse spatial carrier decay length $`\beta `$, ranging from as low as 0.02 ร
<sup>-1</sup> to as high as around one ร
<sup>-1</sup>, the different values most probably reflecting differences in charge transfer for different base sequences .
Recently, Barton and Zewail (BZ) used femtosecond spectroscopy to measure the rates of the DNA charge transfer process. An unusual two-step decay process was observed with characteristic time scales of 5 and 75 ps respectively. Ab-initio molecular-orbital calculations find that DNA has a large single particle band-gap, and a transfer integral $`\tau `$ of order $`0.1`$ eV. This would lead to a charge transfer rate $`\tau /h`$ for coherent tunneling that is comparable to that of the linear chain compounds, but that is much too high compared to the rates measured by BZ.
Apart from coherent tunneling , a number of alternative transport mechanisms have been proposed, in particular incoherent, phonon-assisted electron hopping between bases, with the electron wave fully localized on each subsequent base-pair . This would reduce the transfer rate to a typical intra-molecular vibrational frequency (ps<sup>-1</sup>), but this is still much too large to explain the slow second stage step of the decay. It was also suggested that a charged radical could induce a polaronic distortion of the DNA internal structure that might control charge transfer. The explanation proposed by BZ for the long relaxation times is that large amplitude thermal fluctuations of the intercalated photoreceptor sets the rate limiting step for the charge transfer.
The aim of this letter is to construct a model Hamiltonian to treat charge transfer along a chain under conditions of large structural fluctuations, and suggest that thermally induced structural disorder interferes with the $`\pi `$ orbital overlap mediated charge transfer, leading to long relaxation times. To construct this Hamiltonian, we first must discuss the origin of the structural fluctuations. Figure 1 shows an example of a typical DNA configuration obtained by a molecular dynamics (MD) simulation . The relative orientation of neighboring bases along DNA is characterized by a set of collective variables such as the relative roll and twist angles (R and T), and the relative slide displacement (see Fig.1). Long time MD simulations of DNA lead to typical RMS fluctuation angles for R and T of order 5 and 9 degrees in the ps to ns time-window, while the mean base-pair spacing also shows large amplitude fluctuations . Structural fluctuations in the ps to ns time-window have been observed experimentally as dynamic Stokes shifts in the fluorescence spectrum of the DNA. The local fluctuations are extraordinarily strong compared with those due to thermally excited phonon modes in crystalline linear chain materials. The unusual โsoftnessโ of the R,S, and T variables is also reflected by the fact that their mean values vary greatly depending on base-pair sequence .
In our simplified model, we include only two collective modes. The first mode is an angular variable $`\theta (t)`$, which is that relative rotation angle of the two bases which couples most efficiently to the $`\pi `$-orbital tunnel matrix element. Next, the displacement variable $`y(t)`$ represents that collective mode which couples most efficiently to the
on-site energy of the radical. This second form of coupling provides the necessary mechanism for energy transfer between the charge and the the thermal reservoir required for hopping transport along a random sequence of base-pairs with different on-site energies. Both $`\theta (t)`$ and $`y(t)`$ are treated as classical harmonic variables that are coupled to a heat bath of oscillators. As a result, $`\theta (t)`$ and $`y(t)`$ obey -in the absence of the radical- the following Langevin equations:
$`M\ddot{y}+\gamma _y\dot{y}+M\mathrm{\Omega }_y^2y=\eta _y(t),`$ (1)
$`I\ddot{\theta }+\gamma _\theta \dot{\theta }+I\mathrm{\Omega }_\theta ^2\theta =\eta _\theta (t).`$ (2)
In Eq. 2, $`I`$ is the reduced moment of inertia for the relative rotation of the two adjacent bases, $`\mathrm{\Omega }_\theta `$ is the oscillator frequency of the rotation mode, while $`M`$ and $`\mathrm{\Omega }_y`$ are the reduced mass and natural frequency of the displacement mode. The values of $`M`$,$`I`$, $`\mathrm{\Omega }_\theta `$,$`\mathrm{\Omega }_y`$, and the damping coefficients are obtained by comparing the Fourier power spectra of $`y(t)`$ and $`\theta (t)`$ obtained from Eq. 2, to power spectra of MD simulations of DNA . From a typical long time (10 ns) MD series, we find oscillation periods of $`2\pi /\mathrm{\Omega }`$ of order 1-10 ps, a (large) mass $`M`$ of order 1-10 kDalton, (1 Dalton = 1 a.m.u.) a moment of inertia $`I`$ of order $`10^2k_B\mathrm{\Omega }_\theta ^2`$, and relaxation times comparable to the oscillation period (i.e. the slow modes are close to critical damping). The amplitudes of the white-noise variables $`\eta (t)`$ follow from the fluctuation-dissipation theorem for classical variables. We will neglect mode coupling between different pairs of adjacent base-pairs. The two modes are then coupled to a one-dimensional, tight-binding Hamiltonian for single particle charge transport:
$`H={\displaystyle \underset{i}{}}\{\tau (\theta _{i,i+1})(c_{i+1}^{}c_i+c_i^{}c_{i+1})+ฯต_ic_i^{}c_i`$ (3)
$`+{\displaystyle \frac{1}{2}}I_i(\dot{\theta }_{i,i+1}^2+\mathrm{\Omega }_{\theta ,i}^2\theta _{i,i+1}^2)`$ (4)
$`+{\displaystyle \frac{1}{2}}M_i(\dot{y}_i^2+\mathrm{\Omega }_{y,i}^2(y_i+y_{0,i}c_i^{}c_i)^2)\}+H_{bath}(\{\theta ,y\}).`$ (5)
In Eq. 5, $`ฯต_i`$ is the on-site electronic energy. The distance $`y_{0,i}`$ is the change in the equilibrium value of the $`y`$ variable of the i-th base when the particle localizes on that site, while $`M\mathrm{\Omega }^2y_0^2`$ is the typical deformation energy.
Certain limiting cases of this general Hamiltonian are familiar from studies of one-dimensional charge transport. For uniform $`ฯต_i`$ and for fixed $`\theta `$, $`H`$ is the Hamiltonian of a tight-binding polaron . For fixed $`\theta `$ and and $`y`$ and random $`ฯต_i`$, $`H`$ is the Anderson Hamiltonian for localization in one-dimension. For the case of DNA, we assume that site-to-site differences in the value of $`ฯต_i`$ are of order $`0.1`$eV based on the sequence dependent differences in the ionization potential . Next, the transfer integral $`\tau (\theta )`$ will be assumed to be small compared to the thermal energy $`k_BT`$ for $`\theta `$ near a special value, denoted by $`\theta ^{}`$, the โrapid decay stateโ. Finally, the characteristic interaction energy $`M\mathrm{\Omega }^2y_0^2`$ between the charged radical and the on-site structural variable $`y`$ is assumed to be large compared to $`k_BT`$ (e.g. due to electrostatic effects) and of order $`ฯต`$.
In this unusual, high temperature, strong coupling regime, the transfer integral $`\tau (\theta )`$ is the lowest energy scale. Under these conditions, particle motion described by $`H`$ is indeed dominated by incoherent hopping from site to site. We first restrict ourselves to the case of a particle which resides at site $`A`$ at time $`t=0`$ and then hops to the neighboring site $`B`$ with a different on-site energy. For fixed $`\theta `$, the transition rate $`\mathrm{\Gamma }(\theta )`$ for incoherent charge transfer between $`A`$ and $`B`$ can be computed by applying the method of Garg, Onuchic and Ambegaokar (GOA) . In Fig.2 we show the two potential energy surfaces $`V^+(y)=A|H|A`$ and $`V^{}(y)=B|H|B`$ for the particle respectively on the $`A`$ and $`B`$ sites as a function of $`y`$. Efficient transfer between the two potential energy surfaces takes place nearly exclusively at the crossing points $`y^{}`$ where $`V^+(y^{})=V^{}(y^{})`$, shown in Fig. 2. Note that an energy barrier $`E_r`$ must be overcome to reach this crossing point. In the high temperature, strong coupling limit, the on-site probability decays exponentially with a rate:
$$\mathrm{\Gamma }(\theta )\frac{\tau (\theta )^2}{\mathrm{}}\left[\frac{\pi }{E_rk_BT}\right]^{1/2}exp(E_f/k_BT),$$
(6)
where the energy scale $`E_r`$ defined in Fig.2, depends on
the difference between the on-site energies. The validity condition for Eq. 6 is that $`\tau (\theta )`$ must be small compared to $`(E_rk_BT)^{1/2}`$. By itself, Eq. 6 does not account for a two stage decay, as observed by BZ. In a typical ensemble of radical sites (with the same on-site energies) there must be significant heterogeneity, with less likely states characterized by low charge transfer rates. This heterogeneity is incorporated by demanding that $`\tau (\theta )`$ is appreciable only for $`\theta `$ near the special value $`\theta ^{}`$, while tunneling plays no role for different $`\theta `$ values . If $`\theta `$ undergoes large amplitude thermal fluctuations, then there should be considerable heterogeneity for the transfer rates.
Assume then that, at time $`t=0`$, an ensemble of particles is prepared on the $`A`$ site, with the $`\theta (0)`$ variable obeying the Boltzmann distribution. We define the probability density $`P(\theta ,t)d\theta `$ to be the fraction of radicals at time $`t`$ that are still on site $`A`$, and whose $`\theta `$ values is in the range between $`\theta `$ and $`\theta +d\theta `$. For the overdamped case of Eq. 2, $`P(\theta ,t)`$ obeys the following equation:
$`{\displaystyle \frac{P(\theta ,t)}{t}}={\displaystyle \frac{k_BT}{\gamma _\theta }}{\displaystyle \frac{^2P(\theta ,t)}{\theta ^2}}+{\displaystyle \frac{1}{\gamma _\theta }}{\displaystyle \frac{}{\theta }}[I\mathrm{\Omega }_\theta ^2\theta P(\theta ,t)]`$ (7)
$`\mathrm{\Gamma }(\theta )P(\theta ,t).`$ (8)
We can now use Eq. 8 to discuss the decay rate. For $`\mathrm{\Gamma }(\theta ^{})`$ large compared to the thermal equilibration rate $`\tau _\theta ^1`$ two forms of decay are encountered .
(i) Early stage decay: At times $`t=0`$, a certain fraction of the oscillators has an energy exceeding $`E=(1/2)I(\mathrm{\Omega }_\theta \theta ^{})^2`$. These oscillators will pass through the $`\theta =\theta ^{}`$ point within a time of order max$`\{\tau _\theta ,\mathrm{\Omega }_\theta ^1\}`$. When this happens, there is a finite probability for charge transfer to take place. On a time scale of order max$`\{\tau _\theta ,\mathrm{\Omega }_\theta ^1\}`$, these high energy oscillators are removed from the probability distribuition.
(ii) Late stage decay: After the high energy oscillators have been removed from the distribution, further decay requires energy diffusion along the oscillator scale from lower energies towards $`E(\theta ^{})=(1/2)I(\mathrm{\Omega }_\theta \theta ^{})^2`$. Once the energy of an oscillator reaches this value, efficient charge transfer takes place. After a standard, but lenghty, analysis of Eq. 8 we find that the late stage decay rate is:
$$k_{LATE}=\tau _\theta \mathrm{\Omega }_\theta ^2\left[\frac{\theta ^{}e^{\frac{1}{2}\beta I\mathrm{\Omega }_\theta ^2\theta _{}^{}{}_{}{}^{2}}}{_{\mathrm{}}^{\mathrm{}}e^{\frac{1}{2}\beta I\mathrm{\Omega }_\theta ^2\theta _{}^{}{}_{}{}^{2}}๐\theta }\right].$$
(9)
The factor in front of Eq. 9 is of the same order as the early stage decay rate while the term in brackets is of the order of the thermal probability that $`\theta `$ exceeds $`\theta ^{}`$. The second stage decay rate strongly increases with increasing temperature, while the early stage decay rate is not expected to be strongly temperature dependent although the fraction of sites that exhibits early stage decay should be strongly temperature dependent (of the same order of the term in brackets of Eq. 9).
The time dependence of $`P(\theta ,t)`$ would be consistent with the observations of BZ if this thermal probability of the $`\theta ^{}`$ state is of order $`10^2`$. In that case, one percent of the sites would show rapid decay with time scales of order ps, while the remainder would show decay slowed down by a factor of $`10^2`$. We only treated here the nearest neighbor hopping process. Charge transport over longer distances described by our Hamiltonian reduces - under the assumed conditions - to a classical one-dimensional diffusion in a random medium with site specific transfer rates. The transport properties of such systems have been extensively discussed elsewhere .
In conclusion, we propose that charge transport along DNA proceeds by classical diffusion with high amplitude thermal fluctuations providing the rate limiting step for the site-to-site charge transfer. If correct, charge transport along DNA would have unique characteristics as compared to linear chain compounds. Since the radical severly deforms the local structure it might be considered as a polaron in the strong coupling limit but, unlike polaronic transport, hopping is controlled by thermal fluctuations. Indeed, Eq. 9 predicts that the charge transfer rate should strongly increase with temperature which is consitent with the observations of BZ. A better description of the mechanism proposed in this paper for charge transfer along DNA would be to consider it as a repeated sequence of reversible oxydation-reduction reactions. The site-to-site charge transfer would be viewed as a chemical reaction dominated by a transition state where the collective variables $`y`$ and $`\theta `$ assume a special value ($`y^{}`$ and $`\theta ^{}`$ respectively). We are not aware of any of the linear chain compounds exhibiting this curious form of charge transfer. On the other hand, a recent single-molecule optical study of a particular reversible oxydation-reduction reaction did report two-stage non-exponential behavior but with decay rates much lower than those measured by BZ (in the range of 1 sec<sup>-1</sup>).
The higher rates of charge transfer in DNA would be due to the fact that the molecular motion of the bases still is significantly restrained by the backbone. If the present analysis is appropriate, then charge transport in DNA occupies a unique position intermediate between charge transport in solid-state materials and charge transport in biochemical charge transfer reactions.
Finally, the proposed Hamiltonian obviously incorporates a number of rather serious simplifications. We include the collective modes only in a schematic way. We require a large structural on-site distortion of a site by the particle (of order $`0.1`$ eV), but this is likely to require an anharmonic description of the collective modes. We did not include coupling of modes of adjacent pairs, the double-stranded nature of DNA with the possibility of inter-chain charge transfer or effects related to the tertiary structure, i.e. the coiling of the duplex.
We would like to thank J.Barton and F.Pincus for stimulating and useful discussions. We would like to thank J.Onuchic for providing a copy of his thesis and for discussions. RB would like to thank the Rothschild Foundation and NSF Grant DMR-9407741 for financial support.
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# 1 Introduction
## 1 Introduction
The aim of this paper is to demonstrate some peculiarities of generic free massless fields in anti-de Sitter (AdS) space of an arbitrary space-time dimension $`d`$. The main conclusion will be that, in general, an irreducible $`AdS_d`$ massless field does not classify according to irreducible representations of the flat space massless little algebra $`o(d2)`$, but reduces to a certain set of irreducible flat space massless fields. The pattern of necessary flat-space massless fields will be given. Another (rather unexpected) manifestation of this fact is that not every massless field in flat space admits a deformation to $`AdS_d`$ with the same number of degrees of freedom, since it is impossible to keep all the flat space gauge symmetries unbroken in the $`AdS_d`$ space. This phenomenon does not take place, though, for all those types of massless fields that appear in the usual low-energy massless sectors of the superstring models and supergravities, because it holds only for the representations of the space-time symmetries described by non-rectangular Young diagrams. For the same reason it cannot be observed in $`AdS_4`$ higher spin gauge theories (all massless fields in $`AdS_4`$ are described by one-row Young diagrams). The effect discussed in this paper takes place for $`d6`$.
In superstring theory all types of representations appear at the higher massive levels. The study of higher spin gauge theory has two main motivations (see e.g. ): Firstly to overcome the well-known barrier of $`N8`$ in $`d=4`$ supergravity models and, secondly, to investigate if there is a most symmetric phase of superstring theory that leads to the usual string theory as a result of a certain spontaneous breakdown of higher spin gauge symmetries. These two motivations lead in fact into the same direction because, as shown for the $`d=4`$ case higher spin gauge theories require infinite collections of higher spin gauge fields with infinitely increasing spins. Another important feature discovered in is that gauge invariant higher spin interactions require the cosmological constant $`\lambda ^2`$ of the background AdS space to be non-zero to compensate the extra length dimensions carried by the higher derivative interactions required by the higher spin gauge symmetries. (In this perspective $`\lambda `$ plays the rรดle analogous to $`\alpha ^{}`$ in superstring theory). The fact that higher spin theories require an AdS background was regarded as rather surprising until it was realized that it plays a distinguished rรดle in the superstring theory as well .
To investigate a possible relationship between the superstring theory and higher spin theories one has to build the higher spin gauge theory in higher dimensions, $`d>4`$ (e.g. $`d=10,11,\mathrm{}`$). A conjecture on the possible form of the higher spin symmetries and equations of motion for higher spin spin gauge fields was made in as a certain generalization of the $`d=4`$ results which were proved to describe interactions of all $`d=4`$ massless fields.
A generalization like this to higher dimensions is not straightforward because of the use of certain auxiliary twistor type variables. As a starting point, it is therefore important to analyze more carefully the notion of a general massless field in $`AdS_d`$. This is the main goal of this paper.
Another motivation comes from the flat space analysis of certain massless (nonsupersymmetric) triplets in $`d=11`$ in ,, the dimension of M-theory, where it was shown that there exists an infinite collection of triplets of higher spin fields having equal numbers of bosonic and fermionic degrees of freedom. These triplets show some remarkable properties. For the first four Dynkin indices the bosonic and the fermionic numbers match up. This phenomenon is rather special for 11 dimensions and follows from the fact that the little group $`SO(9)`$ is an equal-rank subgroup of $`F_4`$, bringing in exceptional groups into the picture. If these triplets have anything to do with higher spin gauge field theories and/or M-theory it is an interesting question whether it is possible to extend the analysis of to the AdS case. That question in fact triggered this investigation.
## 2 Massless Unitary Representations in $`AdS_d`$
$`AdS_d`$ is a $`d`$-dimensional space-time with signature $`(d1,1)`$ and the group of motions $`SO(d1,2)`$. It is most useful to identify $`AdS_d`$ with the universal covering space of the appropriate hyperboloid embedded into Minkowski space-time with signature
$`(d1,2)`$. Physically meaningful relativistic fields in $`AdS_d`$ are classified according to the lowest weight unitary representations of $`o(d1,2)`$. Unitarity implies compatibility with quantum mechanics, while lowest weight of a unitary representation guarantees that the energy is bounded from below.
### 2.1 General Facts
The commutation relations of $`o(d1,2)`$ are
$$[M_{\widehat{m}\widehat{n}},M_{\widehat{k}\widehat{l}}]=i(\eta _{\widehat{n}\widehat{k}}M_{\widehat{m}\widehat{l}}\eta _{\widehat{m}\widehat{k}}M_{\widehat{n}\widehat{l}}\eta _{\widehat{n}\widehat{l}}M_{\widehat{m}\widehat{k}}+\eta _{\widehat{m}\widehat{l}}M_{\widehat{n}\widehat{k}}),$$
(1)
where $`\eta _{\widehat{n}\widehat{m}}=(,+,\mathrm{}+,)`$ is the flat metric in the $`(d1,2)`$ space, $`\widehat{m},\widehat{n},\widehat{k},\widehat{l}=0รทd`$. The generators $`M_{\widehat{m}\widehat{n}}`$ are Hermitian. Let us choose the following basis in the algebra:
$$t_\pm ^a=\frac{1}{2}(M_0{}_{}{}^{a}\pm iM_d{}_{}{}^{a}),$$
(2)
$$E=M_{0d},L^{ab}=iM^{ab},$$
(3)
where $`a,b=1รทd1`$. The commutation relations (1) take the form
$$[E,t_\pm ^a]=\pm t_\pm ^a,$$
(4)
$$[t_{}^a,t_+^b]=\frac{1}{2}(E\delta ^{ab}L^{ab}),$$
(5)
$$[L_{ab},t_{\pm c}]=\delta _{bc}t_{\pm a}\delta _{ac}t_{\pm b},$$
(6)
$$[L_{ab},L_{ce}]=\delta _{bc}L_{ae}\delta _{ac}L_{be}\delta _{be}L_{ac}+\delta _{ae}L_{bc}$$
(7)
with all other commutators vanishing. The hermiticity conditions are
$$E^{}=E,(t_\pm ^a)^{}=t_{}^a,(L^{ab})^{}=L^{ab}.$$
(8)
The generators $`E`$ and $`L^{ab}`$ can be identified with the energy and angular momenta, respectively, and span the Lie algebra $`o(2)o(d1)`$ of the maximal compact subgroup of the AdS group $`SO(2,d1)`$. The non-compact generators $`t_\pm ^a`$ are combinations of $`AdS`$ translations and Lorentz boosts. The commutation relations are explicitly $`Z`$-graded with $`t_\pm ^a`$ having grade $`\pm 1`$. $`E`$ is the grading operator. The lowest weight unitary representations are now constructed in the standard fashion (for a review of the $`AdS_4`$ case see e.g. and ) starting with the vacuum space $`|E_0,๐ฌ`$ that is annihilated by $`t_{}^a`$
$$t_{}^a|E_0,๐ฌ=0$$
(9)
and forms a unitary representation of the compact subalgebra $`o(2)o(d1)`$ that means in particular that
$$E|E_0,s=E_0|E_0,๐ฌ.$$
(10)
Here $`๐ฌ`$ denotes the type of representation of $`o(d1)`$ carried by the vacuum space: $`๐ฌ=(s_1,\mathrm{}s_\nu )`$ with $`\nu =[\frac{d1}{2}]`$. $`๐ฌ`$ is a generalized spin characterizing the representation. In terms of Young tableaux $`s_i`$ is the number of cells in the $`i`$-th row of the Young tableaux. Since we are talking about representations of orthogonal algebras, the corresponding tensors are traceless. We will here not discuss the self-dual representations that can be singled out with the aid of the Levi-Civita symbol<sup>1</sup><sup>1</sup>1The self-dual and antiself-dual representations which appear for odd $`d`$ are usually distinguished by a sign of the $`s_{\frac{d1}{2}}`$.. Note that $`๐ฌ`$ describes a finite-dimensional representation of $`o(d1)`$. Field-theoretically this corresponds to a finite-component field carrying a finite spin.
The full representation of the Lie algebra $`o(d1,2)`$, denoted in $`D(E_0,๐ฌ)`$, is spanned by the vectors of the form
$$t_+^{a_1}\mathrm{}t_+^{a_k}|E_0,๐ฌ$$
(11)
for all $`k`$. The states with fixed $`k`$ are called level-$`k`$ states. Note that states with pairwise different $`k`$ are orthogonal as a consequence of (4). As a result, the analysis of unitarity can be reformulated as the check of the positivity of the norms of the finite-dimensional subspaces at every level. For sufficiently large generic $`E_0`$ it is intuitively clear that the representation $`D(E_0,๐ฌ)`$ is irreducible and unitary, i.e. all norms are strictly positive. Such representations are identified with massive representations of $`AdS_d`$.
Let us emphasize that the elements of the module (11) can be identified with the modes of a one-particle state in the corresponding free quantum field theory. The elements of the $`AdS_d`$ algebra $`o(d1,2)`$ are then realized as bilinears in the quantum fields. Note that at the level of equations of motion the light-cone field-theoretical realization of generic $`AdS_d`$ massive representations for arbitrary $`E_0`$ and $`๐ฌ`$ has been developed recently in .
We see that massive states are classified by the parameter $`E_0`$ (which is the analog of mass) and a representation of $`o(d1)`$. This picture is in agreement with the standard description of massive relativistic fields in flat space-time in terms of the Wigner little group $`SO(d1)`$.
If $`E_0`$ gets sufficiently small, the norm cannot stay positive as is most obvious from the following consequence of (5)
$$[t_{}{}_{a}{}^{},t_+^a]=\frac{d1}{2}E,$$
(12)
which implies that for negative $`E_0`$ some level 1 states cannot have positive norm for a positive-definite vacuum subspace. There is therefore a boundary of the unitarity region $`E_0=E_0(๐ฌ)>0`$ such that some states acquire negative norm for $`E_0<E_0(๐ฌ)`$. Obviously, these states should have zero norm for $`E_0=E_0(๐ฌ)`$.
Starting from the inside of the unitarity region decreasing $`E_0`$ for some fixed $`๐ฌ`$ one approaches the boundary of the unitarity region, $`E_0=E_0(๐ฌ)`$. Some zero-norm vectors then appear for $`E_0=E_0(๐ฌ)`$. These necessarily should have vanishing scalar product to any other state. (Otherwise, one can build a negative norm state that is in contradiction with the assumption that we are at the boundary of the unitarity region). Therefore, the zero-norm states form an invariant subspace called a singular submodule. By factoring out this subspace one is left with a unitary representation which is โshorterโ than the generic massive representation.
The resulting โshortenedโ unitary representations correspond either to massless fields or to singletons and doubletons , identified with the conformal fields at the boundary of the $`AdS`$ space . The fact that a singular submodule can be factored out admits an interpretation as some sort of a gauge symmetry (true gauge symmetry for the case of massless fields, or independence of bulk degrees of freedom for singletons and doubletons). For this reason we choose this definition of masslessness for all fields in $`AdS_d`$ except for the scalar and spinor massless matter fields which are not associated with any gauge symmetry principle and singletons<sup>2</sup><sup>2</sup>2Singleton-type fields live at the boundary of $`AdS_d`$ and cannot be interpreted as bulk massless fields. Presumably, all singletons except for scalar and spinor correspond to maximally antisymmetrized representations of the $`AdS_d`$ algebra equivalent to their duals. In particular, this is true for the second rank antisymmetric tensor representation in the $`AdS_5`$ case that can be identified with the field strength of the Yang-Mills fields of the $`N=4`$ $`SYM`$ at the boundary of $`AdS_5`$..
The analysis of positive definiteness of the scalar product of level-1 states was done in for an arbitrary even $`d`$ and an arbitrary type of representation of $`๐ฌ`$ carried by the vacuum state. The final result is
$$E_0E_0(๐ฌ),$$
(13)
$$E_0(๐ฌ)=s_1+dt_12,$$
(14)
where $`t_1`$ is the number of rows of the maximal length $`s_1`$, i.e.
$$s_1=\mathrm{}=s_{t_11}=s_{t_1}>s_{t_1+1}s_{t_1+2}\mathrm{}s_\nu .$$
It can be shown that the same result is true for odd $`d`$ (provided that $`s_\nu `$ is replaced by $`|s_\nu |`$). Note that this bound for $`d=4`$ was originally found in . For the case $`d=5`$ see and references therein.
Massless representations are โshorterโ than massive ones classified according to the parameter $`E_0`$ (equivalent of mass) and a representation of $`o(d1)`$. In flat space, massless fields are classified according to the representation of the massless little group $`SO(d2)`$. The question we address here is whether or not the shortening in $`AdS_d`$ can be interpreted in terms of irreducible representations of $`SO(d2)`$. We will show that the answer is no for a generic representation. This will be demonstrated both at the algebraic level using the language of singular vectors and at the field-theoretical level focusing on the simplest nontrivial massless field with $`๐ฌ=(2,1,0,\mathrm{},0)`$. The most important field-theoretical conclusion is that a generic irreducible massless field in $`AdS_d`$ decomposes into a collection of massless fields in the flat limit. In that sense, a massless field in the $`AdS_d`$ is generically โless masslessโ than elementary massless field in flat space. An important consequence of this fact is that not every massless field in flat space can be deformed into AdS geometry.
### 2.2 Singular Vectors
It is useful to reformulate the problem in terms of singular vectors. Since the energy $`E`$ is bounded from below for the whole representation, the singular submodule spanned by zero-norm states is itself a lowest weight representation. Therefore it contains at least one nontrivial subspace $`|E_0^{},๐ฌ^{}`$ that has the properties analogous to (9) and (10),
$$t_{}^a|E_0^{},๐ฌ^{}=0$$
(15)
and
$$E|E_0^{},๐ฌ^{}=E_0^{}|E_0^{},๐ฌ^{},$$
(16)
i.e. it forms some irreducible representation $`๐ฌ^{}`$ of $`o(d1)`$. Obviously,
$$E_0^{}=E_0+k$$
(17)
if $`|E_0^{},๐ฌ^{}`$ belongs to the level-$`k`$ subspace.
Such spaces $`|E_0^{},๐ฌ^{}`$ we will call singular vacuum spaces while any of their elements will be called a singular vector<sup>3</sup><sup>3</sup>3This terminology is very closely although not exactly coinciding with that used for the Verma module construction when irreducible vacuum subspaces are one-dimensional because the grade zero subalgebra, namely the Cartan subalgebra, is Abelian.. Clearly, singular vacuum spaces form representations of the algebra $`o(2)o(d1)`$ and therefore decompose into a direct sum of irreducible representations of $`o(d1)`$ on different levels. The standard situation is with a single irreducible singular vacuum space. The singular module as a whole then has the structure
$$t_+^{a_1}\mathrm{}t_+^{a_k}|E_0^{},๐ฌ^{}.$$
(18)
The factorization to a unitary irreducible representation is equivalent to identifying all these vectors to zero.
Let us now consider the example of a vacuum state $`|E_0,๐ฌ`$ with $`๐ฌ=(s_1,s_2,0,\mathrm{},0)`$ corresponding to a two-row Young diagram,
(19)
with $`0s_2s_1`$.
This means that $`|E_0,๐ฌ`$ can be realized as a tensor $`v_{a_1\mathrm{}a_{s_2},b_1\mathrm{}b_{s_1}}(E_0)`$ which is symmetric both in the indices $`a`$ and in the indices $`b`$, satisfying the antisymmetry property
$$v_{a_2\mathrm{}a_{s_2}\{b_{s_1+1},b_1\mathrm{}b_{s_1}\}}(E_0)=0,$$
(20)
which implies that symmetrization over any $`s_1+1`$ indices $`a`$ and/or $`b`$ gives zero. The tensor is traceless, which means that contraction of any two indices with the $`o(d1)`$ invariant flat metric $`\eta _{ab}`$ gives zero. Taking (20) into account it is enough to require
$$\eta ^{b_1b_2}v_{a_1\mathrm{}a_{s_2},b_1\mathrm{}b_{s_1}}(E_0)=0.$$
(21)
The level one states are
$$t_+^cv_{a_1\mathrm{}a_{s_2},b_1\mathrm{}b_{s_1}}(E_0).$$
(22)
These states form a reducible representation of $`o(d1)`$ (the tensor product of the vector representation with the representation (19)). For generic $`d`$, $`s_2`$ and $`s_1`$ it contains five irreducible components: two Young diagrams with one cell less (index $`c`$ is contracted to either one of the indices $`a`$ or one of the indices $`b`$) and three diagrams with one cell more: adding one cell to the first, second or an (additional) third row. Our problem therefore is to check whether any of these irreducible representations can be a singular vacuum space, i.e.
$$t_{}^e\mathrm{\Pi }_\alpha (t_+^cv_{a_1\mathrm{}a_{s_2},b_1\mathrm{}b_{s_1}}(E_0))0$$
(23)
for some $`E_0`$ ($`\mathrm{\Pi }_\alpha `$ is a projector to one or another irreducible component in (22)). Let us consider the two representations with cells cut. The appropriate projections are given by the following formulae describing irreducible $`o(d1)`$ tensors
$$v_{a_1\mathrm{}a_{s_2},b_1\mathrm{}b_{s_11}}^1=t_+^c\{v_{a_1\mathrm{}a_{s_2},b_1\mathrm{}b_{s_11}c}(E_0)+\frac{s_2}{s_1s_2+1}v_{ca_1\mathrm{}a_{s_21},b_1\mathrm{}b_{s_11}a_{s_2}}(E_0)\}$$
(24)
(symmetrizations within each of the groups of indices $`a`$ and $`b`$ are assumed) and
$$v_{a_1\mathrm{}a_{s_21},b_1\mathrm{}b_{s_1}}^2=t_+^cv_{ca_1\mathrm{}a_{s_21},b_1\mathrm{}b_{s_1}}(E_0).$$
(25)
Elementary computations give that
$`t_cv_{a_1\mathrm{}a_{s_2},b_1\mathrm{}b_{s_11}}^1`$ $`=`$ $`{\displaystyle \frac{1}{2}}(E_0(d+s_13))(v_{a_1\mathrm{}a_{s_2},b_1\mathrm{}b_{s_11}c}(E_0)`$ (26)
$`+`$ $`{\displaystyle \frac{s_2}{s_1s_2+1}}v_{a_1\mathrm{}a_{s_21}c,b_1\mathrm{}b_{s_11}a_{s_2}}(E_0))`$
and
$$t_cv_{a_1\mathrm{}a_{s_21},b_1\mathrm{}b_{s_1}}^2=\frac{1}{2}(E_0(d+s_24))v_{a_1\mathrm{}a_{s_21}c,b_1\mathrm{}b_{s_1}}(E_0).$$
(27)
As a result, singular vectors appear at
$$E_0^1=d+s_13$$
(28)
and
$$E_0^2=d+s_24.$$
(29)
A few comments are now in order.
The cases $`s_1=s_2`$ and $`s_1>s_2`$ are different because $`v^10`$ for $`s_1=s_2`$ as a consequence of the antisymmetry property (20). For $`s_1>s_2`$ both $`v^1`$ and $`v^2`$ are nontrivial.
As expected, the values of the โsingularโ energies (28) and (29) are in agreement with the general analysis of , where it was also shown that only the representation resulting from the factorization of the singular submodule with the highest energy of a singular vector is unitary (this fact is natural from the singular vector description: unitarity can be preserved only when the boundary of the unitarity region is approached; this implies the highest $`E_0`$). We therefore conclude that unitary massless particles appear for
$$E_0=d+s_13fors_1>s_2$$
(30)
and
$$E_0=d+s_14fors_1=s_2.$$
(31)
The analysis in terms of singular vectors is simple enough but can be simplified further with the aid of the technique proposed in with the tensors corresponding to various Young diagrams realized as certain subspaces of an appropriate Fock space. This technique is explained in section 4.1 since we will use it in the field-theoretical part of the paper. We here use the tensor language to make most clear the interpretation in terms of the representations of the massless little algebra $`o(d2)`$. One can analogously investigate singular spaces with the boxes added to make sure that they have negative energies and therefore do not play a rรดle in our analysis.
We expect that the analysis of higher levels does not affect our conclusions. One reason is that the appearance of singular vectors at higher levels within the unitarity region would imply existence of higher spin gauge fields with gauge transformations having more than one derivative acting on a gauge parameter. The general analysis of massless fields in flat space-time of an arbitrary dimension shows that this does not take place.
## 3 Flat Space Pattern of $`AdS`$ Massless fields
Let $`A_{a_1\mathrm{}a_{s_1},b_1\mathrm{}b_{s_2},\mathrm{}}`$ be an irreducible tensor of $`o(d1)`$ ($`a,b\mathrm{}=1รทd1`$) of a specific symmetry type. Let $`n^a`$ be a nonzero vector of $`o(d1)`$. $`o(d2)`$ can then be identified with the stability subalgebra of $`o(d1)`$ that leaves $`n^a`$ invariant. If the tensor $`A_{a_1\mathrm{}a_{s_1},b_1\mathrm{}b_{s_2},\mathrm{}}`$ is orthogonal to $`n^a`$ with respect to all possible contractions of indices
$$n^{a_1}A_{a_1\mathrm{}a_{s_1},b_1\mathrm{}b_{s_2},\mathrm{}}=0,n^{b_1}A_{a_1\mathrm{}a_{s_1},b_1\mathrm{}b_{s_2},\mathrm{}}=0\mathrm{}$$
(32)
then it describes a representation of $`o(d2)`$ of the same symmetry pattern. If some contractions with $`n^a`$ are nonzero, one can decompose the $`o(d1)`$ tensor $`A_{a_1\mathrm{}a_{s_1},b_1\mathrm{}b_{s_2},\mathrm{}}`$ into irreducible representations of $`o(d2)`$ with the aid of the projection operators constructed from $`n^a`$ or, in other words, performing dimensional reduction. For example, for a vector,
$$A^a=A^a+A^a,A^a=A^a\frac{n^an_b}{n_cn^c}A^b,A^a=\frac{n^an_b}{n_cn^c}A^b.$$
(33)
The analysis of singular vectors in section 2.2 admits a similar interpretation. Indeed, let us interpret $`t_+^a`$ as a vector $`n^a`$ analogous to the momentum operator in flat-space field theory (note that the operators $`t_+^a`$ commute with themselves). The fact that a singular vector appears means that some contractions of $`t_+^a`$ with the vacuum vector decouple from the spectrum and therefore are equivalent to zero. If all possible contractions of $`t_+^a`$ would decouple this would mean that the $`o(d1)`$ vacuum representation would reduce to the $`o(d2)`$ of the same symmetry. Since the energies (28) and (29) are different this cannot be true simultaneously. Therefore, when a singular vector is present, the โreduced representationโ is effectively smaller than an irreducible representation of $`o(d1)`$ but may be larger than the corresponding irreducible representation of $`o(d2)`$, containing a number of irreducible representations of $`o(d2)`$.
Let us note that the energy in $`AdS_d`$ is measured in units of the inverse $`AdS`$ radius $`\lambda `$ that was set equal to unity in our analysis. Reintroducing $`\lambda `$ and taking the flat limit $`\lambda 0`$, all energies of singular vectors tend to zero. This means that in the flat limit different singular vectors may decouple simultaneously and therefore a natural possibility consists of the flat space reduction to one (totally $`t_+^a`$ orthogonal) representation of $`o(d2)`$, in agreement with the standard analysis of massless representations of the Poincareโ algebra. However such massless representations of the Poincareโ algebra may not admit a deformation to a representation of the AdS algebra with $`\lambda 0`$. At the field-theoretical level this means that it will not be possible to preserve all necessary gauge symmetries for $`\lambda 0`$. This phenomenon is demonstrated in the field theoretical example in section 4.2. The deformation will be possible however, if one starts with an appropriate collection of massless fields in flat space dictated by the โincompleteโ dimensional reduction via decoupling of singular vectors. The main aim of this section is to formulate a conjecture on the pattern of massless fields in flat space compatible with the deformation to $`AdS_d`$.
Let us consider an arbitrary Young diagram with row lengths $`s_1s_2s_3\mathrm{}`$. For our analysis it is more convenient to build Young diagrams not from rows as elementary entities but from rectangular blocks of an arbitrary height $`t`$ and length $`s`$
(34)
In other words, a block is a Young diagram composed of $`t`$ rows of equal length $`s`$ (equivalently, from $`s`$ columns of equal height $`t`$). A general Young diagram is a combination of blocks with decreasing lengths (equivalently, heights)
(35)
In these terms, a Young diagram $`Y[(s_i,t_i)]`$ is described by a set of pairs of positive integers $`(s_i,t_i)`$ with $`s_1>s_2>s_3>\mathrm{}s_p>0`$ and arbitrary $`t_i`$ such that
$$\underset{i=1}{\overset{p}{}}t_i\frac{1}{2}(d1).$$
(36)
In other words, $`s_1`$ is the maximal row length in the Young diagram, while $`t_1`$ is the number of rows of the length $`s_1`$. $`s_2`$ is the maximal row length of the remaining rows and $`t_2`$ is the number of rows of length $`s_2`$. Note that an elementary block $`Y[(s,t)]`$ is described in these terms by a single pair of integers $`(s,t)`$.
Let us now address the question what is the result of a dimensional reduction to one dimension less of a general diagram $`Y[(s_i,t_i)]`$. Every cell can be identified with some vector index. It can either be aligned along $`n^a`$ or along the $`d2`$ perpendicular directions. In the first case we cancel a cell, while in the second case we keep it. There cannot be more than $`s_1`$ indices along $`n^a`$ because symmetrization with respect to more than $`s`$ indices gives identically zero by the definition of a Young diagram. But any number of indices from 0 to $`s_1`$ can be chosen to take the extra value $`d1`$. Therefore any number of cells from 0 to $`s_1`$ can be canceled. Of course only such cancelings are allowed that result in a Young diagram. (If not, any resulting tensor is identically zero.)
Let us consider some examples.
First consider a representation (e.g. tensor $`T`$) described by the Young diagram $`Y[(s,t)]`$ which is itself an elementary block. Since components of tensors along $`n_a`$ are automatically symmetrized because the tensor $`N^{a_1a_2\mathrm{}a_k}=n^{a_1}n^{a_2}\mathrm{}n^{a_k}`$ is totally symmetric, we can use the symmetry properties of the Young diagrams to reduce any contraction with $`N^{a_1a_2\mathrm{}a_k}`$ of the original tensor to a contraction of $`N^{a_1a_2\mathrm{}a_k}`$ with the bottom row of the block. As a result, dimensional reduction will lead to a number of tensors resulting from cutting an arbitrary number of cells in the bottom row of the block; every tensor appears once (note that the condition that the tensor is traceless does not affect this analysis since the reduced $`o(d2)`$ tensors are also assumed to be traceless).
(37)
In other words, the dimensional reduction of the block $`Y[(s,t)]`$ gives rise to the following representations of $`o(d2)`$: $`Y[(s,t)]`$ (all indices are orthogonal to $`n^a`$), $`Y[(s,t1)]`$ (a maximal possible number $`s`$ of indices is contracted) and all diagrams $`Y[(s,t1);(s_1,1)]`$ which consist of two blocks with the bottom block having an arbitrary length $`0<s_1<s`$ and height 1.
Now, consider a representation $`T`$ described by a Young diagram $`Y[(s_1,t_1);(s_2,t_2)]`$ composed from two blocks.
(38)
Again, dimensional reduction means that one can cut some cells from the bottom rows of the upper and lower blocks (all cuts inside a block are equivalent by the properties of the Young diagrams to cutting its bottom line). But now, one cannot cut an arbitrary number of boxes in the top block because the cut line cannot be shorter than the length of the second block $`s_2`$. (Such tensors vanish identically.) One can, however, take away an arbitrary number of cells from the bottom line of the second box. The rule therefore is: take away an arbitrary number $`n_1`$ such that $`s_1s_2n_10`$ from the bottom line of the top block and take away an arbitrary number $`n_2`$ such that $`s_2n_20`$ of cells from the bottom line of the bottom block
(39)
Note that with this prescription we have $`n_1+n_2s_1`$ in accordance with the general argument that one cannot cut a number of cells exceeding the maximal row length in the Young diagram. The pattern of the dimensionally reduced representation therefore consists of four-block diagrams $`Y[(s_1,t_11);(s_1^{},1);(s_2,t_21);(s_2^{},1)]`$ with arbitrary integers $`s_1^{}`$ and $`s_2^{}`$ such that
$$s_1>s_1^{}>s_2>s_2^{}>0$$
(40)
and their degenerate versions described by the three-block diagrams
$`Y[(s_1,t_11);(s_2,t_2);(s_2^{},1)]`$, $`Y[(s_1,t_1);(s_2,t_21);(s_2^{},1)]`$, $`Y[(s_1,t_11);(s_1^{},1);(s_2,t_2)]`$, $`Y[(s_1,t_11);(s_1^{},1);(s_2,t_21)]`$ and two-block diagrams $`Y[(s_1,t_1);(s_2,t_2)]`$,
$`Y[(s_1,t_11);(s_2,t_2)]`$, $`Y[(s_1,t_1);(s_2,t_21)]`$, $`Y[(s_1,t_11);(s_2,t_2+1)]`$.
Analogously one proceeds for Young diagrams built from a larger number of blocks. The final result is that the dimensional reduction of a general Young diagram to one dimension less consists of the Young diagrams of the form (every diagram appears once; see e.g. and references therein):
(41)
One is allowed to take away any numbers $`n_i`$ of cells from the $`i^{th}`$ block provided that
$$0n_is_is_{i+1}$$
(42)
(with the convention that $`s_j`$ corresponding to the โnext to lastโ block equals zero).
Let us now formulate the final result concerning a pattern of flat space massless fields that admit a unitary deformation to $`AdS_d`$.
Conjecture. Consider a $`AdS_d`$ massless field characterized by a diagram (35). Take away any number of cells $`n_i`$ satisfying the conditions (42) of the bottom lines of all blocks except for the top one, i.e. require $`n_1=0`$. Any diagram that appears as a result describes some irreducible representation of the massless little algebra $`o(d2)`$ corresponding to some flat space massless field that should be present in the full set compatible with the $`AdS_d`$ geometry and unitarity.
Note that massless fields corresponding to arbitrary irreducible representations of flat space massless little algebra $`o(d2)`$ were considered in .
A few comments are now in order.
The role of the upper block is singled out by the unitarity condition: only singular vectors corresponding to contractions of $`t_+^a`$ to the upper block have maximal energies and describe unitary representations . Therefore canceling out a box from the upper block corresponds to pure gauge (i.e. singular vector) components that decouple. Non-unitary (i.e. ghost containing) sets of fields can be obtained by a similar procedure with one of the lower blocks remaining untouched instead of the top one as in the unitary case.
For Young diagrams being themselves elementary blocks (i.e. $`t_i=0`$ for $`i>0`$) only one $`o(d2)`$ representation appears, described by the same block. This means that elementary block massless fields in $`AdS_d`$ classify according to irreducible representations of $`o(d2)`$ as in the flat case (i.e. no additional massless fields should be added to deform to $`AdS_d`$). In fact all examples of massless fields that appear in supergravity and low energy string theory are described by elementary blocks (specifically, either by single rows, or by single columns). That is why the phenomenon discussed in this paper was not observed before. Note also that for the well studied case of lower dimensions $`d4`$, only block-type massless representations are nontrivial (propagating) and therefore this phenomenon does not occur either.
The spectrum of flat-space massless fields to which an elementary $`AdS`$ massless field decomposes is non degenerate, i.e. all the representations of $`o(d2)`$ are pairwise different.
Some standard massless (gauge) fields may be needed as ingredients of $`AdS`$ massless fields with nontrivial diagrams. For example, for the representation $`Y[(2,1),(1,1)]`$ a graviton-type flat space massless field $`Y[(2,1)]`$ will be present (see example in section 4.2). It is tempting to speculate that this may correspond to a nontrivial deformation of gravity to the $`AdS`$ geometry in the presence of other fields. Analogously one can find a totally symmetric spin $`s_12`$ field corresponding to the diagram $`Y[(s_1,1)]`$ among the fields resulting from the decomposition of the $`AdS_d`$ field $`Y[(s_1,1);(s_2,1)]`$, $`0s_2s_1`$. This is to say that $`AdS_d`$ massless field corresponding to $`Y[(s_1,1);(s_2,1)]`$ decomposes into the following irreps of $`o(d2)`$ algebra
$$Y[(s_1,1);(s_2,1)]_{AdS}Y[(s_1,1)]\underset{s=1}{\overset{s_2}{}}Y[(s_1,1);(s,1)],$$
(43)
where each term under summation appears just once.
An important consequence of the analysis of this section is that totally antisymmetric gauge tensors (i.e. differential forms) corresponding to the diagrams Y\[(1,t)\] can never appear as a result of a decomposition of a certain irreducible (unitary) $`AdS_d`$ massless field in the flat limit. The space of differential forms (including the spin one gauge fields) therefore is closed with respect to the deformation to $`AdS_d`$.
## 4 Field Theoretical Example
Now let us explain what happens from the field-theoretic perspective. In fact, it has been observed already in at the level of equations of motion in the Lorentz gauge that some of the redundant gauge symmetries expected in the flat-space description are absent in the AdS case. Here we analyze the problem at the Lagrangian level focusing on the explicit comparison with the flat-space limit.
### 4.1 Flat Space
Let us consider the simplest nontrivial example of a massless field having the symmetry properties of a non-block diagram with three cells $`Y[(2,1);(1,1)]`$
(44)
In flat space the massless field of this symmetry type was described in . To begin with, let us reformulate the results of these authors in a somewhat different, although equivalent way. We take the representation with the field $`\mathrm{\Phi }_{m_1m_2,n}`$ being a symmetric tensor in $`m_1`$ and $`m_2`$, satisfying the condition that full symmetrization with respect to all three indices gives zero
$$\mathrm{\Phi }_{\{m_1m_2,m_3\}}=0.$$
(45)
(The authors of used an equivalent representation with explicit antisymmetry in two indices). The Lagrangian can be chosen to be of the form
$`L`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Phi }_{m_1m_2,n}\mathrm{}\mathrm{\Phi }^{m_1m_2,n}\mathrm{\Phi }_{mm_1,n}^{m_1}_{m_2}\mathrm{\Phi }^{mm_2,n}{\displaystyle \frac{1}{2}}\mathrm{\Phi }_{m_1m_2,n_1}^{n_1}_{n_2}\mathrm{\Phi }^{m_1m_2,n_2}`$ (46)
$``$ $`{\displaystyle \frac{3}{4}}\mathrm{\Phi }^{m_1}{}_{m_1,n}{}^{}\mathrm{}\mathrm{\Phi }_{m_2}{}_{}{}^{m_2,n}+{\displaystyle \frac{3}{4}}\mathrm{\Phi }^m{}_{m,n}{}^{}_{m_1}^{}_{m_2}\mathrm{\Phi }^{m_1m_2,n}+{\displaystyle \frac{3}{4}}\mathrm{\Phi }^{m_1m_2,n}_{m_1}_{m_2}\mathrm{\Phi }^m_{m,n}`$
$`+`$ $`{\displaystyle \frac{3}{4}}\mathrm{\Phi }^{m_1}{}_{m_1,n_1}{}^{}_{}^{n_1}_{n_2}\mathrm{\Phi }_{m_2}{}_{}{}^{m_2,n_2}.`$
The corresponding action is invariant under the gauge transformations
$`\delta _{as}\mathrm{\Phi }^{m_1m_2,n}={\displaystyle \frac{1}{2}}(^{m_1}\mathrm{\Lambda }_{as}^{m_2n}+^{m_2}\mathrm{\Lambda }_{as}^{m_1n}),`$ (47)
$`\delta _{sym}\mathrm{\Phi }^{m_1m_2,n}={\displaystyle \frac{1}{2}}(^{m_1}\mathrm{\Lambda }_{sym}^{m_2n}+^{m_2}\mathrm{\Lambda }_{sym}^{m_1n})^n\mathrm{\Lambda }_{sym}^{m_1m_2}`$ (48)
with antisymmetric gauge parameter $`\mathrm{\Lambda }_{as}^{mn}(x)`$ and symmetric gauge parameter $`\mathrm{\Lambda }_{sym}^{mn}(x)`$,
$$\mathrm{\Lambda }_{as}^{mn}(x)=\mathrm{\Lambda }_{as}^{nm}(x),\mathrm{\Lambda }_{sym}^{mn}(x)=\mathrm{\Lambda }_{sym}^{nm}(x).$$
(49)
In it was proved that the Lagrangian (46) describes a physical massless field in flat space corresponding to the irreducible representation $`Y[(2,1);(1,1)]`$ of the massless little algebra $`o(d2)`$. Because this Lagrangian is fixed by the gauge transformations one can say that it is the gauge invariance with respect to the both gauge transformations (47) and (48) that ensures irreducibility of the massless field.
Let us now introduce a notation that simplifies computation. In curved space-time it is convenient to use fiberwise fields
$$\mathrm{\Phi }^{m_1m_2m_3}e_{\underset{ยฏ}{m}_1}{}_{}{}^{m_1}e_{\underset{ยฏ}{m}_2}^{}{}_{}{}^{m_2}e_{\underset{ยฏ}{m}_3}^{}{}_{}{}^{m_3}\mathrm{\Phi }_{}^{\underset{ยฏ}{m}_1\underset{ยฏ}{m}_2\underset{ยฏ}{m}_3},$$
(50)
where $`e_{\underset{ยฏ}{m}}^m`$ is the vielbein of an appropriate space-time (e.g., Minkowski or AdS) <sup>4</sup><sup>4</sup>4Tangent space (fiber) indices $`m,n`$ and target space (base) indices $`\underset{ยฏ}{m},\underset{ยฏ}{n}`$ take the values $`0,1,\mathrm{}d1`$..
It is most convenient to formulate the action in terms of the following Fock-type generating function
$$|\mathrm{\Phi }=\frac{1}{\sqrt{2}}\mathrm{\Phi }_{m_1m_2,n}\alpha _1^{m_1}\alpha _1^{m_2}\alpha _2^n|0,$$
(51)
where $`\alpha _A^m`$ and $`\overline{\alpha }_B^n`$ are auxiliary creation and annihilation operators
$$[\overline{\alpha }_A^m,\alpha _B^n]=\eta ^{mn}\delta _{AB},[\alpha _A^m,\alpha _B^n]=0,[\overline{\alpha }_A^m,\overline{\alpha }_B^n]=0,$$
(52)
and $`|0`$ is a Fock vacuum
$$\overline{\alpha }_A^m|0=0.$$
(53)
The indices $`A,B,C,E=1,2`$ label two sets of oscillators.
The fact that we deal with the Young diagram (44) is equivalent to imposing the following constraints on the generating function $`|\mathrm{\Phi }`$
$`N_{11}|\mathrm{\Phi }=2|\mathrm{\Phi },`$ (56)
$`N_{22}|\mathrm{\Phi }=|\mathrm{\Phi },`$
$`N_{12}|\mathrm{\Phi }=0,`$
where we use the notation
$$N_{AB}\alpha _A^m\overline{\alpha }_{Bm},P_{AB}\alpha _A^m\alpha _{Bm},\overline{P}_{AB}\overline{\alpha }_A^m\overline{\alpha }_{Bm}.$$
(57)
The constraints (56) and (56) tell us that the oscillators $`\alpha _1^m`$ and $`\alpha _2^m`$ occur twice and once, respectively, on the right hand side of eq.(51). The constraint (56) is equivalent to the condition (45).
The Lorentz covariant derivative for the representation $`|\mathrm{\Phi }`$ takes the form
$$D_{\underset{ยฏ}{m}}_{\underset{ยฏ}{m}}+\frac{1}{2}\omega _{\underset{ยฏ}{m}}{}_{mn}{}^{}M_{}^{mn},M^{mn}=\underset{A=1,2}{}(\alpha _A^m\overline{\alpha }_A^n\alpha _A^n\overline{\alpha }_A^m),$$
(58)
where $`\omega _{\underset{ยฏ}{m}}^{mn}`$ is the Lorentz connection of space, while $`M^{mn}`$ forms a representation of the Lorentz algebra $`so(d1,1)`$. In the sequel we will often use the notation
$$D_A\alpha _A{}_{}{}^{m}D_{m}^{},\overline{D}_A\overline{\alpha }_A{}_{}{}^{m}D_{m}^{},D_me_m{}_{}{}^{\underset{ยฏ}{m}}D_{\underset{ยฏ}{m}}^{}.$$
(59)
The flat-space Lagrangian (46) now takes the form
$$L=\frac{1}{2}\mathrm{\Phi }|\mathrm{}D_1\overline{D}_1D_2\overline{D}_2\frac{3}{4}P_{11}\mathrm{}\overline{P}_{11}+\frac{3}{4}(P_{11}\overline{D}_1^2+D_1^2\overline{P}_{11}+P_{11}D_2\overline{D}_2\overline{P}_{11})|\mathrm{\Phi },$$
(60)
where $`D_A`$, $`\overline{D}_B`$ and $`\mathrm{}=D^mD_m`$ are defined via (58) and (59) with the flat space vielbein and Lorentz connection
$$e_{\underset{ยฏ}{m}}{}_{}{}^{n}=\delta _{\underset{ยฏ}{m}}^n,\omega _{\underset{ยฏ}{m}}{}_{}{}^{kn}=0.$$
The Lagrangian (60) is invariant under two gauge symmetries generated by the antisymmetric gauge parameter $`\mathrm{\Lambda }_{as}^{mn}`$ and the symmetric gauge parameter $`\mathrm{\Lambda }_{sym}^{mn}`$ which can be conveniently described as Fock vectors
$$|\mathrm{\Lambda }_{as}=\frac{1}{\sqrt{2}}\mathrm{\Lambda }_{asmn}\alpha _1^m\alpha _2^n|0,|\mathrm{\Lambda }_{sym}=\frac{1}{\sqrt{2}}\mathrm{\Lambda }_{symmn}\alpha _1^m\alpha _1^n|0.$$
Now the gauge transformations (47), (48) take the form
$`\delta _{as}|\mathrm{\Phi }=D_1|\mathrm{\Lambda }_{as},`$ (61)
$`\delta _{sym}|\mathrm{\Phi }=({\displaystyle \frac{1}{2}}D_1N_{21}D_2)|\mathrm{\Lambda }_{sym}.`$ (62)
### 4.2 Gauge Symmetries in $`AdS_d`$
Let us now analyze the situation in the AdS case. The Lorentz covariant derivatives (58) are no longer commuting but satisfy the commutation relationships
$$[D_{\underset{ยฏ}{m}},D_{\underset{ยฏ}{n}}]=\lambda ^2M_{\underset{ยฏ}{m}\underset{ยฏ}{n}},$$
(63)
$$D_{\underset{ยฏ}{m}}\alpha _{\underset{ยฏ}{n}A}D_{\underset{ยฏ}{n}}\alpha _{\underset{ยฏ}{m}A}=D_{\underset{ยฏ}{m}}\overline{\alpha }_{\underset{ยฏ}{n}A}D_{\underset{ยฏ}{n}}\overline{\alpha }_{\underset{ยฏ}{m}A}=0$$
(64)
with the convention
$$\alpha _{\underset{ยฏ}{n}A}=e_{\underset{ยฏ}{n}}{}_{}{}^{n}\alpha _{nA}^{},\overline{\alpha }_{\underset{ยฏ}{n}A}=e_{\underset{ยฏ}{n}}{}_{}{}^{n}\overline{\alpha }_{nA}^{}.$$
(65)
The condition (64) is just the standard zero torsion condition
$$D_{\underset{ยฏ}{m}}e_{\underset{ยฏ}{n}}{}_{}{}^{a}D_{\underset{ยฏ}{n}}e_{\underset{ยฏ}{m}}{}_{}{}^{a}=0$$
(66)
while (63) is the equation of the AdS space. The covariant DโAlembertian is
$$๐^2D_m^2+\omega _m{}_{}{}^{mn}D_{n}^{},$$
(67)
where the second term accounts for $`D_m`$ being rotated as a tangent vector. With these conventions the covariant derivatives $`D_A`$ and $`\overline{D}_B`$ satisfy a number of useful relationships summarized in the Appendix.
As in the flat case we will analyze gauge symmetries with totally symmetric and totally antisymmetric gauge parameters $`\mathrm{\Lambda }_{sym}^{mn}`$ and $`\mathrm{\Lambda }_{as}^{mn}`$. It is sometimes convenient to combine them into a gauge parameter $`\mathrm{\Lambda }_{m,n}`$ having no definite symmetry properties
$$|\mathrm{\Lambda }=\mathrm{\Lambda }_{m,n}\alpha _1^m\alpha _2^n|0.$$
The symmetric and antisymmetric parts can be singled out as
$$|\mathrm{\Lambda }_{sym}|S=N_{12}|\mathrm{\Lambda },$$
(68)
$$|\mathrm{\Lambda }_{as}(1\frac{1}{2}N_{12}N_{21})|\mathrm{\Lambda }.$$
(69)
The gauge transformation $`\delta |\mathrm{\Phi }`$ which respects the constraints (56)-(56) is
$$\delta _\mathrm{\Lambda }|\mathrm{\Phi }=D_1|\mathrm{\Lambda }D_2|\mathrm{\Lambda }_{sym}.$$
(70)
It can equivalently be written as a combination of the gauge transformations with symmetric and antisymmetric gauge parameters
$$\delta _\mathrm{\Lambda }|\mathrm{\Phi }=\delta _{as}|\mathrm{\Phi }+\delta _{sym}|\mathrm{\Phi }$$
(71)
with the gauge transformations of the form (61) and (62) but now with the derivatives $`D_1`$ and $`D_2`$ as in $`AdS_d`$.
Next we analyze whether there exists a Lagrangian that generalizes (46) (equivalently, (60)) to $`AdS_d`$. The most general deformation of (60) to the AdS case without higher derivatives is of the form
$`L^{\mathrm{\Phi }\mathrm{\Phi }}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Phi }|๐^2f\lambda ^2D_1\overline{D}_1D_2\overline{D}_2+{\displaystyle \frac{3}{4}}P_{11}(๐^2+g\lambda ^2)\overline{P}_{11}`$ (72)
$`+`$ $`{\displaystyle \frac{3}{4}}(P_{11}\overline{D}_1^2+D_1^2\overline{P}_{11}+P_{11}D_2\overline{D}_2\overline{P}_{11})|\mathrm{\Phi },`$
where $`f`$ and $`g`$ are arbitrary parameters. A straightforward but rather tedious computation with the use of the identities collected in the Appendix leads to the following result
$`\delta S^{\mathrm{\Phi }\mathrm{\Phi }}`$ $`=`$ $`\lambda ^2{\displaystyle _{AdS_d}}e[\mathrm{\Phi }|(f+3)D_1+{\displaystyle \frac{3}{2}}(d5g)P_{11}\overline{D}_1|\mathrm{\Lambda }`$ (73)
$`+`$ $`\mathrm{\Phi }|(63df)D_2+3(d3)P_{21}\overline{D}_1+{\displaystyle \frac{3}{2}}(1dg)P_{21}D_1\overline{P}_{11})|S].`$
From this expression it is clear that the freedom in the parameters $`f`$ and $`g`$ is not enough to warrant an action invariant under both types of symmetries. The best one can do is to find a Lagrangian invariant either with respect to the gauge symmetry with the parameter $`|\mathrm{\Lambda }_{as}`$ or the one with $`|\mathrm{\Lambda }_{sym}`$. Note that at the level of equations of motion in the Lorentz gauge an analogous phenomenon was observed for redundant gauge symmetries in . Since, according to the general analysis of unitary representations of $`AdS_d`$ in , the case with gauge invariance with respect to $`|\mathrm{\Lambda }_{sym}`$ does not lead to unitary dynamics, we focus on the Lagrangian possessing the $`|\mathrm{\Lambda }_{as}`$ invariance. From (73) it is obvious that this is achieved by setting
$$f=3,g=d5,$$
(74)
since the antisymmetric part of the gauge parameter enters only via the first term. Thus, we set
$`L^{\mathrm{\Phi }\mathrm{\Phi }}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Phi }|๐^2+3\lambda ^2D_1\overline{D}_1D_2\overline{D}_2+{\displaystyle \frac{3}{4}}P_{11}(๐^2+\lambda ^2(d5))\overline{P}_{11}`$ (75)
$`+`$ $`{\displaystyle \frac{3}{4}}(P_{11}\overline{D}_1^2+D_1^2\overline{P}_{11}+P_{11}D_2\overline{D}_2\overline{P}_{11})|\mathrm{\Phi }.`$
Because one of the gauge symmetries is lost, the Lagrangian (75) describes more degrees of freedom than the original flat-space Lagrangian we started with. This is in agreement with the general conclusion of Sect.3 that physical d.o.f. of massless AdS fields may not be described by an irreducible representation of $`o(d2)`$. The conjecture of Sect.3 suggests that the flat space and the AdS dynamics can match only once one starts with specific (reducible) collections of fields in flat space. From the general analysis of Sect.3 it follows that, in order to make the AdS deformation consistent for the case under consideration, one has to add a massless spin two field analogous to a graviton field.
Let us therefore introduce the field $`\chi ^{mn}`$ symmetric in indices $`m,n`$, described by the Fock vector
$$|\chi \chi _{mn}\alpha _1^m\alpha _1^n|0.$$
Since this field should describe a massless spin 2 field in the flat limit, it has its own gauge symmetry with the gauge parameter
$$|\xi \xi _m\alpha _1^m|0.$$
(76)
The idea is that starting from the sum of the free Lagrangians $`L^{\mathrm{\Phi }\mathrm{\Phi }}+L^{\chi \chi }`$ one should add cross terms $`L^{\mathrm{\Phi }\chi }`$ which (i) reestablish all (appropriately deformed by $`\lambda `$-dependent terms) gauge symmetries with the parameters $`|\mathrm{\Lambda }`$ and $`|\xi `$ and (ii) tend to zero in the flat limit. It turns out that this is indeed possible. The final result is that the action
$$S=_{AdS_d}e[L^{\mathrm{\Phi }\mathrm{\Phi }}+L^{\mathrm{\Phi }\chi }+L^{\chi \chi }],$$
(77)
where
$`L^{\mathrm{\Phi }\chi }`$ $`=`$ $`{\displaystyle \frac{3}{2}}(d3)\lambda \mathrm{\Phi }|2D_2+2P_{12}\overline{D}_1+P_{11}D_2\overline{P}_{11}|\chi ,`$ (78)
$`L^{\chi \chi }`$ $`=`$ $`{\displaystyle \frac{3}{2}}(d3)\chi |๐^2+d\lambda ^2D_1\overline{D}_1{\displaystyle \frac{1}{2}}P_{11}(๐^2+\lambda ^2)\overline{P}_{11}`$ (79)
$`+`$ $`{\displaystyle \frac{1}{2}}(P_{11}\overline{D}_1^2+D_1^2\overline{P}_{11})|\chi `$
is invariant under the gauge transformations of the form
$`\delta |\mathrm{\Phi }=D_1|\mathrm{\Lambda }D_2|S+\lambda (P_{12}P_{11}N_{21})|\xi ,`$ (80)
$`\delta |\chi =D_1|\xi +\lambda |S.`$ (81)
Note that the Lagrangian $`L^{\mathrm{\Phi }\chi }`$ is proportional to $`\lambda `$ and tends to zero in the flat limit. Therefore, as expected, the action reduces in the flat limit to the sum of two actions for the irreducible fields. In the AdS case, however, the cross term $`L^{\mathrm{\Phi }\chi }`$ becomes nontrivial so that the system does not decompose into a sum of elementary subsystems. Another comment is that according to (81) the field $`|\chi `$ becomes a Stueckelberg field that can be gauged away for $`\lambda 0`$. The resulting gauge fixed action is nothing but the action (75) invariant under the gauge symmetry with antisymmetric gauge parameter. Therefore it describes properly the irreducible AdS representation. The gauge fixing $`|\chi =0`$ is impossible however for $`\lambda =0`$. This is why the naive flat limit of the action (75) describes not two fields but only one in agreement with . This phenomenon can be interpreted as some sort of nonanalyticity of the flat limit exhibited already at the free field level.
We expect that one can analogously find a deformation of the flat space Lagrangian (60) to $`AdS_d`$ by adding an antisymmetric second rank gauge tensor. This would correspond to keeping the symmetry with symmetric parameters. However, because the corresponding representation of the AdS algebra is not unitary, the resulting gauge invariant Lagrangian is expected to have a wrong relative sign of the kinetic terms of the elementary flat-space Lagrangian (i.e., incompatible with unitarity). A similar phenomenon is expected to be true for more complicated Young diagrams: only the sets of fields predicted in Sect.3 will have all signs of kinetic terms of the flat-space Lagrangians correct.
The equations of motion for the fields $`|\mathrm{\Phi }`$ and $`|\chi `$ that follow from the Lagrangian (77) can be reduced to the form
$`\left(๐^2D_1\overline{D}_1D_2\overline{D}_2+{\displaystyle \frac{1}{2}}D_1^2\overline{P}_{11}+D_2D_1\overline{P}_{12}\lambda ^2P_{11}\overline{P}_{11}2\lambda ^2P_{12}\overline{P}_{12}+3\lambda ^2\right)|\mathrm{\Phi }`$
$`+\lambda \left((d3)(D_1N_{21}2D_2)P_{12}\overline{D}_1+P_{11}N_{21}\overline{D}_1+(P_{12}D_1P_{11}D_2)\overline{P}_{11}\right)|\chi =0,`$ (82)
$$(๐^2D_1\overline{D}_1+\frac{1}{2}D_1^2\overline{P}_{11}\frac{1}{2}P_{11}\overline{P}_{11}+d\lambda ^2)|\chi +\lambda (\overline{D}_2D_1\overline{P}_{12})|\mathrm{\Phi }=0.$$
(83)
By imposing the Lorentz gauge $`\overline{D}_i|\mathrm{\Phi }=0`$, the tracelessness condition $`\overline{P}_{ij}|\mathrm{\Phi }=0`$ and the condition $`|\chi =0`$, we are left with
$$(๐^2+3\lambda ^2)|\mathrm{\Phi }=0.$$
(84)
There is a leftover symmetry with the parameter $`|S`$ satisfying certain differential conditions. Taking into account that one can identify the Stueckelberg field $`|\chi `$ with $`|S`$ one can derive these conditions from the equations of motion for $`|\chi `$ in the gauge $`\overline{D}_i|\chi =0`$, $`\overline{P}_{ij}|\chi =0`$
$$(๐^2+d\lambda ^2)|S=0.$$
(85)
Let us compare these results with the equations for the gauge field corresponding to the $`AdS_d`$ massless representation $`D(E_0,๐ฌ)`$
$$(๐^2E_0(E_0+\lambda d\lambda )+\lambda ^2\underset{A=1,2}{}s_A)|\mathrm{\Phi }=0,$$
(86)
and the conditions on the leftover gauge parameter $`|S`$
$$\left(๐^2\lambda ^2(s_22)(s_23+d)+\lambda ^2\underset{A=1,2}{}s_A\lambda ^2\right)|S=0$$
(87)
found in . For the case under consideration the $`E_0`$ and $`๐ฌ`$ are
$$E_0=\lambda (d1),๐ฌ=(2,1,0,\mathrm{},0).$$
Plugging these values into (86) (87) we indeed arrive at the equations (84) and (85).
## 5 Conclusions
We have shown that generic irreducible massless (gauge) fields in $`AdS_d`$ in the flat limit decompose into nontrivial sets of irreducible flat space massless fields. These sets are, however, smaller than the result of a dimensional reduction to one less dimension of a corresponding massive field. In that sense $`AdS_d`$ massless fields are โless masslessโ than flat space massless fields. We made a conjecture on the pattern of the flat-space reduction of a generic $`AdS_d`$ massless field. From this conjecture it follows that there is a unique nontrivial situation when a flat space spin two massless field appears as a result of a nontrivial reduction of $`AdS_d`$ massless field with mixed symmetry properties. This example has been considered in detail. It is tempting to speculate that there may exist some new version of gravity associated with this type of field.
On the other hand we have argued that totally antisymmetric tensors can never result from the flat limit decomposition of other types of $`AdS_d`$ unitary representations. In other words, the space of differential forms is closed with respect to the flat space limit decomposition.
An interesting problem for the future is to generalize these results to the supersymmetric cases to analyze generic $`AdS`$ supermultiplets in higher dimensions and, in particular, in $`AdS_{11}`$. Another problem is to consider the multiplets occurring in to see how they group themselves in the case of $`AdS`$.
## Acknowledgments
The work of R. Metsaev and M. Vasiliev is supported in part by INTAS, Grant No.96-0538 and Grant No.99-0590 and by the RFBR Grant No.99-02-16207. R. Metsaev is also supported by DOE/ER/01545-787. M. Vasiliev was also supported by NFR F-FU 08115-347. The work of L. Brink is supported by NFR F 650-19981268/2000. R. Metsaev and M. Vasiliev would like to thank for hospitality at Chalmers University. M. Vasiliev would also like to express his gratitude to Prof. H.Nicolai for the hospitality at the MPI fรผr Gravitationsphysik, Albert Einstein Institute where some part of this work was done. The authors are grateful to O.Gelfond for the help in drawing of Young diagrams.
## Appendix. Algebra of commutators
In this appendix we collect some formulas that are used in the computations of the section 4.2.
$$[\overline{D}_A,D_B]=\delta _{AB}(๐^2+\lambda ^2\underset{C}{}N_{CC})+\lambda ^2(1d)N_{BA}+\lambda ^2\underset{C}{}(P_{BC}\overline{P}_{AC}N_{CA}N_{BC}),$$
(88)
$$[D_A,D_B]=\lambda ^2\underset{C}{}(P_{BC}N_{AC}P_{AC}N_{BC}),$$
$$[\overline{D}_A,\overline{D}_B]=\lambda ^2\underset{C}{}(N_{CB}\overline{P}_{AC}N_{CA}\overline{P}_{BC}),$$
$$[๐^2,D_A]=\lambda ^2(1d)D_A+2\lambda ^2\underset{C}{}(P_{AC}\overline{D}_CD_CN_{AC}),$$
$$[๐^2,\overline{D}_A]=\lambda ^2(d1)\overline{D}_A+2\lambda ^2\underset{C}{}(N_{CA}\overline{D}_CD_C\overline{P}_{AC}),$$
where $`๐^2`$ in (88) is a covariant DโAlembertian operator (67). These formulas can be derived by straightforward but sometimes lengthy calculation. The derivation of the following relationships:
$$[N_{AB},N_{CE}]=\delta _{BC}N_{AE}\delta _{AE}N_{CB},$$
$$[\overline{D}_A,N_{BC}]=\delta _{AB}\overline{D}_C,[N_{AB},D_C]=\delta _{BC}D_A,$$
$$[\overline{D}_A,P_{BC}]=\delta _{AB}D_C+\delta _{AC}D_B,$$
$$[\overline{P}_{AB},D_C]=\delta _{BC}\overline{D}_A+\delta _{AC}\overline{D}_B,$$
$$[\overline{P}_{AB},P_{CE}]=\delta _{BC}N_{EA}+\delta _{BE}N_{CA}+\delta _{AC}N_{EB}+\delta _{AE}N_{CB}+d(\delta _{BC}\delta _{AE}+\delta _{BE}\delta _{AC}).$$
is elementary.
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# SCALAR LEVIN-TYPE SEQUENCE TRANSFORMATIONSInvited article for J. Comput. Appl. Math.
## 1 Introduction
In applied mathematics and the numerate sciences, extrapolation methods are often used for the convergence acceleration of slowly convergent sequences or series and for the summation of divergent series. For an introduction to such methods, and also further information that cannot be covered here, see the books of Brezinski and Redivo Zaglia and Wimp and the work of Weniger and Homeier , but also the books of Baker , Baker and Graves-Morris , Brezinski , Graves-Morris , Graves-Morris, Saff and Varga , Khovanskii , Lorentzen and Waadeland , Nikishin and Sorokin , Petrushev and Popov , Ross , Saff and Varga , Wall , Werner and Buenger and Wuytack .
For the discussion of extrapolation methods, one considers a sequence $`\{\{s_n\}\}=\{\{s_0,s_1,\mathrm{}\}\}`$ with elements $`s_n`$ or the terms $`a_n=s_ns_{n1}`$ of a series $`_{j=0}^{\mathrm{}}a_j`$ with partial sums $`s_n=_{j=0}^na_j`$ for large $`n`$. A common approach is to rewrite $`s_n`$ as
$$s_n=s+R_n$$
(1)
where $`s`$ is the limit (or antilimit in the case of divergence) and $`R_n`$ is the remainder or tail. The aim then is to find a new sequence $`\{\{s_n^{}\}\}`$ such that
$$s_n^{}=s+R_n^{},R_n^{}/R_n0\text{ for }n\mathrm{}.$$
(2)
Thus, the sequence $`\{\{s_n^{}\}\}`$ converges faster to the limit $`s`$ (or diverges less violently) than $`\{\{s_n\}\}`$.
To find the sequence $`\{\{s_n^{}\}\}`$, i.e., to construct a sequence transformation $`\{\{s_n^{}\}\}=๐ฏ(\{\{s_n\}\})`$, one needs asymptotic information about the $`s_n`$ or the terms $`a_n`$ for large $`n`$, and hence about the $`R_n`$. This information then allows to eliminate the remainder at least asymptotically, for instance by substracting the dominant part of the remainder. Either such information is obtained by a careful mathematical analysis of the behavior of the $`s_n`$ and/or $`a_n`$, or it has to be extracted numerically from the values of a finite number of the $`s_n`$ and/or $`a_n`$ by some method that ideally can be proven to work for a large class of problems.
Suppose that one knows quantities $`\omega _n`$ such that $`R_n/\omega _n=O(1)`$ for $`n\mathrm{}`$, for instance
$$\underset{n\mathrm{}}{lim}R_n/\omega _n=c0$$
(3)
where $`c`$ is a constant. Such quantities are called remainder estimates. Quite often, such remainder estimates can be found with relatively low effort but the exact value of $`c`$ is often quite hard to calculate. Then, it is rather natural to rewrite the rest as $`R_n=\omega _n\mu _n`$ where $`\mu _nc`$. The problem is how to describe or model the $`\mu _n`$. Suppose that one has a system of known functions $`\psi _j(n)`$ such that $`\psi _0(n)=1`$ and $`\psi _{j+1}=o(\psi _j(n))`$ for $`j_0`$. An example of such a system is $`\psi _j(n)=(n+\beta )^j`$ for some $`\beta _+`$. Then, one may model $`\mu _n`$ as a linear combination of the $`\psi _j(n)`$ according to
$$\mu _n\underset{j=0}{\overset{\mathrm{}}{}}c_j\psi _j(n)\text{for }n\mathrm{},$$
(4)
whence the problem sequence is modelled according to
$$s_ns+\omega _n\underset{j=0}{\overset{\mathrm{}}{}}c_j\psi _j(n).$$
(5)
The idea now is to eliminate the leading terms of the remainder with the unknown constants $`c_j`$ up to $`j=k1`$, say. Thus, one uses a model sequence with elements
$$\sigma _m=\sigma +\omega _m\underset{j=0}{\overset{k1}{}}c_j\psi _j(m),m_0$$
(6)
and calculates $`\sigma `$ exactly by solving the system of $`k+1`$ equations resulting for $`m=n,n+1,\mathrm{},n+k`$ for the unknowns $`\sigma `$ and $`c_j`$, $`j=0,\mathrm{},k1`$. The solution for $`\sigma `$ is a ratio of determinants (see below) and may be denoted symbolically as
$$\sigma =T(\sigma _n,\mathrm{},\sigma _{n+k};\omega _n,\mathrm{},\omega _{n+k};\psi _j(n),\mathrm{},\psi _j(n+k)).$$
(7)
The resulting sequence transformation is
$$๐ฏ(\{\{s_n\}\},\{\{\omega _n\}\})=\{\{๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\})\}\}$$
(8)
with
$$๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\})=T(s_n,\mathrm{},s_{n+k};\omega _n,\mathrm{},\omega _{n+k};\psi _j(n),\mathrm{},\psi _j(n+k)).$$
(9)
It eliminates the leading terms of the asymptotic expansion (5). The model sequences (6) are in the kernel of the sequence transformation $`๐ฏ`$, defined as the set of all sequences such that $`๐ฏ`$ reproduces their (anti)limit exactly.
A somewhat more general approach is based on model sequences of the form
$$\sigma _n=\sigma +\underset{j=1}{\overset{k}{}}c_jg_j(n),n_0,k.$$
(10)
Virtually all known sequence transformations can be derived using such model sequences. This leads to the $`๐`$ algorithm as described below in Section 3.1. Also, some further important examples of sequence transformations are described in Section 3.
However, the introduction of remainder estimates proved to be an important theoretical step since it allows to make use of asymptotic information of the remainder easily. The most prominent of the resulting sequence transformations $`๐ฏ(\{\{s_n\}\},\{\{\omega _n\}\})`$ is the Levin transformation that corresponds to the asymptotic system of functions given by $`\psi _j(n)=(n+\beta )^j`$, and thus, to Poincare-type expansions of the $`\mu _n`$. But also other systems are of importance, like $`\psi _j(n)=1/(n+\beta )_j`$ leading to factorial series, or $`\psi _j(n)=t_n^j`$ corresponding to Taylor expansions of $`t`$โdependent functions at the abscissae $`t_n`$ that tend to zero for large $`n`$. The question which asymptotic system is best, cannot be decided generally. The answer to this question depends on the extrapolation problem. To obtain efficient extrapolation procedures for large classes of problems requires to use various asymptotic systems, and thus, a larger number of different sequence transformations. Also, different choices of $`\omega _n`$ lead to different variants of such transformations. Levin has pioneered this question and introduced three variants that are both simple and rather successful for large classes of problems. These variants and some further ones will be discussed. The question which variant is best, also cannot be decided generally. There are, however, a number of results that favor certain variants for certain problems. For example, for Stieltjes series, the choice $`\omega _n=a_{n+1}`$ can be theoretically justified (see Appendix A.).
Thus, we will focus on sequence transformations that involve an auxiliary sequence $`\{\{\omega _n\}\}`$. To be more specific, we consider transformations of the form $`๐ฏ(\{\{s_n\}\},\{\{\omega _n\}\})=\{\{๐ฏ{}_{n}{}^{(k)}\}\}`$ with
$$๐ฏ{}_{n}{}^{(k)}=\frac{{\displaystyle \underset{j=0}{\overset{k}{}}}\lambda {}_{n,j}{}^{(k)}{\displaystyle \frac{s_{n+j}}{\omega _{n+j}}}}{{\displaystyle \underset{j=0}{\overset{k}{}}}\lambda {}_{n,j}{}^{(k)}{\displaystyle \frac{1}{\omega _{n+j}}}}$$
(11)
This will be called a Levin-type transformation. The known sequence transformations that involve remainder estimates, for instance the $`๐`$, $`๐ฎ`$, and $``$ transformation of Weniger , the $`W`$ algorithm of Sidi , and the $`๐ฅ`$ transformation of Homeier with its many special cases like the important $`{}_{p}{}^{}๐`$ transformations , are all of this type. Interestingly, also the $``$, $``$, and $`๐ฆ`$ transformations of Homeier for the extrapolation of orthogonal expansions are of this type although the $`\omega _n`$ in some sense cease to be remainder estimates as defined in Eq. (3).
The Levin transformation was also generalized in a different way by Levin and Sidi who introduced the $`d^{(m)}`$ transformations. This is an important class of transformations that would deserve a thorough review itself. This, however, is outside the scope of the present review. We collect some important facts regarding this class of transformations in Section 3.2.
Levin-type transformations as defined in Eq. (11) have been used for the solution of a large variety of problems. For instance, Levin-type sequence transformations have be applied for the convergence acceleration of infinite series representations of molecular integrals , for the calculation of the lineshape of spectral holes , for the extrapolation of cluster- and crystal-orbital calculations of one-dimensional polymer chains to infinite chain length , for the calculation of special functions , for the summation of divergent and acceleration of convergent quantum mechanical perturbation series , for the evaluation of semiinfinite integrals with oscillating integrands and Sommerfeld integral tails , and for the convergence acceleration of multipolar and orthogonal expansions and Fourier series . This list is clearly not complete but sufficient to demonstrate the possibility of successful application of these transformations.
The outline of this survey is as follows: After listing some definitions and notations, we discuss some basic sequence transformations in order to provide some background information. Then, special definitions relevant for Levin-type sequence transformations are given, including variants obtained by choosing specific remainder estimates $`\omega _n`$. After this, important examples of Levin-type sequence transformations are introduced. In Section 5, we will discuss approaches for the construction of Levin-type sequence transformations, including model sequences, kernels and annihilation operators, and also the concept of hierarchical consistency. In Section 6, we derive basic properties, those of limiting transformations and discuss the application to power series. In Section 7, results on convergence acceleration are presented, while in Section 8, results on the numerical stability of the transformations are provided. Finally, we discuss guidelines for the application of the transformations and some numerical examples in Section 9.
## 2 Definitions and Notations
### 2.1 General Definitions
#### 2.1.1 Sets
##### Natural Numbers:
$$=\{1,2,3,\mathrm{}\},_0=\{0\}$$
(12)
##### Integer Numbers:
$$=\{0,1,2,3,\mathrm{}\}$$
(13)
##### Real Numbers and Vectors:
$$\begin{array}{cc}& =\{x:x\text{ real}\},\hfill \\ & _+=\{x:x>0\},\hfill \\ & ^n=\{(x_1,\mathrm{},x_n)|x_j,j=1,\mathrm{},n\}\hfill \end{array}$$
(14)
##### Complex Numbers:
$$\begin{array}{cc}& =\{z=x+\text{i}y:x,y,\text{i}^2=1\}\hfill \\ & ^n=\{(z_1,\mathrm{},z_n)|z_j,j=1,\mathrm{},n\}\hfill \end{array}$$
(15)
For $`z=x+\text{i}y`$, real and imaginary parts are denoted as $`x=\mathrm{}(z)`$, $`y=\mathrm{}(z)`$. We use $`๐`$ to denote $``$ or $``$.
##### Vectors with nonvanishing components:
$$๐ฝ^n=\{(z_1,\mathrm{},z_n)|z_j,z_j0,j=1,\mathrm{},n\}$$
(16)
##### Polynomials:
$$^k=\{P:z\underset{j=0}{\overset{k}{}}c_jz^j|z,(c_0,\mathrm{},c_k)๐^{k+1}\}$$
(17)
##### Sequences:
$$๐^๐=\{\{\{s_0,s_1,\mathrm{},s_n,\mathrm{}\}\}|s_n๐,n_0\}$$
(18)
##### Sequences with nonvanishing terms:
$$๐^๐=\{\{\{s_0,s_1,\mathrm{},s_n,\mathrm{}\}\}|s_n0,s_n๐,n_0\}$$
(19)
#### 2.1.2 Special Functions and Symbols
##### Gamma function
\[58, p. 1\]:
$$\mathrm{\Gamma }(z)=_0^{\mathrm{}}t^{z1}\mathrm{exp}(t)๐t(z_+)$$
(20)
##### Factorial:
$$n!=\mathrm{\Gamma }(n+1)=\underset{j=1}{\overset{n}{}}j$$
(21)
##### Pochhammer Symbol
\[58, p. 2\]:
$$(a)_n=\frac{\mathrm{\Gamma }(a+n)}{\mathrm{\Gamma }(a)}=\underset{j=1}{\overset{n}{}}(a+j1)$$
(22)
##### Binomial Coefficients
\[1, p. 256, Eq. (6.1.21)\]:
$$\left(\genfrac{}{}{0pt}{}{z}{w}\right)=\frac{\mathrm{\Gamma }(z+1)}{\mathrm{\Gamma }(w+1)\mathrm{\Gamma }(zw+1)}$$
(23)
##### Entier Function:
$$[[x]]=\mathrm{max}\{j:jx,x\}$$
(24)
### 2.2 Sequences, Series and Operators
#### 2.2.1 Sequences and Series
For Stieltjes series see Appendix A.
##### Scalar Sequences
with Elements $`s_n`$, Tail $`R_n`$, and Limit $`s`$:
$$\{\{s_n\}\}=\{\{s_n\}\}_{n=0}^{\mathrm{}}=\{\{s_0,s_1,s_2,\mathrm{}\}\}๐^๐,R_n=s_ns,\underset{n\mathrm{}}{lim}s_n=s.$$
(25)
If the sequence is not convergent but summable to $`s`$, $`s`$ is called the Antilimit. The $`n`$-th element $`s_n`$ of a sequence $`\sigma =\{\{s_n\}\}๐^๐`$ is also denoted by $`\sigma _n`$. A sequence is called a *constant sequence*, if all elements are constant, i.e., if there is a $`c๐`$ such that $`s_n=c`$ for all $`n_0`$, in which case it is denoted by $`\{\{c\}\}`$. The constant sequence $`\{\{0\}\}`$ is called the *zero sequence*.
##### Scalar Series
with Terms $`a_j๐`$, Partial Sums $`s_n`$, Tail $`R_n`$, and Limit/Antilimit $`s`$:
$$s=\underset{j=0}{\overset{\mathrm{}}{}}a_j,s_n=\underset{j=0}{\overset{n}{}}a_j,R_n=\underset{j=n+1}{\overset{\mathrm{}}{}}a_j=s_ns$$
(26)
We say that $`\widehat{a}_n`$ are Kummer-related to the $`a_n`$ with limit or antilimit $`\widehat{s}`$ if $`\widehat{a}_n=\mathrm{}\widehat{s}_{n1}`$ satisfy $`a_n\widehat{a}_n`$ for $`n\mathrm{}`$ and $`\widehat{s}`$ is the limit (or antilimit) of $`\widehat{s}_n=_{j=0}^n\widehat{a}_j`$.
##### Scalar Power Series
in $`z`$ with coefficients $`c_j๐`$, Partial Sums $`f_n(z)`$, Tail $`R_n(z)`$, and Limit/Antilimit $`f(z)`$:
$$f(z)=\underset{j=0}{\overset{\mathrm{}}{}}c_jz^j,f_n(z)=\underset{j=0}{\overset{n}{}}c_jz^j,R_n(z)=\underset{j=n+1}{\overset{\mathrm{}}{}}c_jz^j=f(z)f_n(z)$$
(27)
#### 2.2.2 Types of Convergence
Sequences $`\{\{s_n\}\}`$ satisfying the equation
$$\underset{n\mathrm{}}{lim}(s_{n+1}s)/(s_ns)=\rho $$
(28)
are called *linearly convergent* if $`0<\left|\rho \right|<1`$, logarithmically convergent for $`\rho =1`$ and *hyperlinearly convergent* for $`\rho =0`$. For $`\left|\rho \right|>1`$, the sequence diverges.
A sequence $`\{\{u_n\}\}`$ accelerates a sequence $`\{\{v_n\}\}`$ to $`s`$ if
$$\underset{n\mathrm{}}{lim}(u_ns)/(v_ns)=0.$$
(29)
If $`\{\{v_n\}\}`$ converges to $`s`$ then we also say that $`\{\{u_n\}\}`$ converges faster than $`\{\{v_n\}\}`$.
A sequence $`\{\{u_n\}\}`$ accelerates a sequence $`\{\{v_n\}\}`$ to $`s`$ with order $`\alpha >0`$ if
$$(u_ns)/(v_ns)=O(n^\alpha ).$$
(30)
If $`\{\{v_n\}\}`$ converges to $`s`$ then we also say that $`\{\{u_n\}\}`$ converges faster than $`\{\{v_n\}\}`$ with order $`\alpha `$.
#### 2.2.3 Operators
##### Annihilation Operator:
An operator $`๐:๐^๐๐`$ is called an annihilation operator for a given sequence $`\{\{\tau _n\}\}`$ if it satisfies
$$\begin{array}{c}๐(\{\{s_n+zt_n\}\})=๐(\{\{s_n\}\})+z๐(\{\{t_n\}\})\text{ for all }\{\{s_n\}\}๐^๐,\{\{t_n\}\}๐^๐,z๐\hfill \\ ๐(\{\{\tau _n\}\})=0\hfill \end{array}$$
(31)
##### Forward Difference Operator:
$$\begin{array}{cc}& \mathrm{}_mg(m)=g(m+1)g(m),\mathrm{}_mg_m=g_{m+1}g_m\hfill \\ & \mathrm{}_m^k=\mathrm{}_m\mathrm{}_m^{k1}\hfill \\ & \mathrm{}=\mathrm{}_n\hfill \\ & \mathrm{}^kg_n=\underset{j=0}{\overset{k}{}}(1)^{kj}\left(\genfrac{}{}{0pt}{}{k}{j}\right)g_{n+j}\hfill \end{array}$$
(32)
##### Generalized Difference Operator
$`_n^{(k)}`$ for given quantities $`\delta {}_{n}{}^{(k)}0`$:
$${}_{n}{}^{(k)}=(\delta {}_{n}{}^{(k)})^1\mathrm{}$$
(33)
##### Generalized Difference Operator
$`\stackrel{~}{}_n^{(k)}`$ for given quantities $`\zeta {}_{n}{}^{(k)}0`$:
$$\stackrel{~}{}{}_{n}{}^{(k)}=(\zeta {}_{n}{}^{(k)})^1\mathrm{}^2$$
(34)
##### Generalized Difference Operator
$`{}_{n}{}^{(k)}[\alpha ]`$ for given quantities $`\mathrm{\Delta }{}_{n}{}^{(k)}0`$:
$${}_{n}{}^{(k)}[\alpha ]f_n=(\mathrm{\Delta }{}_{n}{}^{(k)})^1(f_{n+2}2\mathrm{cos}\alpha f_{n+1}+f_n)$$
(35)
##### Generalized Difference Operator
$`{}_{n}{}^{(k)}[\zeta ]`$ for given quantities $`\stackrel{~}{\mathrm{\Delta }}{}_{n}{}^{(k)}0`$:
$${}_{n}{}^{(k)}[\zeta ]f_n=(\stackrel{~}{\mathrm{\Delta }}{}_{n}{}^{(k)})^1(\zeta {}_{n+k}{}^{(2)}f_{n+2}^{}+\zeta {}_{n+k}{}^{(1)}f_{n+1}^{}+\zeta {}_{n+k}{}^{(0)}f_{n}^{})$$
(36)
##### Weighted Difference Operators
for given $`P^{(k1)}^{k1}`$:
$$๐ฒ{}_{n}{}^{(k)}=๐ฒ{}_{n}{}^{(k)}[P^{(k1)}]=\mathrm{}^kP^{(k1)}(n)$$
(37)
##### Polynomial Operators
$`๐ซ`$for given $`P^{(k)}^k`$: Let $`P^{(k)}(x)=_{j=0}^kp{}_{j}{}^{(k)}x_{}^{j}`$. Then put
$$๐ซ[P^{(k)}]g_n=\underset{j=0}{\overset{k}{}}p{}_{j}{}^{(k)}g_{n+j}^{}.$$
(38)
##### Divided Difference Operator:
For given $`\{\{x_n\}\}`$ and $`k,n_0`$, put
$$\begin{array}{cc}& \mathrm{}{}_{n}{}^{(k)}[\{\{x_n\}\}](f(x))=\mathrm{}{}_{n}{}^{(k)}(f(x))=f[x_n,\mathrm{},x_{n+k}]=\underset{j=0}{\overset{k}{}}f(x_{n+j})\underset{\genfrac{}{}{0pt}{}{i=0}{ij}}{\overset{k}{}}\frac{1}{x_{n+j}x_{n+i}}\hfill \\ & \mathrm{}{}_{n}{}^{(k)}[\{\{x_n\}\}]g_n=\mathrm{}{}_{n}{}^{(k)}g_{n}^{}=\underset{j=0}{\overset{k}{}}g_{n+j}\underset{\genfrac{}{}{0pt}{}{i=0}{ij}}{\overset{k}{}}\frac{1}{x_{n+j}x_{n+i}}\hfill \end{array}$$
(39)
## 3 Some Basic Sequence Transformations
### 3.1 $`๐`$ Algorithm
Putting for sequences $`\{\{y_n\}\}`$ and $`\{\{g_j(n)\}\}`$, $`j=1,\mathrm{},k`$
$$E{}_{n}{}^{(k)}[\{\{y_n\}\};\{\{g_j(n)\}\}]=\left|\begin{array}{ccc}y_n& \mathrm{}& y_{n+k}\\ g_1(n)& \mathrm{}& g_1(n+k)\\ \mathrm{}& \mathrm{}& \mathrm{}\\ g_k(n)& \mathrm{}& g_k(n+k)\end{array}\right|$$
(40)
one may define the sequence transformation
$$๐{}_{n}{}^{(k)}(\{\{s_n\}\})=\frac{E{}_{n}{}^{(k)}[\{\{s_n\}\};\{\{g_j(n)\}\}]}{E{}_{n}{}^{(k)}[\{\{1\}\};\{\{g_j(n)\}\}]}$$
(41)
As is plain using Cramerโs rule, we have $`๐{}_{n}{}^{(k)}(\{\{\sigma _n\}\})=\sigma `$ if the $`\sigma _n`$ satisfy Eq. (10). Thus, the sequence transformation yields the limit $`\sigma `$ exactly for model sequences (10).
The sequence transformation $`๐`$ is known as the $`๐`$ algorithm or also as BrezinskiโHรฅvieโProtocol \[102, Sec. 10\] after two of its main investigators, Hรฅvie and Brezinski . A good introduction to this transformation is also given in the book of Brezinski and Redivo Zaglia \[14, Sec. 2.1\]. Compare also Ref. .
Numerically, the computation of the $`๐{}_{n}{}^{(k)}(\{\{s_n\}\})`$ can be performed recursively using either the algorithm of Brezinski \[14, p. 58f\]
$`๐{}_{n}{}^{(0)}(\{\{s_n\}\})=s_n,g{}_{0,i}{}^{(n)}=g_i(n),n_0,i`$
$`๐{}_{n}{}^{(k)}(\{\{s_n\}\})=๐{}_{n}{}^{(k1)}(\{\{s_n\}\}){\displaystyle \frac{๐{}_{n+1}{}^{(k1)}(\{\{s_n\}\})๐{}_{n}{}^{(k1)}(\{\{s_n\}\})}{g{}_{k1,k}{}^{(n+1)}g_{k1,k}^{(n)}}}g_{k1,k}^{(n)}`$
$`g{}_{k,i}{}^{(n)}=g{}_{k1,i}{}^{(n)}{\displaystyle \frac{g{}_{k1,i}{}^{(n+1)}g_{k1,i}^{(n)}}{g{}_{k1,k}{}^{(n+1)}g_{k1,k}^{(n)}}}g{}_{k1,k}{}^{(n)},i=k+1,k+2,\mathrm{}`$
(42)
or the algorithm of Ford and Sidi that requires additionally the quantities $`g_{k+1}(n+j)`$, $`j=0,\mathrm{},k`$ for the computation of $`๐{}_{n}{}^{(k)}(\{\{s_n\}\})`$. The algorithm of Ford and Sidi involves the quantities
$$\mathrm{\Psi }_{k,n}(u)=\frac{E{}_{n}{}^{(k)}[\{\{u_n\}\};\{\{g_j(n)\}\}]}{E{}_{n}{}^{(k)}[\{\{g_{k+1}(n)\}\};\{\{g_j(n)\}\}]}$$
(43)
for any sequence $`\{\{u_0,u_1,\mathrm{}\}\}`$ where the $`g_i(n)`$ are not changed even if they depend on the $`u_n`$ and the $`u_n`$ are changed. Then we have
$$๐{}_{n}{}^{(k)}(\{\{s_n\}\})=\frac{\mathrm{\Psi }{}_{k}{}^{(n)}(s)}{\mathrm{\Psi }{}_{k}{}^{(n)}(1)},$$
(44)
and the $`\mathrm{\Psi }`$ are calculated recursively via
$$\mathrm{\Psi }_{k,n}(u)=\frac{\mathrm{\Psi }_{k1,n+1}(u)\mathrm{\Psi }_{k1,n}(u)}{\mathrm{\Psi }_{k1,n+1}(g_{k+1})\mathrm{\Psi }_{k1,n}(g_{k+1})}.$$
(45)
Of course, for $`g_j(n)=\omega _n\psi _{j1}(n)`$, i.e., in the context of sequences modelled via the expansion (5), the $`๐`$ algorithm may be used to obtain an explicit representation for any Levin-type sequence transformation of the form (compare Eq. (9))
$$๐ฏ{}_{n}{}^{(k)}=T(s_n,\mathrm{},s_{n+k};\omega _n,\mathrm{},\omega _{n+k};\psi _j(n),\mathrm{},\psi _j(n+k))$$
(46)
as ratio of two determinants
$$๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\})=\frac{E{}_{n}{}^{(k)}[\{\{s_n/\omega _n\}\};\{\{\psi _{j1}(n)\}\}]}{E{}_{n}{}^{(k)}[\{\{1/\omega _n\}\};\{\{\psi _{j1}(n)\}\}]}.$$
(47)
This follows from the identity
$$\frac{E{}_{n}{}^{(k)}[\{\{s_n\}\};\{\{\omega _n\psi _{j1}(n)\}\}]}{E{}_{n}{}^{(k)}[\{\{1\}\};\{\{\omega _n\psi _{j1}(n)\}\}]}=\frac{E{}_{n}{}^{(k)}[\{\{s_n/\omega _n\}\};\{\{\psi _{j1}(n)\}\}]}{E{}_{n}{}^{(k)}[\{\{1/\omega _n\}\};\{\{\psi _{j1}(n)\}\}]}$$
(48)
that is an easy consequence of usual algebraic manipulations of determinants.
### 3.2 The $`d^{(m)}`$ transformations
As noted in the introduction, the $`d^{(m)}`$ transformations were introduced by Levin and Sidi as a generalization of the $`u`$ variant of the Levin transformation . We describe a slightly modified variant of the $`d^{(m)}`$ transformations :
Let $`s_r`$, $`r=0,1,\mathrm{}`$ be a real or complex sequence with limit or antilimit $`s`$ and terms $`a_0=s_0`$ and $`a_r=s_rs_{r1}`$, $`r=1,2,\mathrm{}`$ such that $`s_r=_{r=0}^ra_j`$, $`r=0,1,\mathrm{}`$. For given $`m`$ and $`\xi _l_0`$ with $`l_0`$ and $`0\xi _0<\xi _1<\xi _2<\mathrm{}`$ and $`\nu =(n_1,\mathrm{},n_m)`$ with $`n_j_0`$ the $`d^{(m)}`$ transformation yields a table of approximations $`s_\nu ^{(m,j)}`$ for the (anti-)limit $`s`$ as solution of the linear system of equations
$$s_{\xi _l}=s{}_{\nu }{}^{(m,j)}+\underset{k=1}{\overset{m}{}}(\xi _l+\alpha )^k[\mathrm{\Delta }^{k1}a_{\xi _l}]\underset{i=0}{\overset{n_k}{}}\frac{\overline{\beta }_{ki}}{(\xi _l+\alpha )^i},jlj+N,$$
(49)
with $`\alpha >0`$, $`N=_{k=1}^mn_k`$ and the $`N+1`$ unknowns $`s_\nu ^{(m,j)}`$ and $`\overline{\beta }_{ki}`$. The $`[\mathrm{\Delta }^ka_j]`$ are defined via $`[\mathrm{\Delta }^0a_j]=a_j`$ and $`[\mathrm{\Delta }^ka_j]=[\mathrm{\Delta }^{k1}a_{j+1}][\mathrm{\Delta }^{k1}a_j]`$, $`k=1,2,\mathrm{}`$. In most cases, all $`n_k`$ are chosen equal and one puts $`\nu =(n,n,\mathrm{},n)`$. Apart from the value of $`\alpha `$, only the input of $`m`$ and of $`\xi _{\mathrm{}}`$ is required from the user. As transformed sequence, often one chooses the elements $`s_{(n,\mathrm{},n)}^{(m,0)}`$ for $`n=0,1,\mathrm{}`$. The $`u`$ variant of the Levin transformation is obtained for $`m=1`$, $`\alpha =\beta `$ and $`\xi _l=l`$. The definition above differs slightly from the original one and was given in Ref. with $`\alpha =1`$.
Ford and Sidi have shown, how these transformations can be calculated recursively with the $`๐^{(m)}`$ algorithms . The $`d^{(m)}`$ transformations are the best known special cases of the *generalised Richardson Extrapolation process* (GREP) as defined by Sidi .
The $`d^{(m)}`$ transformations are derived by asymptotic analysis of the remainders $`s_rs`$ for $`r\mathrm{}`$ for the family $`\stackrel{~}{B}^{(m)}`$ of sequences $`\{\{a_r\}\}`$ as defined in Ref. . For such sequences, the $`a_r`$ satisfy a difference equation of order $`m`$ of the form
$$a_r=\underset{k=1}{\overset{m}{}}p_k(r)\mathrm{\Delta }^ka_r.$$
(50)
The $`p_k(r)`$ satisfy the asymptotic relation
$$p_k(r)r^{i_k}\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\frac{p_k\mathrm{}}{r^{\mathrm{}}}\text{for }r\mathrm{}.$$
(51)
The $`i_k`$ are integers satisfying $`i_kk`$ for $`k=1,\mathrm{},m`$. This family of sequences is very large. But still, Levin and Sidi could prove \[54, Theorem 2\] that under mild additional assumptions, the remainders for such sequences satisfy
$$s_rs\underset{k=1}{\overset{m}{}}r^{j_k}(\mathrm{\Delta }^{k1}a_r)\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\frac{\beta _k\mathrm{}}{r^{\mathrm{}}}\text{for }r\mathrm{}.$$
(52)
The $`j_k`$ are integers satisfying $`j_kk`$ for $`k=1,\mathrm{},m`$. A corresponding result for $`m=1`$ was proven by Sidi \[71, Theorem 6.1\].
The system (49) now is obtained by truncation of the expansions at $`\mathrm{}=n_n`$, evaluation at $`r=\xi _l`$, and some further obvious substitutions.
The introduction of suitable $`\xi _l`$ was shown to improve the accuracy and stability in difficult situations considerably .
### 3.3 Shanks Transformation and Epsilon Algorithm
An important special case of the $`๐`$ algorithm is the choice $`g_j(n)=\mathrm{}s_{n+j1}`$ leading to the Shanks transformation
$$e_k(s_n)=\frac{E{}_{n}{}^{(k)}[\{\{s_n\}\};\{\{\mathrm{}s_{n+j1}\}\}]}{E{}_{n}{}^{(k)}[\{\{1\}\};\{\{\mathrm{}s_{n+j1}\}\}]}$$
(53)
Instead of using one of the recursive schemes for the $`๐`$ algorithms, the Shanks transformation may be implemented using the epsilon algorithm that is defined by the recursive scheme
$$\begin{array}{cc}\hfill ฯต_1^{(n)}& =0,ฯต_0^{(n)}=s_n,\hfill \\ \hfill ฯต_{k+1}^{(n)}& =ฯต_{k1}^{(n+1)}+\mathrm{\hspace{0.17em}1}/[ฯต_k^{(n+1)}ฯต_k^{(n)}].\hfill \end{array}$$
(54)
The relations
$$ฯต_{2k}^{(n)}=e_k(s_n),ฯต_{2k+1}^{(n)}=\mathrm{\hspace{0.33em}1}/e_k(\mathrm{}s_n)$$
(55)
hold and show that the elements $`ฯต_{2k+1}^{(n)}`$ are only auxiliary quantities.
The kernel of the Shanks transformation $`e_k`$ is given by sequences of the form
$$s_n=s+\underset{j=0}{\overset{k1}{}}c_j\mathrm{}s_{n+j}.$$
(56)
See also \[14, Theorem 2.18\].
Additionally, one can use the Shanks transformation โ and hence the epsilon algorithm โ to compute the upper half of the Padรฉ table according to
$$e_k(f_n(z))=[n+k/k]_f(z),(k0,n0).$$
(57)
where
$$f_n(z)=\underset{j=0}{\overset{n}{}}c_jz^j$$
(58)
are the partial sums of a power series of a function $`f(z)`$. Padรฉ approximants of $`f(z)`$ are rational functions in $`z`$ given as ratio of two polynomials $`p_{\mathrm{}}^{(\mathrm{})}`$ and $`q_m^{(m)}`$ according to
$$[\mathrm{}/m]_f(z)=p_{\mathrm{}}(z)/q_m(z),$$
(59)
where the Taylor series of $`f`$ and $`[\mathrm{}/m]_f`$ are identical to the highest possible power of $`z`$, i.e.
$$f(z)p_{\mathrm{}}(z)/q_m(z)=O(z^{\mathrm{}+m+1}).$$
(60)
Methods for the extrapolation of power series will be treated later.
### 3.4 Aitken Process
The special case $`ฯต{}_{2}{}^{(n)}=e_1(s_n)`$ is identical to the famous $`\mathrm{\Delta }^2`$ method of Aitken
$$s{}_{n}{}^{(1)}=s_n\frac{(s_{n+1}s_n)^2}{s_{n+2}2s_{n+1}+s_n}$$
(61)
with kernel
$$s_n=s+c(s_{n+1}s_n),n_0.$$
(62)
Iteration of the $`\mathrm{\Delta }^2`$ method yields the iterated Aitken process
$$\begin{array}{cc}& ๐{}_{n}{}^{(0)}=s_n,\hfill \\ & ๐{}_{n}{}^{(k+1)}=๐{}_{n}{}^{(k)}\frac{(๐{}_{n+1}{}^{(k)}๐{}_{n}{}^{(k)})^2}{๐{}_{n+2}{}^{(k)}2๐{}_{n+1}{}^{(k)}+๐_n^{(k)}}.\hfill \end{array}$$
(63)
The iterated Aitken process and the epsilon algorithm accelerate linear convergence and can sometimes be applied successfully for the summation of alternating divergent series.
### 3.5 Overholt Process
The Overholt process is defined by the recursive scheme
$$\begin{array}{cc}\hfill V{}_{n}{}^{(0)}(\{\{s_n\}\})=& s_n,\hfill \\ \hfill V{}_{n}{}^{(k)}(\{\{s_n\}\})=& \frac{(\mathrm{}s_{n+k1})^kV{}_{n+1}{}^{(k1)}(\{\{s_n\}\})(\mathrm{}s_{n+k})^kV{}_{n}{}^{(k1)}(\{\{s_n\}\})}{(\mathrm{}s_{n+k1})^k(\mathrm{}s_{n+k})^k}\hfill \end{array}$$
(64)
for $`k`$ and $`n_0`$. It is important for the convergence acceleration of fixed point iterations.
## 4 Levin-Type Sequence Transformations
### 4.1 Definitions for Levin-Type Transformations
A set $`\mathrm{\Lambda }^{(k)}=\{\lambda {}_{n,j}{}^{(k)}๐|n_0,\mathrm{\hspace{0.25em}0}jk\}`$ is called a *coefficient set of order k* with $`k`$ if $`\lambda {}_{n,k}{}^{(k)}0`$ for all $`n_0`$. Also, $`\mathrm{\Lambda }=\{\mathrm{\Lambda }^{(k)}|k\}`$ is called *coefficient set*. Two coefficient sets $`\mathrm{\Lambda }=\{\{\lambda {}_{n,j}{}^{(k)}\}\}`$ and $`\widehat{\mathrm{\Lambda }}=\{\{\widehat{\lambda }{}_{n,j}{}^{(k)}\}\}`$ are called *equivalent*, if for all $`n`$ and $`k`$, there is a constant $`c{}_{n}{}^{(k)}0`$ such that $`\widehat{\lambda }{}_{n,j}{}^{(k)}=c{}_{n}{}^{(k)}\lambda _{n,j}^{(k)}`$ for all $`j`$ with $`0jk`$.
For each coefficient set $`\mathrm{\Lambda }^{(k)}=\{\lambda {}_{n,j}{}^{(k)}|n_0,0jk\}`$ of order $`k`$, one may define a *Levin-Type sequence transformation of order $`k`$* by
$$\begin{array}{cccc}\hfill ๐ฏ[\mathrm{\Lambda }^{(k)}]:& ๐^๐\times ๐^{(k)}\hfill & \hfill & ๐^๐\hfill \\ \hfill :& (\{\{s_n\}\},\{\{\omega _n\}\})\hfill & \hfill & \{\{s_n^{}\}\}=๐ฏ[\mathrm{\Lambda }^{(k)}](\{\{s_n\}\},\{\{\omega _n\}\})\hfill \end{array}$$
(65)
with
$$s_n^{}=๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\})=\frac{{\displaystyle \underset{j=0}{\overset{k}{}}}\lambda {}_{n,j}{}^{(k)}{\displaystyle \frac{s_{n+j}}{\omega _{n+j}}}}{{\displaystyle \underset{j=0}{\overset{k}{}}}\lambda {}_{n,j}{}^{(k)}{\displaystyle \frac{1}{\omega _{n+j}}}}$$
(66)
and
$$๐^{(k)}=\{\{\{\omega _n\}\}๐^๐:\underset{j=0}{\overset{k}{}}\lambda {}_{n,j}{}^{(k)}/\omega _{n+j}0\text{ for all }n_0\}$$
(67)
We call $`๐ฏ[\mathrm{\Lambda }]=\{๐ฏ[\mathrm{\Lambda }^{(k)}]|k\}`$ the Levin-type sequence transformation corresponding to the coefficient set $`\mathrm{\Lambda }=\{\mathrm{\Lambda }^{(k)}|k\}`$. We write $`๐ฏ^{(k)}`$ and $`๐ฏ`$ instead of $`๐ฏ[\mathrm{\Lambda }^{(k)}]`$ and $`๐ฏ[\mathrm{\Lambda }]`$, respectively, whenever the coefficients $`\lambda _{n,j}^{(k)}`$ are clear from the context. Also, if two coefficient sets $`\mathrm{\Lambda }`$ and $`\widehat{\mathrm{\Lambda }}`$ are equivalent, they give rise to the same sequence transformation, i.e., $`๐ฏ[\mathrm{\Lambda }]=๐ฏ[\widehat{\mathrm{\Lambda }}]`$, since
$$\frac{{\displaystyle \underset{j=0}{\overset{k}{}}}\widehat{\lambda }{}_{n,j}{}^{(k)}{\displaystyle \frac{s_{n+j}}{\omega _{n+j}}}}{{\displaystyle \underset{j=0}{\overset{k}{}}}\widehat{\lambda }{}_{n,j}{}^{(k)}{\displaystyle \frac{1}{\omega _{n+j}}}}=\frac{{\displaystyle \underset{j=0}{\overset{k}{}}}\lambda {}_{n,j}{}^{(k)}{\displaystyle \frac{s_{n+j}}{\omega _{n+j}}}}{{\displaystyle \underset{j=0}{\overset{k}{}}}\lambda {}_{n,j}{}^{(k)}{\displaystyle \frac{1}{\omega _{n+j}}}}\text{for }\widehat{\lambda }{}_{n,j}{}^{(k)}=c{}_{n}{}^{(k)}\lambda _n^{(k)}$$
(68)
with arbitrary $`c{}_{n}{}^{(k)}0`$.
The number $`๐ฏ_n^{(k)}`$ are often arranged in a two-dimensional table
$$\begin{array}{cccc}๐ฏ_0^{(0)}\hfill & ๐ฏ_0^{(1)}\hfill & ๐ฏ_0^{(2)}\hfill & \mathrm{}\hfill \\ ๐ฏ_1^{(0)}\hfill & ๐ฏ_1^{(1)}\hfill & ๐ฏ_1^{(2)}\hfill & \mathrm{}\hfill \\ ๐ฏ_2^{(0)}\hfill & ๐ฏ_2^{(1)}\hfill & ๐ฏ_2^{(2)}\hfill & \mathrm{}\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \end{array}$$
(69)
that is called the $`๐ฏ`$ table. The transformations $`๐ฏ^{(k)}`$ thus correspond to columns, i.e., to following vertical paths in the table. The numerators and denominators such that $`๐ฏ{}_{n}{}^{(k)}=N{}_{n}{}^{(k)}/D_n^{(k)}`$ also are often arranged in analogous $`N`$ and $`D`$ tables.
Note that for fixed $`N`$, one may also define a transformation
$$๐ฏ_N:\{\{s_{n+N}\}\}\{\{๐ฏ{}_{N}{}^{(k)}\}\}_{k=0}^{\mathrm{}}.$$
(70)
This corresponds to horizontal paths in the $`๐ฏ`$ table. These are sometimes called diagonals, because rearranging the table in such a way that elements with constant values of $`n+k`$ are members of the same row, $`๐ฏ_N^{(k)}`$ for fixed $`N`$ correspond to diagonals of the rearranged table.
For a given coefficient set $`\mathrm{\Lambda }`$ define the *moduli* by
$$\mu {}_{n}{}^{(k)}=\underset{0jk}{\mathrm{max}}\{|\lambda {}_{n,j}{}^{(k)}|\}$$
(71)
and the *characteristic polynomials* by
$$\mathrm{\Pi }{}_{n}{}^{(k)}^k:\mathrm{\Pi }{}_{n}{}^{(k)}(z)=\underset{j=0}{\overset{k}{}}\lambda {}_{n,j}{}^{(k)}z_{}^{j}$$
(72)
for $`n_0`$ and $`k`$.
Then, $`๐ฏ[\mathrm{\Lambda }]`$ is said to be *in normalized form* if $`\mu {}_{n}{}^{(k)}=1`$ for all $`k`$ and $`n_0`$. Is is said to be *in subnormalized form* if for all $`k`$ there is a constant $`\stackrel{~}{\mu }^{(k)}`$ such that $`\mu {}_{n}{}^{(k)}\stackrel{~}{\mu }^{(k)}`$ for all $`n_0`$.
Any Levin-type sequence transformation $`๐ฏ[\mathrm{\Lambda }]`$ can rewritten in normalized form. To see this, use
$$c{}_{n}{}^{(k)}=1/\mu _n^{(k)}$$
(73)
in Eq. (68). Similarly, each Levin-type sequence transformation can be rewritten in (many different) subnormalized forms.
A Levin-type sequence transformation of order $`k`$ is said to be *convex* if $`\mathrm{\Pi }{}_{n}{}^{(k)}(1)=0`$ for all $`n`$ in $`_0`$. Equivalently, it is convex if $`\{\{1\}\}๐^{(k)}`$, i.e., if the transformation vanishes for $`\{\{s_n\}\}=\{\{c\omega _n\}\}`$, $`c๐`$. Also, $`๐ฏ[\mathrm{\Lambda }]`$ is called convex, if $`๐ฏ[\mathrm{\Lambda }^{(k)}]`$ is convex for all $`k`$. We will see that this property is important for ensuring convergence acceleration for linearly convergent sequences.
A given Levin-type transformation $`๐ฏ`$ can also be rewritten as
$$๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\})=\underset{j=0}{\overset{k}{}}\gamma {}_{n,j}{}^{(k)}(\stackrel{}{\omega }_n)s_{n+j},\stackrel{}{\omega }_n=(\omega _n,\mathrm{},\omega _{n+k}),$$
(74)
with
$$\gamma {}_{n,j}{}^{(k)}(\stackrel{}{\omega }_n)=\frac{\lambda _{n,j}^{(k)}}{\omega _{n+j}}\left[\underset{j^{}=0}{\overset{k}{}}\frac{\lambda _{n,j^{}}^{(k)}}{\omega _{n+j^{}}}\right]^1,\underset{j=0}{\overset{k}{}}\gamma {}_{n,j}{}^{(k)}(\stackrel{}{\omega }_n)=1.$$
(75)
Then, one may define *stability indices* by
$$๐ช{}_{n}{}^{(k)}(๐ฏ)=\underset{j=0}{\overset{k}{}}|\gamma {}_{n,j}{}^{(k)}(\stackrel{}{\omega }_n)|1.$$
(76)
Note that any sequence transformation $`๐ฌ`$
$$๐ฌ{}_{n}{}^{(k)}=\underset{j=0}{\overset{k}{}}q{}_{n,j}{}^{(k)}s_{n+j}^{}$$
(77)
with
$$\underset{j=0}{\overset{k}{}}q{}_{n,j}{}^{(k)}=1$$
(78)
can formally be rewritten as a Levin-type sequence transformation according to $`๐ฌ{}_{n}{}^{(k)}=๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\})`$ with coefficients $`\lambda {}_{n,j}{}^{(k)}=\omega _{n+j}q{}_{n,j}{}^{(k)}\rho _n^{(k)}`$ where the validity of Eq. (78) requires to set
$$\rho {}_{n}{}^{(k)}=\underset{j=0}{\overset{k}{}}\lambda {}_{n,j}{}^{(k)}/\omega _{n+j}.$$
(79)
If for given $`k`$ and for a transformation $`๐ฏ[\mathrm{\Lambda }^{(k)}]`$ the following limits exist and have the values
$$\underset{n\mathrm{}}{lim}\lambda {}_{n,j}{}^{(k)}=\stackrel{}{๐}_j^{(k)}$$
(80)
for all $`0jk`$, and if $`\stackrel{}{\Lambda }^{(k)}`$ is a coefficient set of order $`k`$ which means that at least the limit $`\stackrel{}{๐}_k^{(k)}`$ does not vanish, then a limiting transformation $`\stackrel{}{๐ฏ}[\stackrel{}{\Lambda }{}_{}{}^{(k)}]`$ exists where $`\stackrel{}{\Lambda }{}_{}{}^{(k)}=\{\stackrel{}{๐}{}_{j}{}^{(k)}\}`$. More explicitly, we have
$$\begin{array}{cccc}\hfill \stackrel{}{๐ฏ}[\mathrm{\Lambda }^{(k)}]:& ๐^๐\times \stackrel{}{๐}^{(k)}\hfill & \hfill & ๐^๐\hfill \\ \hfill :& (\{\{s_n\}\},\{\{\omega _n\}\})\hfill & \hfill & \{\{s_n^{}\}\}\hfill \end{array}$$
(81)
with
$$s_n^{}=\stackrel{}{๐ฏ}{}_{}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\})=\frac{{\displaystyle \underset{j=0}{\overset{k}{}}}\stackrel{}{๐}{}_{j}{}^{(k)}{\displaystyle \frac{s_{n+j}}{\omega _{n+j}}}}{{\displaystyle \underset{j=0}{\overset{k}{}}}\stackrel{}{๐}{}_{j}{}^{(k)}{\displaystyle \frac{1}{\omega _{n+j}}}}$$
(82)
and
$$\stackrel{}{๐}{}_{}{}^{(k)}=\{\{\{\omega _n\}\}๐^๐:\underset{j=0}{\overset{k}{}}\stackrel{}{๐}{}_{j}{}^{(k)}/\omega _{n+j}0\text{ for all }n_0\}.$$
(83)
Obviously, this limiting transformation itself is a Levin-type sequence transformation and automatically is given in subnormalized form.
#### 4.1.1 Variants of Levin-Type Transformations
For the following, assume that $`\beta >0`$ is an arbitrary constant, $`a_n=\mathrm{}s_{n1}`$, and $`\widehat{a}_n`$ are Kummer-related to the $`a_n`$ with limit or antilimit $`\widehat{s}`$ (cp. Section 2.2.1).
A *variant* of a Levin-Type sequence transformation $`๐ฏ`$ is obtained by a particular choice $`\omega _n`$. For $`\omega _n=f_n(\{\{s_n\}\})`$, the transformation $`๐ฏ`$ is nonlinear in the $`s_n`$. In particular, we have :
##### t Variant
$${}_{}{}^{t}\omega _{n}^{}=\mathrm{}s_{n1}=a_n:{}_{}{}^{t}๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\})=๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{{}_{}{}^{t}\omega _{n}^{}\}\})$$
(84)
##### u Variant
$${}_{}{}^{u}\omega _{n}^{}=(n+\beta )\mathrm{}s_{n1}=(n+\beta )a_n:{}_{}{}^{u}๐ฏ{}_{n}{}^{(k)}(\beta ,\{\{s_n\}\})=๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{{}_{}{}^{u}\omega _{n}^{}\}\})$$
(85)
##### v Variant
$${}_{}{}^{v}\omega _{n}^{}=\frac{\mathrm{}s_{n1}\mathrm{}s_n}{\mathrm{}^2s_{n1}}=\frac{a_na_{n+1}}{a_na_{n+1}}:{}_{}{}^{v}๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\})=๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{{}_{}{}^{v}\omega _{n}^{}\}\})$$
(86)
##### $`\stackrel{~}{t}`$ Variant
$${}_{}{}^{\stackrel{~}{t}}\omega _{n}^{}=\mathrm{}s_n=a_{n+1}:{}_{}{}^{\stackrel{~}{t}}๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\})=๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{{}_{}{}^{\stackrel{~}{t}}\omega _{n}^{}\}\})$$
(87)
##### lt Variant
$${}_{}{}^{lt}\omega _{n}^{}=\widehat{a}_n:{}_{}{}^{lt}๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\})=๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{{}_{}{}^{lt}\omega _{n}^{}\}\})$$
(88)
##### lu Variant
$${}_{}{}^{lu}\omega _{n}^{}=(n+\beta )\widehat{a}_n:{}_{}{}^{lu}๐ฏ{}_{n}{}^{(k)}(\beta ,\{\{s_n\}\})=๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{{}_{}{}^{lu}\omega _{n}^{}\}\})$$
(89)
##### lv Variant
$${}_{}{}^{lv}\omega _{n}^{}=\frac{\widehat{a}_n\widehat{a}_{n+1}}{\widehat{a}_n\widehat{a}_{n+1}}:{}_{}{}^{lv}๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\})=๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{{}_{}{}^{lv}\omega _{n}^{}\}\})$$
(90)
##### l$`\stackrel{~}{t}`$ Variant
$${}_{}{}^{l\stackrel{~}{t}}\omega _{n}^{}=\widehat{a}_{n+1}:{}_{}{}^{l\stackrel{~}{t}}๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\})=๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{{}_{}{}^{l\stackrel{~}{t}}\omega _{n}^{}\}\})$$
(91)
##### K Variant
$${}_{}{}^{K}\omega _{n}^{}=\widehat{s}_n\widehat{s}:{}_{}{}^{K}๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\})=๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{{}_{}{}^{K}\omega _{n}^{}\}\})$$
(92)
The K variant of a Levin-type transformation $`๐ฏ`$ is linear in the $`s_n`$. This holds also for the lt, lu, lv and l$`\stackrel{~}{t}`$ variants.
### 4.2 Important Examples of Levin-Type Sequence Transformations
In this section, we present important Levin-type sequence transformations. For each transformation, we give the definition, recursive algorithms and some background information.
#### 4.2.1 $`๐ฅ`$ Transformation
The $`๐ฅ`$ transformation was derived and studied by Homeier. Although the $`๐ฅ`$ transformation was derived by hierarchically consistent iteration of the simple transformation
$$s_n^{}=s_{n+1}\omega _{n+1}\frac{\mathrm{}s_n}{\mathrm{}\omega _n}$$
(93)
it was possible to derive an explicit formula for its kernel as is discussed lated. It may be defined via the recursive scheme
$$\begin{array}{cc}& N{}_{n}{}^{(0)}=s_n/\omega _n,D{}_{n}{}^{(0)}=1/\omega _n,\hfill \\ & N{}_{n}{}^{(k)}={}_{n}{}^{(k1)}N{}_{n}{}^{(k1)},D{}_{n}{}^{(k)}={}_{n}{}^{(k1)}D{}_{n}{}^{(k1)},\hfill \\ & ๐ฅ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\},\{\delta {}_{n}{}^{(k)}\})=N{}_{n}{}^{(k)}/D{}_{n}{}^{(k)},\hfill \end{array}$$
(94)
where the generalized difference operator defined in Eq. (33) involves quantities $`\delta {}_{n}{}^{(k)}0`$ for $`k_0`$. Special cases of the $`๐ฅ`$ transformation result from corresponding choices of the $`\delta _n^{(k)}`$. These are summarized in Table LABEL:tabJtransSpecial.
Using generalized difference operators $`_n^{(k)}`$, we also have the representation \[36, Eq. (38)\]
$$๐ฅ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\},\{\delta {}_{n}{}^{(k)}\})=\frac{{}_{n}{}^{(k1)}{}_{n}{}^{(k2)}\mathrm{}{}_{n}{}^{(0)}[s_n/\omega _n]}{{}_{n}{}^{(k1)}{}_{n}{}^{(k2)}\mathrm{}{}_{n}{}^{(0)}[1/\omega _n]}.$$
(95)
The $`๐ฅ`$ transformation may also be computed using the alternative recursive schemes
$$\begin{array}{cc}& \widehat{D}{}_{n}{}^{(0)}=1/\omega _n,\widehat{N}{}_{n}{}^{(0)}=s_n/\omega _n,\hfill \\ & \widehat{D}{}_{n}{}^{(k)}=\mathrm{\Phi }{}_{n}{}^{(k1)}\widehat{D}{}_{n+1}{}^{(k1)}\widehat{D}{}_{n}{}^{(k1)},k,\hfill \\ & \widehat{N}{}_{n}{}^{(k)}=\mathrm{\Phi }{}_{n}{}^{(k1)}\widehat{N}{}_{n+1}{}^{(k1)}\widehat{N}{}_{n}{}^{(k1)},k,\hfill \\ & ๐ฅ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\},\{\delta {}_{n}{}^{(k)}\})=\frac{\widehat{N}_n^{(k)}}{\widehat{D}_n^{(k)}}\hfill \end{array}$$
(96)
with
$$\mathrm{\Phi }{}_{n}{}^{(0)}=1,\mathrm{\Phi }{}_{n}{}^{(k)}=\frac{\delta {}_{n}{}^{(0)}\delta {}_{n}{}^{(1)}\mathrm{}\delta _n^{(k1)}}{\delta {}_{n+1}{}^{(0)}\delta {}_{n+1}{}^{(1)}\mathrm{}\delta _{n+1}^{(k1)}},k.$$
(97)
and
$$\begin{array}{c}\stackrel{~}{D}{}_{n}{}^{(0)}=1/\omega _n,\stackrel{~}{N}{}_{n}{}^{(0)}=s_n/\omega _n,\hfill \\ \stackrel{~}{D}{}_{n}{}^{(k)}=\stackrel{~}{D}{}_{n+1}{}^{(k1)}\mathrm{\Psi }{}_{n}{}^{(k1)}\stackrel{~}{D}{}_{n}{}^{(k1)},k,\hfill \\ \stackrel{~}{N}{}_{n}{}^{(k)}=\stackrel{~}{N}{}_{n+1}{}^{(k1)}\mathrm{\Psi }{}_{n}{}^{(k1)}\stackrel{~}{N}{}_{n}{}^{(k1)},k,\hfill \\ ๐ฅ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\},\{\delta {}_{n}{}^{(k)}\})=\frac{\stackrel{~}{N}_n^{(k)}}{\stackrel{~}{D}_n^{(k)}}\hfill \end{array}$$
(98)
with
$$\mathrm{\Psi }{}_{n}{}^{(0)}=1,\mathrm{\Psi }{}_{n}{}^{(k)}=\frac{\delta {}_{n+k}{}^{(0)}\delta {}_{n+k1}{}^{(1)}\mathrm{}\delta _{n+1}^{(k1)}}{\delta {}_{n+k1}{}^{(0)}\delta {}_{n+k2}{}^{(1)}\mathrm{}\delta _n^{(k1)}},k.$$
(99)
The quantities $`\mathrm{\Psi }_n^{(k)}`$ should not be mixed up with the $`\mathrm{\Psi }_{k,n}(u)`$ as defined in Eq. (43).
As shown in , the coefficients for the algorithm (96) that are defined via $`\widehat{D}{}_{n}{}^{(k)}=_{j=0}^k\lambda {}_{n,j}{}^{(k)}/\omega _{n+j}`$, satisfy the recursion
$$\lambda {}_{n,j}{}^{(k+1)}=\mathrm{\Phi }{}_{n}{}^{(k)}\lambda {}_{n+1,j1}{}^{(k)}\lambda _{n,j}^{(k)}$$
(100)
with starting values $`\lambda {}_{n,j}{}^{(0)}=1`$. This holds for all $`j`$ if we define $`\lambda {}_{n,j}{}^{(k)}=0`$ for $`j<0`$ or $`j>k`$. Because $`\mathrm{\Phi }{}_{n}{}^{(k)}0`$, we have $`\lambda {}_{n,k}{}^{(k)}0`$ such that $`\{\lambda {}_{n,j}{}^{(k)}\}`$ is a coefficient set for all $`k_0`$.
Similarly, the coefficients for the algorithm (98) that are defined via $`\stackrel{~}{D}{}_{n}{}^{(k)}=_{j=0}^k\stackrel{~}{\lambda }{}_{n,j}{}^{(k)}/\omega _{n+j}`$, satisfy the recursion
$$\stackrel{~}{\lambda }{}_{n,j}{}^{(k+1)}=\stackrel{~}{\lambda }{}_{n+1,j1}{}^{(k)}\mathrm{\Psi }{}_{n}{}^{(k)}\stackrel{~}{\lambda }_{n,j}^{(k)}$$
(101)
with starting values $`\stackrel{~}{\lambda }{}_{n,j}{}^{(0)}=1`$. This holds for all $`j`$ if we define $`\stackrel{~}{\lambda }{}_{n,j}{}^{(k)}=0`$ for $`j<0`$ or $`j>k`$. In this case, we have $`\stackrel{~}{\lambda }{}_{n,k}{}^{(k)}=1`$ such that $`\{\stackrel{~}{\lambda }{}_{n,j}{}^{(k)}\}`$ is a coefficient set for all $`k_0`$.
Since the $`๐ฅ`$ transformation vanishes for $`\{\{s_n\}\}=\{\{c\omega _n\}\}`$, $`c๐`$ according to Eq. (95) for all $`k`$, it is convex. This may also be shown by using induction in $`k`$ using $`\lambda {}_{n,1}{}^{(1)}=\lambda {}_{n,0}{}^{(1)}=1`$ and the equation
$$\underset{j=0}{\overset{k+1}{}}\lambda {}_{n,j}{}^{(k+1)}=\mathrm{\Phi }{}_{n}{}^{(k)}\underset{j=0}{\overset{k}{}}\lambda {}_{n+1,j}{}^{(k)}\underset{j=0}{\overset{k}{}}\lambda _{n,j}^{(k)}$$
(102)
that follows from Eq. (100).
Assuming that the limits $`\mathrm{\Phi }_k=lim_n\mathrm{}\mathrm{\Phi }_n^{(k)}`$ exist for all $`k`$ and noting that for $`k=0`$ always $`\mathrm{\Phi }_0=1`$ holds, it follows that there exists a limiting transformation $`\stackrel{}{๐ฅ}[\stackrel{}{\Lambda }]`$ that can be considered as special variant of the $`๐ฅ`$ transformation and with coefficients given explicitly as \[46, Eq. (16)\]
$$\stackrel{}{๐}{}_{j}{}^{(k)}=(1)^{kj}\underset{\genfrac{}{}{0pt}{}{j_0+j_1+\mathrm{}+j_{k1}=j,}{j_0\{0,1\},\mathrm{},j_{k1}\{0,1\}}}{}\underset{m=0}{\overset{k1}{}}(\mathrm{\Phi }_m)^{j_m}.$$
(103)
As characteristic polynomial we obtain
$$\stackrel{}{\Pi }{}_{}{}^{(k)}(z)=\underset{j=0}{\overset{k}{}}\stackrel{}{๐}{}_{j}{}^{(k)}z_{}^{j}=\underset{j=0}{\overset{k1}{}}(\mathrm{\Phi }_jz1).$$
(104)
Hence, the $`\stackrel{}{๐ฅ}`$ transformation is convex since $`\stackrel{}{\Pi }{}_{}{}^{(k)}(1)=0`$ due to $`\mathrm{\Phi }_0=1`$.
##### The $`{}_{p}{}^{}๐`$ Transformation.
This is the special case of the $`๐ฅ`$transformation corresponding to
$$\delta {}_{n}{}^{(k)}=\frac{1}{\left(n+\beta +(p1)k\right)_2},$$
(105)
or to \[46, Eq. (18)\]<sup>1</sup><sup>1</sup>1The equation in contains an error.
$$\mathrm{\Phi }{}_{n}{}^{(k)}=\{\begin{array}{cc}\left(\frac{n+\beta +2}{p1}\right)_k/\left(\frac{n+\beta }{p1}\right)_k\hfill & \text{for }p1\hfill \\ & \\ \left(\frac{n+\beta +2}{n+\beta }\right)^k\hfill & \text{for }p=1\hfill \end{array}$$
(106)
or to
$$\mathrm{\Psi }{}_{n}{}^{(k)}=\{\begin{array}{cc}\left(\frac{n+\beta +k1}{p2}\right)_k/\left(\frac{n+\beta +k+1}{p2}\right)_k\hfill & \text{for }p2\hfill \\ & \\ \left(\frac{n+\beta +k1}{n+\beta +k+1}\right)^k\hfill & \text{for }p=2\hfill \end{array}$$
(107)
that is,
$${}_{p}{}^{}๐{}_{n}{}^{(k)}(\beta ,\{\{s_n\}\},\{\{\omega _n\}\})=๐ฅ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\},\{1/(n+\beta +(p1)k)_2\}).$$
(108)
The limiting transformation $`{}_{p}{}^{}\stackrel{}{๐}`$ of the $`{}_{p}{}^{}๐`$ transformation exists for all $`p`$ and corresponds to the $`\stackrel{}{๐ฅ}`$ transformation with $`\mathrm{\Phi }_k=1`$ for all $`k`$ in $`_0`$. This is exactly the Drummond transformation discussed in Section 4.2.2, i.e., we have
$${}_{p}{}^{}\stackrel{}{๐}{}_{n}{}^{(k)}(\beta ,\{\{s_n\}\},\{\{\omega _n\}\})=๐{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\}).$$
(109)
#### 4.2.2 Drummond Transformation
This transformation was given by Drummond . It was also discussed by Weniger . It may be defined as
$$๐{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\})=\frac{\mathrm{}^k[s_n/\omega _n]}{\mathrm{}^k[1/\omega _n]}$$
(110)
Using the definition (32) of the forward difference operator, the coefficients may be taken as
$$\lambda {}_{n,j}{}^{(k)}=(1)^j\left(\genfrac{}{}{0pt}{}{k}{j}\right),$$
(111)
i.e., independent of $`n`$. As moduli, one has $`\mu {}_{n}{}^{(k)}=\left(\genfrac{}{}{0pt}{}{k}{[[k/2]]}\right)=\stackrel{~}{\mu }^{(k)}`$. Consequently, the Drummond transformation is given in subnormalized form. As characteristic polynomial we obtain
$$\mathrm{\Pi }{}_{n}{}^{(k)}(z)=\underset{j=0}{\overset{k}{}}(1)^j\left(\genfrac{}{}{0pt}{}{k}{j}\right)z^j=(1z)^k.$$
(112)
Hence, the Drummond transformation is convex since $`\mathrm{\Pi }{}_{n}{}^{(k)}(1)=0`$. Interestingly, the Drummond transformation is identical to its limiting transformation:
$$\stackrel{}{๐}{}_{}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\})=๐{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\}).$$
(113)
The Drummond transformation may be computed using the recursive scheme
$$\begin{array}{cc}& N{}_{n}{}^{(0)}=s_n/\omega _n,D{}_{n}{}^{(0)}=1/\omega _n,\hfill \\ & N{}_{n}{}^{(k)}=\mathrm{}N{}_{n}{}^{(k1)},D{}_{n}{}^{(k)}=\mathrm{}D{}_{n}{}^{(k1)},\hfill \\ & ๐{}_{n}{}^{(k)}=N{}_{n}{}^{(k)}/D{}_{n}{}^{(k)}.\hfill \end{array}$$
(114)
#### 4.2.3 Levin Transformation
This transformation was given by Levin . It was also discussed by Weniger . It may be defined as<sup>2</sup><sup>2</sup>2Note that the order of indices is different from that in the literature.
$${}_{n}{}^{(k)}(\beta ,\{\{s_n\}\},\{\{\omega _n\}\})=\frac{(n+\beta +k)^{1k}\mathrm{}^k[(n+\beta )^{k1}s_n/\omega _n]}{(n+\beta +k)^{1k}\mathrm{}^k[(n+\beta )^{k1}/\omega _n]}$$
(115)
Using the definition (32) of the forward difference operator, the coefficients may be taken as
$$\lambda {}_{n,j}{}^{(k)}=(1)^j\left(\genfrac{}{}{0pt}{}{k}{j}\right)(n+\beta +j)^{k1}/(n+\beta +k)^{k1},$$
(116)
The moduli satisfy $`\mu {}_{n}{}^{(k)}\left(\genfrac{}{}{0pt}{}{k}{[[k/2]]}\right)=\stackrel{~}{\mu }^{(k)}`$ for given $`k`$. Consequently, the Levin transformation is given in subnormalized form. As characteristic polynomial we obtain
$$\mathrm{\Pi }{}_{n}{}^{(k)}(z)=\underset{j=0}{\overset{k}{}}(1)^j\left(\genfrac{}{}{0pt}{}{k}{j}\right)z^j(n+\beta +j)^{k1}/(n+\beta +k)^{k1}$$
(117)
Since $`\mathrm{\Pi }{}_{n}{}^{(k)}(1)=0`$ because $`\mathrm{}^{k1}`$ annihilates any polynomial in $`n`$ with degree less than $`k`$, the Levin transformation is convex. The limiting transformation is identical to the Drummond transformation
$$\stackrel{}{}{}_{}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\})=๐{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\}).$$
(118)
The Levin transformation may be computed using the recursive scheme , \[14, Sec. 2.7\]
$$\begin{array}{cc}& N{}_{n}{}^{(0)}=s_n/\omega _n,D{}_{n}{}^{(0)}=1/\omega _n,\hfill \\ & N{}_{n}{}^{(k)}=N{}_{n+1}{}^{(k1)}\frac{(\beta +n)(\beta +n+k1)^{k2}}{(\beta +n+k)^{k1}}N{}_{n}{}^{(k1)},\hfill \\ & D{}_{n}{}^{(k)}=D{}_{n+1}{}^{(k1)}\frac{(\beta +n)(\beta +n+k1)^{k2}}{(\beta +n+k)^{k1}}D{}_{n}{}^{(k1)},\hfill \\ & {}_{n}{}^{(k)}(\beta ,\{\{s_n\}\},\{\{\omega _n\}\})=N{}_{n}{}^{(k)}/D{}_{n}{}^{(k)}.\hfill \end{array}$$
(119)
This is essentially the same as the recursive scheme (98) for the $`๐ฅ`$ transformation with
$$\mathrm{\Psi }{}_{n}{}^{(k)}=\frac{(\beta +n)(\beta +n+k)^{k1}}{(\beta +n+k+1)^k}.$$
(120)
since the Levin transformation is a special case of the $`๐ฅ`$ transformation (see Table LABEL:tabJtransSpecial). Thus, the Levin transformation can also be computed recursively using the scheme (94)
$$\delta {}_{n}{}^{(k)}=\frac{1}{(n+\beta )(n+\beta +k+1)},$$
(121)
or the scheme (96) with
$$\mathrm{\Phi }{}_{n}{}^{(k)}=(n+\beta +k+1)\frac{(n+\beta +1)^{k1}}{(n+\beta )^k}.$$
(122)
#### 4.2.4 Weniger Transformations
Weniger derived sequence transformations related to factorial series. These may be regarded as special cases of the transformation
$$๐{}_{n}{}^{(k)}(\alpha ,\zeta ,\{\{s_n\}\},\{\{\omega _n\}\})=\frac{((\alpha [n+\zeta +k])_{k1})^1\mathrm{}^k[(\alpha [n+\zeta ])_{k1}s_n/\omega _n]}{((\alpha [n+\zeta +k])_{k1})^1\mathrm{}^k[(\alpha [n+\zeta ])_{k1}/\omega _n]}.$$
(123)
In particular, the Weniger $`๐ฎ`$ transformation may be defined as
$$๐ฎ{}_{n}{}^{(k)}(\beta ,\{\{s_n\}\},\{\{\omega _n\}\})=๐{}_{n}{}^{(k)}(1,\beta ,\{\{s_n\}\},\{\{\omega _n\}\})$$
(124)
and the Weniger $``$ transformation as
$${}_{n}{}^{(k)}(\xi ,\{\{s_n\}\},\{\{\omega _n\}\})=๐{}_{n}{}^{(k)}(1,\xi ,\{\{s_n\}\},\{\{\omega _n\}\}).$$
(125)
The parameters $`\beta `$, $`\xi `$, and $`\zeta `$ are taken to be positive real numbers. Weniger considered the $`๐`$ transformation only for $`\alpha >0`$ and thus, he was not considering the $``$ transformation as a special case of the $`๐`$ transformation. He also found that one should choose $`\xi k1`$. In the $`u`$ variant of the $``$ transformation he proposed to choose $`\omega _n=(n\xi )\mathrm{}s_{n1}`$. This variant is denoted as $`{}_{}{}^{u}`$ transformation in the present work.
Using the definition (32) of the forward difference operator, the coefficients may be taken as
$$\lambda {}_{n,j}{}^{(k)}=(1)^j\left(\genfrac{}{}{0pt}{}{k}{j}\right)(\alpha [n+\zeta +j])_{k1}/(\alpha [n+\zeta +k])_{k1}$$
(126)
in the case of the $`๐`$ transformation, as
$$\lambda {}_{n,j}{}^{(k)}=(1)^j\left(\genfrac{}{}{0pt}{}{k}{j}\right)(n+\beta +j)_{k1}/(n+\beta +k)_{k1}$$
(127)
in the case of the $`๐ฎ`$ transformation, and as
$$\lambda {}_{n,j}{}^{(k)}=(1)^j\left(\genfrac{}{}{0pt}{}{k}{j}\right)(n\xi j)_{k1}/(n\xi k)_{k1}$$
(128)
in the case of the $``$ transformation.
The $`๐ฎ`$ transformation in (124) may be computed using the recursive scheme (98) with \[84, Sec. 8.3\]
$$\mathrm{\Psi }{}_{n}{}^{(k)}=\frac{(\beta +n+k)(\beta +n+k1)}{(\beta +n+2k)(\beta +n+2k1)}$$
(129)
The $``$ transformation in (125) may be computed using the recursive scheme (98) with \[84, Sec. 9.3\]
$$\mathrm{\Psi }{}_{n}{}^{(k)}=\frac{\xi +nk+1}{\xi +n+k+1}$$
(130)
The $`๐`$ transformation in (123) may be computed using the recursive scheme (98) with \[87, Eq. (3.3)\]
$$\mathrm{\Psi }{}_{n}{}^{(k)}=(\alpha [\zeta +n]+k2)\frac{(\alpha [n+\zeta +k1])_{k2}}{(\alpha [n+\zeta +k])_{k1}}$$
(131)
Since the operator $`\mathrm{}^k`$ for $`k`$ annihilates all polynomials in $`n`$ of degree smaller than $`k`$, the transformations $`๐ฎ`$, $``$, and $`๐`$ are convex. The moduli satisfy $`\mu {}_{n}{}^{(k)}\left(\genfrac{}{}{0pt}{}{k}{[[k/2]]}\right)=\stackrel{~}{\mu }^{(k)}`$ for given $`k`$. Consequently, the three Weniger transformations are given in subnormalized form.
For $`\alpha \mathrm{}`$, the Levin transformation is obtained from the $`๐`$ transformation . The $`๐ฎ`$ transformation is identical to the $`{}_{3}{}^{}๐`$ transformation. It is also the special case $`x_n=n+\beta `$ of the $``$ transformation. Analogously, the $`๐`$ transformation is obtained for $`x_n=\alpha [\zeta +n]`$. All these Weniger transformations are special cases of the $`๐ฅ`$ transformation (compare Table LABEL:tabJtransSpecial).
The limiting transformation of all these Weniger transformations is the Drummond transformation.
#### 4.2.5 Levin-Sidi Transformations and W Algorithms
As noted above in Section 3.2, the $`d^{(m)}`$ transformations were introduced by Levin and Sidi as a generalization of the $`u`$ variant of the Levin transformation, and these transformations may be implemented recursively using the $`๐^{(m)}`$ algorithms.
The case $`m=1`$ corresponding to the $`d^{(1)}`$ transformation and the $`๐^{(1)}=W`$ algorithm is relevant for the present survey of Levin-type transformations. In the following, the $`k`$-th order transformation $`๐ฏ^{(k)}`$ of Levin-type transformation $`๐ฏ`$as given by the $`W`$ algorithm is denoted by $`W^{(k)}`$ which should not be confused with the $`๐^{(m)}`$ algorithms of Ford and Sidi .
The $`W`$ algorithm was also studied by other authors \[84, Sec. 7.4\], \[14, p. 71f, 116f\] and may be regarded as a special case of the $`๐ฅ`$ transformation . It may be defined as (compare \[78, Theorems 1.1 and 1.2\])
$$\begin{array}{cc}& N{}_{n}{}^{(0)}=\frac{s_n}{\omega _n},D{}_{n}{}^{(0)}=\frac{1}{\omega _n},\hfill \\ & N{}_{n}{}^{(k)}=\frac{N{}_{n+1}{}^{(k1)}N_n^{(k1)}}{t_{n+k}t_n},\hfill \\ & D{}_{n}{}^{(k)}=\frac{D{}_{n+1}{}^{(k1)}D_n^{(k1)}}{t_{n+k}t_n},\hfill \\ & W{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\},\{\{t_n\}\})=N{}_{n}{}^{(k)}/D{}_{n}{}^{(k)},\hfill \end{array}$$
(132)
and computes
$$W{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\},\{\{t_n\}\})=\frac{\mathrm{}{}_{n}{}^{(k)}(s_n/\omega _n)}{\mathrm{}{}_{n}{}^{(k)}(1/\omega _n)}$$
(133)
where the divided difference operators $`\mathrm{}{}_{n}{}^{(k)}=\mathrm{}{}_{n}{}^{(k)}[\{\{t_n\}\}]`$ are used. The $`W`$ algorithm may be used to calculate the Levin transformation on putting $`t_n=1/(n+\beta )`$. Some authors call a linear variant of the $`W`$ algorithm with $`\omega _n=(1)^{n+1}e^{n\zeta q}t_n^\alpha `$ the $`W`$ transformation, while the $`\stackrel{~}{t}`$ variant of the $`W`$ algorithm is sometimes called $`mW`$ transformation .
If $`t_{n+1}/t_n\tau `$ for large $`n`$, one obtains as limiting transformation the $`\stackrel{}{๐ฅ}`$ transformation with $`\mathrm{\Phi }_j=\tau ^j`$ and characteristic polynomial
$$\stackrel{}{\Pi }{}_{}{}^{(k)}(z)=\underset{j=0}{\overset{k1}{}}(z/\tau ^j1).$$
(134)
For the $`d^{(1)}`$ transformation, we write
$$(d^{(1)}){}_{n}{}^{(k)}(\alpha ,\{\{s_n\}\},\{\{\xi _n\}\})=W{}_{n}{}^{(k)}(\{\{s_{\xi _n}\}\},\{\{(\xi _n+\alpha )(s_{\xi _n}s_{\xi _n1})\}\},\{\{1/(\xi _n+\alpha )\}\})$$
(135)
Thus, it corresponds to the variant of the $`W`$ algorithm with remainder estimates chosen as $`(\xi _n+\alpha )(s_{\xi _n}s_{\xi _n1})`$ operating on the subsequence $`\{\{s_{\xi _n}\}\}`$ of $`\{\{s_n\}\}`$ with $`t_n=1/(\xi _n+\alpha )`$. It should be noted that this is not(!) identical to the $`u`$ variant
$${}_{}{}^{u}W{}_{n}{}^{(k)}(\{\{s_{\xi _n}\}\},\{\{1/(\xi _n+\alpha )\}\})=W{}_{n}{}^{(k)}(\{\{s_{\xi _n}\}\},\{\{{}_{}{}^{u}\omega _{n}^{}\}\},\{\{1/(\xi _n+\alpha )\}\}),$$
(136)
neither for $`{}_{}{}^{u}\omega _{n}^{}=(n+\alpha )(s_{\xi _n}s_{\xi _{n1}})`$ nor for $`{}_{}{}^{u}\omega _{n}^{}=(\xi _n+\alpha )(s_{\xi _n}s_{\xi _{n1}})`$, since the remainder estimates are chosen differently in Eq. (135).
The $`d^{(1)}`$ transformation was thoroughly analyzed by Sidi (See and references therein).
#### 4.2.6 Mosig-Michalski Transformation
The Mosig-Michalski transformation โ also known as โweightedโaverages algorithmโ โ was introduced by Mosig and modified later by Michalski who gave the $`\stackrel{~}{t}`$ variant of the transformation the name $`๐ฆ`$ transformation (that is used for a different transformation in the present article(!)), and applied it to the computation of Sommerfeld integrals .
The Mosig-Michalski transformation $`M`$ may be defined via the recursive scheme
$$\begin{array}{c}s{}_{n}{}^{(0)}=s_n,\hfill \\ s{}_{n}{}^{(k+1)}=\frac{s{}_{n}{}^{(k)}+\eta {}_{n}{}^{(k)}s_{n+1}^{(k)}}{1+\eta _n^{(k)}},\hfill \\ M{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\},\{\{x_n\}\})=s_n^{(k)}\hfill \end{array}$$
(137)
for $`n_0`$ and $`k_0`$ where $`\{\{x_n\}\}`$ is an auxiliary sequence with $`lim_n\mathrm{}1/x_n=0`$ such that $`x_{n+\mathrm{}}>x_n`$ for $`\mathrm{}`$ and $`x_0>1`$, i.e., a diverging sequence of monotonously increasing positive numbers, and
$$\eta {}_{n}{}^{(k)}=\frac{\omega _n}{\omega _{n+1}}\left(\frac{x_{n+1}}{x_n}\right)^{2k}.$$
(138)
Putting $`\omega {}_{n}{}^{(k)}=\omega _n/x_n^{2k}`$, $`N{}_{n}{}^{(k)}=s{}_{n}{}^{(k)}/\omega _n^{(k)}`$, and $`D{}_{n}{}^{(k)}=1/\omega _n^{(k)}`$, it is easily seen that the recursive scheme (137) is equivalent to the scheme (94) with
$$\delta {}_{n}{}^{(k)}=\frac{1}{x_n^2}(1\frac{\omega _nx_{n+1}^{2k}}{\omega _{n+1}x_n^{2k}}).$$
(139)
Thus, the Mosig-Michalski transformation is a special case of the $`๐ฅ`$ transformation. Its character as a Levin-type transformation is somewhat formal since the $`\delta _n^{(k)}`$ and, hence, the coefficients $`\lambda _{n,j}^{(k)}`$ depend on the $`\omega _n`$.
If $`x_{n+1}/x_n\xi >1`$ for large $`n`$, then a limiting transformation exists, namely
$`M(\{\{s_n\}\},\{\{\omega _n\}\},\{\{\xi ^{n+1}\}\})`$. It corresponds to the $`\stackrel{}{๐ฅ}`$ transformation with $`\mathrm{\Phi }_k=\xi ^{2k}`$. This may be seen by putting $`\widehat{D}{}_{n}{}^{(k)}=1/\omega _n`$, $`\widehat{N}{}_{n}{}^{(k)}=s{}_{n}{}^{(k)}D_n^{(k)}`$ and $`\mathrm{\Phi }{}_{n}{}^{(k)}=\xi ^{2k}`$ in Eq. (96).
#### 4.2.7 $``$ Transformation
This transformation is seemingly new. It will be derived in a later section. It may be defined as
$${}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\},\{\{x_n\}\})=\frac{\mathrm{}{}_{n}{}^{(k)}((x_n)_{k1}s_n/\omega _n)}{\mathrm{}{}_{n}{}^{(k)}((x_n)_{k1}/\omega _n)}=\frac{x_n^k/(x_n)_{k1}\mathrm{}{}_{n}{}^{(k)}((x_n)_{k1}s_n/\omega _n)}{x_n^k/(x_n)_{k1}\mathrm{}{}_{n}{}^{(k)}((x_n)_{k1}/\omega _n)}$$
(140)
where $`\{\{x_n\}\}`$ is an auxiliary sequence with $`lim_n\mathrm{}1/x_n=0`$ such that $`x_{n+\mathrm{}}>x_n`$ for $`\mathrm{}`$ and $`x_0>1`$, i.e., a diverging sequence of monotonously increasing positive numbers. Using the definition (39) of the divided difference operator $`\mathrm{}{}_{n}{}^{(k)}=\mathrm{}{}_{n}{}^{(k)}[\{\{x_n\}\}]`$, the coefficients may be taken as
$$\lambda {}_{n,j}{}^{(k)}=\frac{(x_{n+j})_{k1}}{(x_n)_{k1}}\underset{\genfrac{}{}{0pt}{}{i=0}{ij}}{\overset{k}{}}\frac{x_n}{x_{n+j}x_{n+i}}=\underset{m=0}{\overset{k2}{}}\frac{x_{n+j}+m}{x_n+m}\left(\frac{x_n}{x_{n+j}}\right)^k\underset{\genfrac{}{}{0pt}{}{i=0}{ij}}{\overset{k}{}}\frac{1}{1x_{n+i}/x_{n+j}}.$$
(141)
Assuming that the following limit exists such that
$$\underset{n\mathrm{}}{lim}\frac{x_{n+1}}{x_n}=\xi >1$$
(142)
holds, we see that one can define a limiting transformation $`\stackrel{}{}^{(k)}`$ with coefficients
$$\stackrel{}{๐}{}_{j}{}^{(k)}=\underset{n\mathrm{}}{lim}\lambda {}_{n,j}{}^{(k)}=\frac{1}{\xi ^j}\underset{\genfrac{}{}{0pt}{}{\mathrm{}=0}{\mathrm{}j}}{\overset{k}{}}\frac{1}{1\xi ^\mathrm{}j}=(1)^k\xi ^{k(k+1)/2}\underset{\genfrac{}{}{0pt}{}{\mathrm{}=0}{\mathrm{}j}}{\overset{k}{}}\frac{1}{\xi ^j\xi ^{\mathrm{}}}$$
(143)
since
$$\underset{m=0}{\overset{k2}{}}\frac{x_{n+j}+m}{x_n+m}\left(\frac{x_n}{x_{n+j}}\right)^k\underset{\genfrac{}{}{0pt}{}{\mathrm{}=0}{\mathrm{}j}}{\overset{k}{}}\frac{1}{1x_{n+\mathrm{}}/x_{n+j}}\xi ^{(k1)j}\xi ^{k(j)}\underset{\genfrac{}{}{0pt}{}{\mathrm{}=0}{\mathrm{}j}}{\overset{k}{}}\frac{\xi ^l}{\xi ^{\mathrm{}}\xi ^j}$$
(144)
for $`n\mathrm{}`$. Thus, the limiting transformation is given by
$$\stackrel{}{}{}_{}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\},\xi )=\frac{{\displaystyle \underset{j=0}{\overset{k}{}}}{\displaystyle \frac{s_{n+j}}{\omega _{n+j}}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\mathrm{}=0}{\mathrm{}j}}{\overset{k}{}}}{\displaystyle \frac{1}{\xi ^j\xi ^{\mathrm{}}}}}{{\displaystyle \underset{j=0}{\overset{k}{}}}{\displaystyle \frac{1}{\omega _{n+j}}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\mathrm{}=0}{\mathrm{}j}}{\overset{k}{}}}{\displaystyle \frac{1}{\xi ^j\xi ^{\mathrm{}}}}}.$$
(145)
Comparison with the definition (39) of the divided difference operators reveals that the limiting transformation can be rewritten as
$$\stackrel{}{}{}_{}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\},\xi )=\frac{\mathrm{}{}_{n}{}^{(k)}[\{\{\xi ^n\}\}](s_n/\omega _n)}{\mathrm{}{}_{n}{}^{(k)}[\{\{\xi ^n\}\}](1/\omega _n)}.$$
(146)
Comparison to Eq. (133) shows that the limiting transformation is nothing but the W algorithm for $`t_n=\xi ^n`$. As characteristic polynomial we obtain
$$\stackrel{}{\Pi }{}_{}{}^{(k)}(z)=\underset{j=0}{\overset{k}{}}z^j\underset{\genfrac{}{}{0pt}{}{\mathrm{}=0}{\mathrm{}j}}{\overset{k}{}}\frac{1}{\xi ^j\xi ^{\mathrm{}}}=\xi ^{k(k+1)/2}\underset{j=0}{\overset{k1}{}}\frac{1z\xi ^j}{\xi ^{j+1}1}.$$
(147)
The last equality is easily proved by induction. Hence, the $`\stackrel{}{}`$ transformation is convex since $`\stackrel{}{\Pi }{}_{}{}^{(k)}(1)=0`$.
As shown in Appendix B, the $``$ transformation may be computed using the recursive scheme
$$\begin{array}{cc}& N{}_{n}{}^{(0)}=\frac{1}{x_n1}\frac{s_n}{\omega _n},D{}_{n}{}^{(0)}=\frac{1}{x_n1}\frac{1}{\omega _n},\hfill \\ & N{}_{n}{}^{(k)}=\frac{(x_{n+k}+k2)N{}_{n+1}{}^{(k1)}(x_n+k2)N_n^{(k1)}}{x_{n+k}x_n},\hfill \\ & D{}_{n}{}^{(k)}=\frac{(x_{n+k}+k2)D{}_{n+1}{}^{(k1)}(x_n+k2)D_n^{(k1)}}{x_{n+k}x_n},\hfill \\ & {}_{n}{}^{(k)}=N{}_{n}{}^{(k)}/D{}_{n}{}^{(k)}.\hfill \end{array}$$
(148)
It follows directly from Eq. (146) and the recursion relation for divided differences that the limiting transformation can be computed via the recursive scheme
$$\begin{array}{cc}& \stackrel{}{๐}{}_{n}{}^{(0)}=\frac{s_n}{\omega _n},\stackrel{}{๐ท}{}_{n}{}^{(0)}=\frac{1}{\omega _n},\hfill \\ & \stackrel{}{๐}{}_{n}{}^{(k)}=\frac{\stackrel{}{๐}{}_{n+1}{}^{(k1)}\stackrel{}{๐}_n^{(k1)}}{\xi ^{(n+k)}\xi ^n},\hfill \\ & \stackrel{}{๐ท}{}_{n}{}^{(k)}=\frac{\stackrel{}{๐ท}{}_{n+1}{}^{(k1)}\stackrel{}{๐ท}_n^{(k1)}}{\xi ^{(n+k)}\xi ^n},\hfill \\ & \stackrel{}{}{}_{n}{}^{(k)}=\stackrel{}{๐}{}_{n}{}^{(k)}/\stackrel{}{๐ท}{}_{n}{}^{(k)}.\hfill \end{array}$$
(149)
#### 4.2.8 $`๐ฅ๐`$ Transformation
This transformation is newly introduced in this article. In Section 5.2.1, it is derived via (asymptotically) hierarchically consistent iteration of the $`๐^{(2)}`$ transformation, i.e., of
$$s_n^{}=\frac{\mathrm{}^2(s_n/\omega _n)}{\mathrm{}^2(1/\omega _n)}.$$
(150)
The $`๐ฅ๐`$ transformation may be defined via the recursive scheme
$$\begin{array}{cc}& N{}_{n}{}^{(0)}=s_n/\omega _n,D{}_{n}{}^{(0)}=1/\omega _n,\hfill \\ & N{}_{n}{}^{(k)}=\stackrel{~}{}{}_{n}{}^{(k1)}N{}_{n}{}^{(k1)},D{}_{n}{}^{(k)}=\stackrel{~}{}{}_{n}{}^{(k1)}D{}_{n}{}^{(k1)},\hfill \\ & ๐ฅ๐{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\},\{\zeta {}_{n}{}^{(k)}\})=N{}_{n}{}^{(k)}/D{}_{n}{}^{(k)},\hfill \end{array}$$
(151)
where the generalized difference operator defined in Eq. (34) involves quantities $`\zeta {}_{n}{}^{(k)}0`$ for $`k_0`$. Special cases of the $`๐ฅ๐`$ transformation result from corresponding choices of the $`\zeta _n^{(k)}`$. From Eq. (151) one easily obtains the alternative representation
$$๐ฅ๐{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\},\{\zeta {}_{n}{}^{(k)}\})=\frac{\stackrel{~}{}{}_{n}{}^{(k1)}\stackrel{~}{}{}_{n}{}^{(k2)}\mathrm{}\stackrel{~}{}{}_{n}{}^{(0)}[s_n/\omega _n]}{\stackrel{~}{}{}_{n}{}^{(k1)}\stackrel{~}{}{}_{n}{}^{(k2)}\mathrm{}\stackrel{~}{}{}_{n}{}^{(0)}[1/\omega _n]}.$$
(152)
Thus, the $`๐ฅ๐^{(k)}`$ is a Levin-type sequence transformation of order $`2k`$.
#### 4.2.9 $``$ Transformation and Generalized $``$ Transformation
The $``$ transformation was introduced by Homeier and used or studied in a series of articles . Target of the $``$ transformation are Fourier series
$$s=A_0/2+\underset{j=1}{\overset{\mathrm{}}{}}\left(A_j\mathrm{cos}(j\alpha )+B_j\mathrm{sin}(j\alpha )\right)$$
(153)
with partial sums $`s_n=A_0/2+_{j=1}^n\left(A_j\mathrm{cos}(j\alpha )+B_j\mathrm{sin}(j\alpha )\right)`$ where the Fourier coefficients $`A_n`$ and $`B_n`$ have asymptotic expansions of the form
$$C_n\rho ^nn^ฯต\underset{j=0}{\overset{\mathrm{}}{}}c_jn^j$$
(154)
for $`n\mathrm{}`$ with $`\rho ๐`$, $`ฯต๐`$ and $`c_00`$.
The $``$ transformation was critized by Sidi as very unstable and useless near singularities of the Fourier series. However, Sidi failed to notice that โ as in the case of the $`d^{(1)}`$ transformation with $`\xi _n=\tau n`$ โ one can apply also the $``$ transformation (and also most other Levin-type sequence transformations) to the subsequence $`\{\{s_{\xi _n}\}\}`$ of $`\{\{s_n\}\}`$. The new sequence elements $`s_{\xi _n}=s_{\tau n}`$ can be regarded as the partial sums of a Fourier series with $`\tau `$โfold frequency. Using this $`\tau `$-fold frequency approach, one can obtain stable and accurate convergence acceleration even in the vicinity of singularities .
The $``$ transformation may be defined as
$$\begin{array}{cc}N_n^{(0)}\hfill & =(n+\beta )^1s_n/\omega _n,D{}_{n}{}^{(0)}=(n+\beta )^1/\omega _n,\hfill \\ N_n^{(k)}\hfill & =(n+\beta )N{}_{n}{}^{(k1)}+(n+2k+\beta )N_{n+2}^{(k1)}\hfill \\ & 2\mathrm{cos}(\alpha )(n+k+\beta )N{}_{n+1}{}^{(k1)},\hfill \\ D_n^{(k)}\hfill & =(n+\beta )D{}_{n}{}^{(k1)}+(n+2k+\beta )D_{n+2}^{(k1)}\hfill \\ & 2\mathrm{cos}(\alpha )(n+k+\beta )D{}_{n+1}{}^{(k1)},\hfill \\ \multicolumn{2}{c}{{}_{n}{}^{(k)}(\alpha ,\beta ,\{\{s_n\}\},\{\{\omega _n\}\})=N{}_{n}{}^{(k)}/D{}_{n}{}^{(k)},}\end{array}$$
(155)
where $`\mathrm{cos}\alpha \pm 1`$ and $`\beta _+`$.
It can also be represented in the explicit form
$${}_{n}{}^{(k)}(\alpha ,\beta ,\{\{s_n\}\},\{\{\omega _n\}\})=\frac{๐ซ[P^{(2k)}(\alpha )][(n+\beta )^{k1}s_n/\omega _n]}{๐ซ[P^{(2k)}(\alpha )][(n+\beta )^{k1}/\omega _n]}$$
(156)
where the $`p{}_{m}{}^{(2k)}(\alpha )`$ and the polynomial $`P^{(2k)}(\alpha )^{(2k)}`$ are defined via
$$P^{(2k)}(\alpha )(x)=(x^22x\mathrm{cos}\alpha +1)^k=\underset{m=0}{\overset{2k}{}}p{}_{m}{}^{(2k)}(\alpha )x^m$$
(157)
and $`๐ซ`$ is the polynomial operator defined in Eq. (38). This shows that the $`^{(k)}`$ transformation is a Levin-type transformation of order $`2k`$. It is not convex.
A subnormalized form is
$${}_{n}{}^{(k)}(\alpha ,\beta ,\{\{s_n\}\},\{\{\omega _n\}\})=\frac{{\displaystyle \underset{m=0}{\overset{2k}{}}}p{}_{m}{}^{(2k)}(\alpha ){\displaystyle \frac{(n+\beta +m)^{k1}}{(n+\beta +2k)^{k1}}}{\displaystyle \frac{s_{n+m}}{\omega _{n+m}}}}{{\displaystyle \underset{m=0}{\overset{2k}{}}}p{}_{m}{}^{(2k)}(\alpha ){\displaystyle \frac{(n+\beta +m)^{k1}}{(n+\beta +2k)^{k1}}}{\displaystyle \frac{1}{\omega _{n+m}}}}.$$
(158)
This relation shows that the limiting transformation
$$\stackrel{}{}{}_{}{}^{(k)}=\frac{๐ซ[P^{(2k)}(\alpha )][s_n/\omega _n]}{๐ซ[P^{(2k)}(\alpha )][1/\omega _n]}$$
(159)
exists, and has characteristic polynomial $`P^{(2k)}(\alpha )`$.
A generalized $``$ transformation was defined by Homeier . It is given in terms of the polynomial $`P^{(k,M)}(๐)^{(kM)}`$ with
$$P^{(k,M)}(๐)(x)=\underset{m=1}{\overset{M}{}}(xe_m)^k=\underset{\mathrm{}=0}{\overset{kM}{}}p{}_{\mathrm{}}{}^{(k,M)}(๐)x^{\mathrm{}}$$
(160)
where $`๐=(e_1,\mathrm{},e_M)๐^M`$ is a vector of constant parameters. Then, the generalized $``$ transformation is defined as
$${}_{n}{}^{(k,M)}(\beta ,\{\{s_n\}\},\{\{\omega _n\}\},๐)=\frac{๐ซ[P^{(k,M)}(๐)][(n+\beta )^{k1}s_n/\omega _n]}{๐ซ[P^{(k,M)}(๐)][(n+\beta )^{k1}/\omega _n]}$$
(161)
This shows that the generalized $`^{(k,M)}`$ is a Levin-type sequence transformation of order $`kM`$. The generalized $``$ transformation can be computed recursively using the scheme
$$\begin{array}{cc}\hfill N_n^{(0)}& =(n+\beta )^1s_n/\omega _n,D{}_{n}{}^{(0)}=(n+\beta )^1/\omega _n,\hfill \\ \hfill N_n^{(k)}& =\underset{j=0}{\overset{M}{}}q_j(n+\beta +jk)N{}_{n+j}{}^{(k1)},\hfill \\ \hfill D_n^{(k)}& =\underset{j=0}{\overset{M}{}}q_j(n+\beta +jk)D{}_{n+j}{}^{(k1)},\hfill \\ \multicolumn{2}{c}{{}_{n}{}^{(k,M)}(\beta ,\{\{s_n\}\},\{\{\omega _n\}\},๐)=\frac{N_n^{(k)}}{D_n^{(k)}}.}\end{array}$$
(162)
Here, the $`q_j`$ are defined by
$$\underset{m=1}{\overset{M}{}}(xe_m)=\underset{j=0}{\overset{M}{}}q_jx^j.$$
(163)
The algorithm (155) is a special case of the algorithm (162). To see this, one observes that $`M=2`$, $`e_1=\mathrm{exp}(\text{i}\alpha )`$ und $`e_2=\mathrm{exp}(\text{i}\alpha )`$ imply $`q_0=q_2=1`$ and $`q_1=2\mathrm{cos}(\alpha )`$.
For $`M=1`$ and $`e_1=1`$, the Levin transformation is recovered.
#### 4.2.10 $``$ Transformation
The $``$ transformation was in a slightly different form introduced by Homeier . It was derived via (asymptotically) hierarchically consistent iteration of the $`^{(1)}`$ transformation, i.e., of
$$s_n^{}=\frac{s_{n+2}/\omega _{n+2}2\mathrm{cos}(\alpha )s_{n+1}/\omega _{n+1}+s_n/\omega _n}{1/\omega _{n+2}2\mathrm{cos}(\alpha )/\omega _{n+1}+1/\omega _n}$$
(164)
For the derivation and an analysis of the properties of the $``$ transformation see . The $``$ transformation may be defined via the recursive scheme
$$\begin{array}{cc}N_n^{(0)}\hfill & =s_n/\omega _n,D{}_{n}{}^{(0)}=1/\omega _n,\hfill \\ N_n^{(k+1)}\hfill & ={}_{n}{}^{(k)}[\alpha ]N{}_{n}{}^{(k)},\hfill \\ D_n^{(k+1)}\hfill & ={}_{n}{}^{(k)}[\alpha ]D{}_{n}{}^{(k)},\hfill \\ \multicolumn{2}{c}{{}_{n}{}^{(k)}(\alpha ,\{\{s_n\}\},\{\{\omega _n\}\},\{\mathrm{\Delta }{}_{n}{}^{(k)}\})=\frac{N_n^{(k)}}{D_n^{(k)}}.}\end{array}$$
(165)
where the generalized difference operator $`{}_{n}{}^{(k)}[\alpha ]`$ defined in Eq. (35) involves quantities $`\mathrm{\Delta }{}_{n}{}^{(k)}0`$ for $`k_0`$. Special cases of the $``$ transformation result from corresponding choices of the $`\mathrm{\Delta }_n^{(k)}`$. From Eq. (165) one easily obtains the alternative representation
$${}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\},\{\mathrm{\Delta }{}_{n}{}^{(k)}\})=\frac{{}_{n}{}^{(k1)}[\alpha ]{}_{n}{}^{(k2)}[\alpha ]\mathrm{}{}_{n}{}^{(0)}[\alpha ][s_n/\omega _n]}{{}_{n}{}^{(k1)}[\alpha ]{}_{n}{}^{(k2)}[\alpha ]\mathrm{}{}_{n}{}^{(0)}[\alpha ][1/\omega _n]}.$$
(166)
Thus, $`^{(k)}`$ is a Levin-type sequence transformation of order $`2k`$. It is not convex.
Put $`\mathrm{\Theta }_n^{(0)}=1`$ and for $`k>0`$ define
$$\mathrm{\Theta }_n^{(k)}=\frac{\mathrm{\Delta }_n^{(0)}\mathrm{}\mathrm{\Delta }_n^{(k1)}}{\mathrm{\Delta }_{n+1}^{(0)}\mathrm{}\mathrm{\Delta }_{n+1}^{(k1)}}.$$
(167)
If for all $`k`$ the limits
$$\underset{n\mathrm{}}{lim}\mathrm{\Theta }_n^{(k)}=\mathrm{\Theta }_k$$
(168)
exist (we have always $`\mathrm{\Theta }_0=1`$), then one can define a limiting transformation $`\stackrel{}{}`$ for large $`n`$. It is a special case of the $``$ transformation according to
$$\stackrel{}{}{}_{n}{}^{(k)}(\alpha ,\{\{s_n\}\},\{\{\omega _n\}\},\{\{\mathrm{\Theta }_k\}\})={}_{n}{}^{(k)}(\alpha ,\{\{s_n\}\},\{\{\omega _n\}\},\{\{(\mathrm{\Theta }_k/\mathrm{\Theta }_{k+1})^n\}\}).$$
(169)
This is a transformation of order $`2k`$. The characteristic polynomials of $`\stackrel{}{}`$ are known to be
$$Q^{(2k)}(\alpha )^{2k}:Q^{(2k)}(\alpha )(z)=\underset{j=0}{\overset{k1}{}}\left[(1z\mathrm{\Theta }_j\mathrm{exp}(i\alpha ))(1z\mathrm{\Theta }_j\mathrm{exp}(i\alpha ))\right].$$
(170)
#### 4.2.11 $`๐ฆ`$ Transformation
The $`๐ฆ`$ transformation was introduced by Homeier in a slightly different form. It was obtained via iteration of the simple transformation
$$s_n^{}=\frac{\zeta _n^{(0)}{\displaystyle \frac{s_n}{\omega _n}}+\zeta _n^{(1)}{\displaystyle \frac{s_{n+1}}{\omega _{n+1}}}+\zeta _n^{(2)}{\displaystyle \frac{s_{n+2}}{\omega _{n+2}}}}{\zeta _n^{(0)}{\displaystyle \frac{1}{\omega _n}}+\zeta _n^{(1)}{\displaystyle \frac{1}{\omega _{n+1}}}+\zeta _n^{(2)}{\displaystyle \frac{1}{\omega _{n+2}}}}$$
(171)
that is exact for sequences of the form
$$s_n=s+\omega _n(cP_n+dQ_n),$$
(172)
where $`c`$ and $`d`$ are arbitrary constants, while $`P_n`$ and $`Q_n`$ are two linearly independent solutions of the three-term recurrence
$$\zeta _n^{(0)}v_n+\zeta _n^{(1)}v_{n+1}+\zeta _n^{(2)}v_{n+2}=0.$$
(173)
The $`๐ฆ`$ transformation may be defined via the recursive scheme
$$\begin{array}{cc}N_n^{(0)}\hfill & =s_n/\omega _n,D{}_{n}{}^{(0)}=1/\omega _n,\hfill \\ N_n^{(k+1)}\hfill & ={}_{n}{}^{(k)}[\zeta ]N{}_{n}{}^{(k)},\hfill \\ D_n^{(k+1)}\hfill & ={}_{n}{}^{(k)}[\zeta ]D{}_{n}{}^{(k)},\hfill \\ \multicolumn{2}{c}{๐ฆ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\},\{\stackrel{~}{\mathrm{\Delta }}{}_{n}{}^{(k)}\},\{\zeta {}_{n}{}^{(j)}\})=\frac{N_n^{(k)}}{D_n^{(k)}}.}\end{array}$$
(174)
where the generalized difference operator $`{}_{n}{}^{(k)}[\zeta ]`$ defined in Eq. (36) involves recursion coefficients $`\zeta _{n+k}^{(j)}`$ with $`j=0,1,2`$ and quantities $`\stackrel{~}{\mathrm{\Delta }}{}_{n}{}^{(k)}0`$ for $`k_0`$. Special cases of the $`๐ฆ`$ transformation for given recursion, i.e., for given $`\zeta _n^{(j)}`$, result from corresponding choices of the $`\stackrel{~}{\mathrm{\Delta }}_n^{(k)}`$. From Eq. (174) one easily obtains the alternative representation
$$๐ฆ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\},\{\stackrel{~}{\mathrm{\Delta }}{}_{n}{}^{(k)}\},\{\zeta {}_{n}{}^{(j)}\})=\frac{{}_{n}{}^{(k1)}[\zeta ]{}_{n}{}^{(k2)}[\zeta ]\mathrm{}{}_{n}{}^{(0)}[\zeta ][s_n/\omega _n]}{{}_{n}{}^{(k1)}[\zeta ]{}_{n}{}^{(k2)}[\zeta ]\mathrm{}{}_{n}{}^{(0)}[\zeta ][1/\omega _n]}.$$
(175)
Thus, $`๐ฆ^{(k)}`$ is a Levin-type sequence transformation of order $`2k`$. It is not convex.
For applications of the $`๐ฆ`$ transformation see .
## 5 Methods for the Construction of Levin-Type Transformations
In this section, we discuss approaches for the construction of Levin-type sequence transformations and point out the relation to their kernel.
### 5.1 Model Sequences and Annihilation Operators
As discussed in the introduction, the derivation of sequence transformations may be based on model sequences. These may be of the form (10) or of the form (6). Here, we consider model sequences of the latter type that involves remainder estimates $`\omega _n`$. As described in Section 3.1, determinantal representations for the corresponding sequence transformations can be derived using Cramerโs rule, and one of the recursive schemes of the $`๐`$ algorithm may be used for the computation. However, for important special choices of the functions $`\psi _j(n)`$, simpler recursive schemes and more explicit representations in the form (11) can be obtained using the annihilation operator approach of Weniger . This approach was also studied by Brezinski and Matos who showed that it leads to a unified derivation of many extrapolation algorithms and related devices and general results about their kernels. Further, we mention the work of Matos who analysed the approach further and derived a number of convergence acceleration results for Levin-type sequence transformations.
In this approach, an annihilation operator $`๐=๐_n^{(k)}`$ as defined in Eq. (31) is needed that annihilates the sequences $`\{\{\psi _j(n)\}\}`$, i.e., such that
$$๐{}_{n}{}^{(k)}(\{\{\psi _j(n)\}\})=0\text{for }j=0,\mathrm{},k1.$$
(176)
Rewriting Eq. (6) in the form
$$\frac{\sigma _n\sigma }{\omega _n}=\underset{j=0}{\overset{k1}{}}c_j\psi _j(n),$$
(177)
and applying $`๐`$ to both sides of this equation, one sees that
$$๐{}_{n}{}^{(k)}\{\{\frac{\sigma _n\sigma }{\omega _n}\}\}=0$$
(178)
This equation may be solved for $`\sigma `$ due to the linearity of $`๐`$. The result is
$$\sigma =\frac{๐{}_{n}{}^{(k)}(\{\{\sigma _n/\omega _n\}\})}{๐{}_{n}{}^{(k)}(\{\{1/\omega _n\}\})}$$
(179)
leading to a sequence transformation
$$๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega \}\})=\frac{๐{}_{n}{}^{(k)}\left(\{\{s_n/\omega _n\}\}\right)}{๐{}_{n}{}^{(k)}\left(\{\{1/\omega _n\}\}\right)}$$
(180)
Since $`๐`$ is linear, this transformation can be rewritten in the form (11), i.e., a Levin-type transformation has been obtained.
We note that this process can be reversed, that is, for each Levin-type sequence transformation $`๐ฏ[\mathrm{\Lambda }^{(k)}]`$ of order $`k`$ there is an annihilation operator, namely the polynomial operator $`๐ซ[\mathrm{\Pi }{}_{n}{}^{(k)}]`$ as defined in Eq. (38) where $`\mathrm{\Pi }_n^{(k)}`$ are the characteristic polynomials as defined in Eq. (72). Using this operator, the defining equation (66) can be rewritten as
$$๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\})=\frac{๐ซ[\mathrm{\Pi }{}_{n}{}^{(k)}](s_n/\omega _n)}{๐ซ[\mathrm{\Pi }{}_{n}{}^{(k)}](1/\omega _n)}$$
(181)
Let $`\varphi _{n,m}^{(k)}`$ for $`m=0,\mathrm{},k1`$ be $`k`$ linearly independent solutions of the linear $`(k+1)`$โterm recurrence
$$\underset{j=0}{\overset{k}{}}\lambda {}_{n,j}{}^{(k)}v_{n+j}^{}=0.$$
(182)
Then $`๐ซ[\mathrm{\Pi }{}_{n}{}^{(k)}]\varphi {}_{n,m}{}^{(k)}=0`$ for $`m=0,\mathrm{},k1`$, i.e., $`๐ซ[\mathrm{\Pi }{}_{n}{}^{(k)}]`$ is an annihilation operator for all solutions of Eq. (182). Thus, all sequences that are annihilated by this operator are linear combinations of the $`k`$ sequences $`\{\{\varphi {}_{n,m}{}^{(k)}\}\}`$.
If $`\{\{\sigma _n\}\}`$ is a sequence in the kernel of $`๐ฏ^{(k)}`$ with (anti)limit $`\sigma `$, we must have
$$\sigma =\frac{๐ซ[\mathrm{\Pi }{}_{n}{}^{(k)}](\sigma _n/\omega _n)}{๐ซ[\mathrm{\Pi }{}_{n}{}^{(k)}](1/\omega _n)}$$
(183)
or after some rearrangement using the linearity of $`๐ซ`$
$$๐ซ[\mathrm{\Pi }{}_{n}{}^{(k)}]\left(\frac{\sigma _n\sigma }{\omega _n}\right)=0.$$
(184)
Hence, we must have
$$\frac{\sigma _n\sigma }{\omega _n}=\underset{m=0}{\overset{k1}{}}c_m\varphi {}_{n,m}{}^{(k)},$$
(185)
or, equivalently
$$\sigma _n=\sigma +\omega _n\underset{m=0}{\overset{k1}{}}c_m\varphi {}_{n,m}{}^{(k)},$$
(186)
for some constants $`c_m`$. Thus, we have determined the kernel of $`๐ฏ^{(k)}`$ that can also be considered as the set of model sequences for this transformation. Thus, we have proved the following theorem:
###### Theorem 1
Let $`\varphi _{n,m}^{(k)}`$ for $`m=0,\mathrm{},k1`$ be the $`k`$ linearly independent solutions of the linear $`(k+1)`$โterm recurrence (182). The kernel of $`๐ฏ[\mathrm{\Lambda }^{(k)}](\{\{s_n\}\},\{\{\omega _n\}\})`$ is given by all sequences $`\{\{\sigma _n\}\}`$ with (anti)limit $`\sigma `$ and elements $`\sigma _n`$ of the form (186) for arbitrary constants $`c_m`$.
We note that the $`\psi _j(n)`$ for $`j=0,\mathrm{},k1`$ can essentially be identified with the $`\varphi _{n,j}^{(k)}`$. Thus, we have determinantal representations for known $`\psi _j(n)`$ as noted above in the context of the E algorithm. See also for determinantal representations of the $`๐ฅ`$ transformations and the relation to its kernel.
Examples of annihilation operators and the functions $`\psi _j(n)`$ that are annihilated are given in Table 2. Examples for the Levin-type sequence transformations that have been derived using the approach of model sequences are discussed in Section 5.1.2.
Note that the annihilation operators used by Weniger were weighted difference operators $`๐ฒ_n^{(k)}`$ as defined in Eq. (37). Homeier discussed operator representations for the $`๐ฅ`$ transformation that are equivalent to many of the annihilation operators and related sequence transformations as given by Brezinski and Matos . The latter have been further discussed by Matos who considered among others Levin-type sequence transformations with constant coefficients, $`\lambda {}_{n,j}{}^{(k)}=const.`$, and with polynomial coefficients $`\lambda {}_{n,j}{}^{(k)}=\lambda _j(n+1)`$, with $`\lambda _j`$, and $`n_0`$, in particular annihilation operators of the form
$$L(u_n)=(\mathrm{\Omega }^l+\lambda _1\mathrm{\Omega }^{l1}+\mathrm{}+\lambda _l)(u_n)$$
(187)
with the special cases
$$L_1(u_n)=(\mathrm{\Omega }\alpha _1)(\mathrm{\Omega }\alpha _2)\mathrm{}(\mathrm{\Omega }\alpha _l)(u_n),(\alpha _i\alpha _j\text{ for all }ij,$$
(188)
and
$$L_2(u_n)=(\mathrm{\Omega }\alpha )^l(u_n)$$
(189)
where
$$\mathrm{\Omega }^r(u_n)=(n+1)_ru_{n+r},n_0$$
(190)
and
$$\stackrel{~}{L}(u_n)=(\pi \alpha _1)(\pi \alpha _2)\mathrm{}(\pi \alpha _l)(u_n)$$
(191)
where
$$\pi (u_n)=(n+1)\mathrm{}u_n,\pi ^r(u_n)=\pi (\pi ^{r1}(u_n)),n_0$$
(192)
and the $`\lambda `$โs and $`\alpha `$โs are constants. Note that $`n`$ is shifted in comparison to where the convention $`n`$ was used. See also Table 2 for the corresponding annihilated functions $`\psi _j(n)`$.
Matos also considered difference operators of the form
$$L(u_n)=\mathrm{}^k+p_{k1}(n)\mathrm{}^{k1}+\mathrm{}+p_1(n)\mathrm{}+p_0(n)$$
(193)
where the functions $`f_j`$ given by $`f_j(t)=p_j(1/t)t^{k+j}`$ for $`j=0,\mathrm{},k1`$ are analytic in the neighborhood of 0. For such operators, there is no explicit formula for the functions that are annihilated. However, the asymptotic behavior of such functions is known . We will later return to such annihilation operators and state some convergence results.
#### 5.1.1 Derivation of the $``$ Transformation
As an example for the application of the annihilation operator approach, we derive the $``$ transformation. Consider the model sequence
$$\sigma _n=\sigma +\omega _n\underset{j=0}{\overset{k1}{}}c_j\frac{1}{(x_n)_j}$$
(194)
that may be rewritten as
$$\frac{\sigma _n\sigma }{\omega _n}=\underset{j=0}{\overset{k1}{}}c_j\frac{1}{(x_n)_j}.$$
(195)
We note that Eq. (194) corresponds to modeling $`\mu _n=R_n/\omega _n`$ as a truncated factorial series in $`x_n`$ (instead as a truncated power series as in the case of the W algorithm). The $`x_n`$ are elements of $`\{\{x_n\}\}`$ an auxiliary sequence $`\{\{x_n\}\}`$ such that $`lim_n\mathrm{}1/x_n=0`$ and also $`x_{n+\mathrm{}}>x_n`$ for $`\mathrm{}`$ and $`x_0>1`$, i.e., a diverging sequence of monotonously increasing positive numbers. To find an annihilation operator for the $`\psi _j(n)=1/(x_n)_j`$, we make use of the fact that the divided difference operator $`\mathrm{}{}_{n}{}^{(k)}=\mathrm{}{}_{n}{}^{(k)}[\{\{x_n\}\}]`$ annihilates polynomials in $`x_n`$ of degree less than $`k`$. Also, we observe that the definition of the Pochhammer symbols entails that
$$(x_n)_{k1}/(x_n)_j=(x_n+j)_{k1j},$$
(196)
is a polynomial of degree less than $`k`$ in $`x_n`$ for $`0jk1`$. Thus, the sought annihilation operator is $`๐=\mathrm{}{}_{n}{}^{(k)}(x_n)_{k1}^{}`$ because
$$\mathrm{}{}_{n}{}^{(k)}(x_n)_{k1}^{}\frac{1}{(x_n)_j}=0,0jk1.$$
(197)
Hence, for the model sequence (194), one can calculate $`\sigma `$ via
$$\sigma =\frac{\mathrm{}{}_{n}{}^{(k)}((x_n)_{k1}\sigma _n/\omega _n)}{\mathrm{}{}_{n}{}^{(k)}((x_n)_{k1}/\omega _n)}$$
(198)
and the $``$ transformation (140) results by replacing $`\sigma _n`$ by $`s_n`$ in the right hand side of Eq. (198).
#### 5.1.2 Important Special Cases
Here, we collect model sequences and annihilation operators for some important Levin-type sequence transformations that were derived using the model sequence approach. For further examples see also . The model sequences are the kernels by construction. In Section 5.2.2, kernels and annihilation operators are stated for important Levin-type transformation that were derived using iterative methods.
##### Levin transformation
The model sequence for $`^{(k)}`$ is
$$\sigma _n=\sigma +\omega _n\underset{j=0}{\overset{k1}{}}c_j/(n+\beta )^j.$$
(199)
The annihilation operator is
$$๐{}_{n}{}^{(k)}=\mathrm{}^k(n+\beta )^{k1}.$$
(200)
##### Weniger transformations
The model sequence for $`๐ฎ^{(k)}`$ is
$$\sigma _n=\sigma +\omega _n\underset{j=0}{\overset{k1}{}}c_j/(n+\beta )_j.$$
(201)
The annihilation operator is
$$๐{}_{n}{}^{(k)}=\mathrm{}^k(n+\beta )_{k1}.$$
(202)
The model sequence for $`^{(k)}`$ is
$$\sigma _n=\sigma +\omega _n\underset{j=0}{\overset{k1}{}}c_j/(n\xi )_j.$$
(203)
The annihilation operator is
$$๐{}_{n}{}^{(k)}=\mathrm{}^k(n\xi )_{k1}.$$
(204)
The model sequence for $`๐^{(k)}`$ is
$$\sigma _n=\sigma +\omega _n\underset{j=0}{\overset{k1}{}}c_j/(\alpha [n+\zeta ])_j.$$
(205)
The annihilation operator is
$$๐{}_{n}{}^{(k)}=\mathrm{}^k(\alpha [n+\zeta ])_{k1}.$$
(206)
##### W algorithm
The model sequence for $`W^{(k)}`$ is
$$\sigma _n=\sigma +\omega _n\underset{j=0}{\overset{k1}{}}c_jt_n^j.$$
(207)
The annihilation operator is
$$๐{}_{n}{}^{(k)}=\mathrm{}{}_{n}{}^{(k)}[\{\{t_n\}\}].$$
(208)
##### $``$ Transformation
The model sequence for $`^{(k)}`$ is
$$\sigma _n=\sigma +\omega _n\left(\mathrm{exp}(i\alpha n)\underset{j=0}{\overset{k1}{}}c_j^+/(n+\beta )^j+\mathrm{exp}(i\alpha n)\underset{j=0}{\overset{k1}{}}c_j^{}/(n+\beta )^j\right).$$
(209)
The annihilation operator is
$$๐{}_{n}{}^{(k)}=๐ซ[P^{(2k)}(\alpha )](n+\beta )^{k1}.$$
(210)
##### Generalized $``$ Transformation
The model sequence for $`^{(k,M)}`$ is
$$\sigma _n=\sigma +\omega _n\underset{m=1}{\overset{M}{}}e_m^n\underset{j=0}{\overset{k1}{}}c_{m,j}(n+\beta )^j.$$
(211)
The annihilation operator is
$$๐{}_{n}{}^{(k)}=๐ซ[P^{(k,M)}(๐)](n+\beta )^{k1}.$$
(212)
### 5.2 Hierarchically Consistent Iteration
As alternative to the derivation of sequence transformations using model sequences and possibly annihilation operators, one may take some simple sequence transformation $`T`$ and iterate it $`k`$ times to obtain a transformation $`T^{(k)}=T\mathrm{}T`$. For the iterated transformation, by construction one has a simple algorithm by construction, but the theoretical analysis is complicated since usually no kernel is known. See for instance the iterated Aitken process where the $`\mathrm{\Delta }^2`$ method plays the rรดle of the simple transformation. However, as is discussed at length in Refs. , there are usually several possibilities for the iteration. Both problems โ unknown kernel and arbitrariness of iteration โ are overcome using the concept of hierarchical consistency that was shown to give rise to powerful algorithms like the $`๐ฅ`$ and the $``$ transformations . The basic idea of the concept is to provide a hierarchy of model sequences such that the simple transformation provides a mapping between neighboring levels of the hierarchy. To ensure the latter, normally one has to fix some parameters in the simple transformation to make the iteration consistent with the hierarchy.
A formal description of the concept is given in the following taken mainly from the literature. . As an example, the concept is later used to derive the $`๐ฅ๐`$ transformation in Section 5.2.1.
Let $`\{\{\sigma _n(\stackrel{}{c},\stackrel{}{p})\}\}_{n=0}^{\mathrm{}}`$ be a simple โbasicโ model sequence that depends on a vector $`\stackrel{}{c}๐^a`$ of constants, and further parameters $`\stackrel{}{p}`$. Assume that its (anti)limit $`\sigma (\stackrel{}{p})`$ exists and is independent of $`\stackrel{}{c}`$. Assume that the basic transformation $`T=T(\stackrel{}{p})`$ allows to compute the (anti)limit exactly according to
$$T(\stackrel{}{p}):\{\{\sigma _n(\stackrel{}{c},\stackrel{}{p})\}\}\{\{\sigma (\stackrel{}{p})\}\}.$$
(213)
Let the hierarchy of model sequences be given by
$$\{\{\{\sigma _n^{(\mathrm{})}(\stackrel{}{c}{}_{}{}^{(\mathrm{})},\stackrel{}{p}^{(\mathrm{})})|\stackrel{}{c}^{(\mathrm{})}๐^{a^{(\mathrm{})}}\}\}\}_{\mathrm{}=0}^L$$
(214)
with $`a^{(\mathrm{})}>a^{(\mathrm{}^{})}`$ for $`\mathrm{}>\mathrm{}^{}`$. Here, $`\mathrm{}`$ numbers the levels of the hierarchy. Each of the model sequences $`\{\{\sigma _n^{(\mathrm{})}(\stackrel{}{c}{}_{}{}^{(\mathrm{})},\stackrel{}{p}^{(\mathrm{})})\}\}`$ depends on an $`a^{(\mathrm{})}`$โdimensional complex vector $`\stackrel{}{c}^{(\mathrm{})}`$ and further parameters $`\stackrel{}{p}^{(\mathrm{})}`$. Assume that the model sequences of lower levels are also contained in those of higher levels: For all $`\mathrm{}<L`$ and all $`\mathrm{}^{}>\mathrm{}`$ and $`\mathrm{}^{}L`$, every sequence $`\{\{\sigma _n^{(\mathrm{})}(\stackrel{}{c}^{(\mathrm{})},\stackrel{}{p}{}_{}{}^{(\mathrm{})})\}\}`$ is assumed to be representable as a model sequence $`\{\{\sigma _n^{(\mathrm{}^{})}(\stackrel{}{c}{}_{}{}^{(\mathrm{}^{})},\stackrel{}{p}{}_{}{}^{(\mathrm{}^{})})\}\}`$ where $`\stackrel{}{c}^{(\mathrm{}^{})}`$ is obtained from $`\stackrel{}{c}^{(\mathrm{})}`$ by the natural injection $`๐^{a^{(\mathrm{})}}๐^{a^{(\mathrm{}^{})}}`$. Assume that for all $`\mathrm{}`$ with $`0<\mathrm{}L`$
$$T(\stackrel{}{p}{}_{}{}^{(\mathrm{})}):\{\{\sigma _n^{(\mathrm{})}(\stackrel{}{c}{}_{}{}^{(\mathrm{})},\stackrel{}{p}{}_{}{}^{(\mathrm{})})\}\}\{\{\sigma _n^{(\mathrm{}1)}(\stackrel{}{c}{}_{}{}^{(\mathrm{}1)},\stackrel{}{p}{}_{}{}^{(\mathrm{}1)})\}\}$$
(215)
is a mapping between neighboring levels of the hierarchy. Composition yields an iterative transformation
$$T^{(L)}=T(\stackrel{}{p}{}_{}{}^{(0)})T(\stackrel{}{p}{}_{}{}^{(1)})\mathrm{}T(\stackrel{}{p}{}_{}{}^{(L)})$$
(216)
This transformation is called โhierarchically consistentโ or โconsistent with the hierarchyโ. It maps model sequences $`\sigma _n^{(\mathrm{})}(\stackrel{}{c}{}_{}{}^{(\mathrm{})},\stackrel{}{p}^{(\mathrm{})})`$ to constant sequences if Eq. (213) holds with
$$\{\{\sigma _n^{(0)}(\stackrel{}{c}{}_{}{}^{(0)},\stackrel{}{p}{}_{}{}^{(0)})\}\}=\{\{\sigma _n(\stackrel{}{c},\stackrel{}{p})\}\}.$$
(217)
If instead of Eq. (215) we have
$$T(\stackrel{}{p}{}_{}{}^{(\mathrm{})})\left(\{\{\sigma _n^{(\mathrm{})}(\stackrel{}{c}{}_{}{}^{(\mathrm{})},\stackrel{}{p}^{(\mathrm{})})\}\}\right)\{\{\sigma _n^{(\mathrm{}1)}(\stackrel{}{c}{}_{}{}^{(\mathrm{}1)},\stackrel{}{p}{}_{}{}^{(\mathrm{}1)})\}\}$$
(218)
for $`n\mathrm{}`$ for all $`\mathrm{}>0`$ then the iterative transformation $`T^{(L)}`$ is called โasymptotically consistent with the hierarchyโ or โasymptotically hierarchy-consistentโ.
#### 5.2.1 Derivation of the $`๐ฅ๐`$ Transformation
The simple transformation is the $`๐^{(2)}`$ transformation
$$s_n^{}=T(\{\{\omega _n\}\})(\{\{s_n\}\})=\frac{\mathrm{}^2(s_n/\omega _n)}{\mathrm{}^2(1/\omega _n)}$$
(219)
depending on the โparametersโ $`\{\{\omega _n\}\}`$, with basic model sequences
$$\frac{\sigma _n}{\omega _n}=\sigma \frac{1}{\omega _n}+(an+b).$$
(220)
The more complicated model sequences of the next level are taken to be
$$\frac{\sigma _n}{\omega _n}=\sigma \frac{1}{\omega _n}+(an+b+(a_1n+b_1)r_n).$$
(221)
Application of $`\mathrm{}^2`$ eliminates the terms involving $`a`$ and $`b`$. The result is
$$\frac{\mathrm{}^2{\displaystyle \frac{\sigma _n}{\omega _n}}}{\mathrm{}^2r_n}=\sigma \frac{\mathrm{}^2{\displaystyle \frac{1}{\omega _n}}}{\mathrm{}^2r_n}++(a_1n+b_1+2a_1\frac{\mathrm{}r_n}{\mathrm{}^2r_n})$$
(222)
for $`\mathrm{}^2r_n0`$. Assuming that for large $`n`$
$$\frac{\mathrm{}r_n}{\mathrm{}^2r_n}=An+B+o(1)$$
(223)
holds, the result is asymptotically of the same *form* as the model sequence in Eq. (220), namely
$$\frac{\sigma _n^{}}{\omega _n^{}}=\sigma \frac{1}{\omega _n^{}}+(a^{}n+b^{}+o(1)).$$
(224)
with renormalized โparametersโ
$$1/\omega _n^{}=\frac{\mathrm{}^2(1/\omega _n)}{\mathrm{}^2r_n}$$
(225)
and obvious identifications for $`a^{}`$ and $`b^{}`$.
We now assume that this mapping between two neighboring levels of the hierarchy can be extended to any two neighboring levels, provided that one introduces $`\mathrm{}`$-dependent quantities, especially $`r_nr_n^{(\mathrm{})}`$ with $`\zeta {}_{n}{}^{(\mathrm{})}=\mathrm{}^2r{}_{n}{}^{(\mathrm{})}0`$, $`s_n/\omega _nN_n^{(\mathrm{})}`$, $`1/\omega _nD_n^{(\mathrm{})}`$ and $`s_n^{}/\omega _n^{}N_n^{(\mathrm{}+1)}`$, $`1/\omega _n^{}D_n^{(\mathrm{}+1)}`$.
Iterating in this way leads to the algorithm (151).
The condition (223) or more generally
$$\frac{\mathrm{}r_n^{(\mathrm{})}}{\mathrm{}^2r_n^{(\mathrm{})}}=A_{\mathrm{}}n+B_{\mathrm{}}+o(1)$$
(226)
for given $`\mathrm{}`$ and for large $`n`$ is satisfied in many cases. For instance, it is satisfied if there are constants $`\beta _{\mathrm{}}0`$, $`\gamma _{\mathrm{}}`$ and $`\delta _{\mathrm{}}0`$ such that
$$\mathrm{}r{}_{n}{}^{(\mathrm{})}\beta _{\mathrm{}}\{\begin{array}{cc}\left(\frac{\delta _{\mathrm{}}+1}{\delta _{\mathrm{}}}\right)^n\hfill & \text{for }\gamma _{\mathrm{}}=0\hfill \\ \left(\frac{\delta _{\mathrm{}}+1}{\gamma _{\mathrm{}}}\right)_n/\left(\frac{\delta _{\mathrm{}}}{\gamma _{\mathrm{}}}\right)_n\hfill & \text{otherwise}\hfill \end{array}.$$
(227)
This is for instance the case for $`r{}_{n}{}^{(\mathrm{})}=n^\zeta _{\mathrm{}}`$ with $`\zeta _{\mathrm{}}(\zeta _{\mathrm{}}1)0`$.
The kernel of $`๐ฅ๐^{(k)}`$ may be found inductively in the following way:
$$\begin{array}{ccc}& \hfill N{}_{n}{}^{(k)}\sigma D_n^{(k)}& =0\hfill \\ \hfill & \hfill \mathrm{}^2(N{}_{n}{}^{(k1)}\sigma D{}_{n}{}^{(k1)})& =0\hfill \\ \hfill & \hfill N{}_{n}{}^{(k1)}\sigma D_n^{(k1)}& =a_{k1}n+b_{k1}\hfill \\ \hfill & \hfill \mathrm{}^2(N{}_{n}{}^{(k2)}\sigma D{}_{n}{}^{(k2)})& =(a_{k1}n+b_{k1})\zeta _n^{(k2)}\hfill \\ \hfill & \hfill N{}_{n}{}^{(k2)}\sigma D{}_{n}{}^{(k2)})& =a_{k2}n+b_{k2}+\underset{j=0}{\overset{n2}{}}\underset{n^{}=0}{\overset{j}{}}(a_{k1}n^{}+b_{k1})\zeta _n^{}^{(k2)}\hfill \end{array}$$
(228)
yielding the result
$$\begin{array}{cc}\hfill N{}_{n}{}^{(0)}\sigma D{}_{n}{}^{(0)}=& a_0n+b_0+\underset{j=0}{\overset{n2}{}}\underset{n_1=0}{\overset{j}{}}\zeta {}_{n_1}{}^{(0)}(a_1n_1+b_1+\mathrm{}\hfill \\ & \mathrm{}(a_{k2}n+b_{k2}+\underset{j_{k2}=0}{\overset{n_{k2}2}{}}\underset{n_{k1}=0}{\overset{j_{k2}}{}}\zeta {}_{n_{k1}}{}^{(k2)}(a_{k1}n_{k1}+b_{k1})).\hfill \end{array}$$
(229)
Here, the definitions $`N{}_{n}{}^{(0)}=\sigma _n/\omega _n`$ and $`D{}_{n}{}^{(0)}=1/\omega _n`$ may be used to obtain the model sequence $`\{\{\sigma _n\}\}`$ for $`๐ฅ๐^{(k)}`$, that may be identified as kernel of that transformation, and also may be regarded as model sequence of the $`k`$-th level according to $`\{\{\sigma _n^{(k)}(\stackrel{}{c}^{(k)},\stackrel{}{p}{}_{}{}^{(k)})\}\}`$ with $`\stackrel{}{c}^{(k)}=(a_0,b_0,\mathrm{},a_{k1},b_{k1})`$ and $`\stackrel{}{p}^{(k)}`$ corresponds to $`\omega {}_{n}{}^{(k)}=1/D_n^{(k)}`$ and the $`\{\zeta {}_{n}{}^{(\kappa )}|0\kappa k2\}`$.
We note this as a theorem:
###### Theorem 2
The kernel of $`๐ฅ๐^{(k)}`$ is given by the set of sequences $`\{\{\sigma _n\}\}`$ such that Eq. (229) holds with $`N{}_{n}{}^{(0)}=\sigma _n/\omega _n`$ and $`D{}_{n}{}^{(0)}=1/\omega _n`$.
#### 5.2.2 Important Special Cases
Here, we give the hierarchies of model sequences for sequence transformations derived via hierarchically consistent iteration.
##### $`๐ฅ`$ transformation
The most prominent example is the $`๐ฅ`$ transformation (actually a large class of transformations). The corresponding hierarchy of model sequences provided by the kernels that are explicitly known according to the following theorem:
###### Theorem 3
The kernel of the $`๐ฅ^{(k)}`$ transformation is given by the sequences $`\{\{\sigma _n\}\}`$ with elements of the form
$$\sigma _n=\sigma +\omega _n\underset{j=0}{\overset{k1}{}}c_j\psi _j(n)$$
(230)
with
$$\begin{array}{cc}\hfill \psi _0(n)& =1\hfill \\ \hfill \psi _1(n)& =\underset{n_1=0}{\overset{n1}{}}\delta _{n_1}^{(0)}\hfill \\ \hfill \psi _2(n)& =\underset{n_1=0}{\overset{n1}{}}\delta {}_{n_1}{}^{(0)}\underset{n_2=0}{\overset{n_11}{}}\delta _{n_2}^{(1)}\hfill \\ & \mathrm{}\hfill \\ \hfill \psi _{k1}(n)& =\underset{n>n_1>n_2>\mathrm{}>n_{k1}}{}\delta {}_{n_1}{}^{(0)}\delta {}_{n_2}{}^{(1)}\mathrm{}\delta _{n_{k1}}^{(k2)}\hfill \end{array}$$
(231)
with arbitrary constants $`c_0,\mathrm{},c_{k1}`$.
##### $``$ transformation
Since the $``$ transformation is a special case of the $`๐ฅ`$ transformation (cp. Table LABEL:tabJtransSpecial and , its kernels (corresponding to the hierarchy of model sequences) are explicitly known according to the following theorem:
###### Theorem 4
\[44, Theorem 8\] The kernel of the $`^{(k)}`$ transformation is given by the sequences $`\{\{\sigma _n\}\}`$ with elements of the form
$$\begin{array}{cc}\hfill \sigma _n& =\sigma +\mathrm{exp}(i\alpha n)\omega _n[d_0+d_1\mathrm{exp}(2i\alpha n)\hfill \\ & +\underset{n_1=0}{\overset{n1}{}}\underset{n_2=0}{\overset{n_11}{}}\mathrm{exp}(2i\alpha (n_1n_2))\left(d_2+d_3\mathrm{exp}(2i\alpha n_2)\right)\mathrm{\Delta }_{n_2}^{(0)}+\mathrm{}\hfill \\ & +\underset{n>n_1>n_2>\mathrm{}>n_{2k2}}{}\mathrm{exp}(2i\alpha [n_1n_2+\mathrm{}+n_{2k3}n_{2k2}])\hfill \\ & \times (d_{2k2}+d_{2k1}\mathrm{exp}(2i\alpha n_{2k2}))\underset{j=0}{\overset{k2}{}}\mathrm{\Delta }_{n_{2j+2}}^{(j)}]\hfill \end{array}$$
(232)
with constants $`d_0,\mathrm{},d_{2k1}`$. Thus, we have $`s={}_{n}{}^{(k^{})}(\alpha ,\{\{s_n\}\},\{\{\omega _n\}\},\{\mathrm{\Delta }{}_{n}{}^{(k)}\})`$ for $`k^{}k`$ for sequences of this form.
### 5.3 A Two-Step Approach
In favorable cases, one may use a two-step approach for the construction of sequence transformations:
Use asymptotic analysis of the remainder $`R_n=s_ns`$ of the given problem to find the adequate model sequence (or hierarchy of model sequences) for large $`n`$.
Use the methods described in Sections 5.1 or 5.2 to construct the sequence transformation adapted to the problem.
This is, of course, a mathematically promising approach. A good example for the two-step approach is the derivation of the $`d^{(m)}`$ transformations by Levin and Sidi (compare also Section 3.2).
But there are two difficulties with this approach.
The first difficulty is a practical one. In many cases, the problems to be treated in applications are simply too complicated to allow to perform Step 1 of the two-step approach.
The second difficulty is a more mathematical one. The optimal system of functions $`f_j(n)`$ used in the asymptotic expansion,
$$s_ns\underset{j=0}{\overset{\mathrm{}}{}}c_jf_j(n)$$
(233)
with $`f_{j+1}(n)=o(f_j(n))`$, i.e., the optimal *asymptotic scale* \[102, p. 2\], is not clear *a priori*. For instance, as the work of Weniger has shown, sequence transformations like the Levin transformation that are based on expansions in powers of $`1/n`$, i.e., the asymptotic scale $`\varphi _j(n)=1/(n+\beta )^j`$, are not always superior to, and even often worse than those based upon factorial series, like Wenigerโs $`๐ฎ`$ transformation that is based on the asymptotic scale $`\psi _j(n)=1/(n+\beta )_j`$. To find an optimal asymptotic scale in combination with nonlinear sequence transformations seems to be an open mathematical problem.
Certainly, the proper choice of remainder estimates is also crucial in the context of Levin-type sequence transformations. See also Section 9.
## 6 Properties of Levin-Type Transformations
### 6.1 Basic Properties
Directly from the definition in Eqs. (65) and (66), we obtain the following theorem. The proof is left to the interested reader.
###### Theorem 5
Any Levin-type sequence transformation $`๐ฏ`$ is quasilinear, i.e., we have
$$๐ฏ{}_{n}{}^{(k)}(\{\{As_n+B\}\},\{\{\omega _n\}\})=A๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\})+B$$
(234)
for arbitrary constants $`A`$ and $`B`$. It is multiplicatively invariant in $`\omega _n`$, i.e., we have
$$๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{C\omega _n\}\})=๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\})$$
(235)
for arbitrary constants $`C0`$.
For a coefficient set $`\mathrm{\Lambda }`$ define the sets $`Y{}_{n}{}^{(k)}[\mathrm{\Lambda }]`$ by
$$Y{}_{n}{}^{(k)}[\mathrm{\Lambda }]=\{(x_0,\mathrm{},x_k)๐ฝ^{k+1}|\underset{j=0}{\overset{k}{}}\lambda {}_{n,j}{}^{(k)}/x_j0\}.$$
(236)
Since $`๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\})`$ for given coefficient set $`\mathrm{\Lambda }`$ depends only on the $`2k+2`$ numbers $`s_n,\mathrm{},s_{n+k}`$ and $`\omega _n,\mathrm{},\omega _{n+k}`$, it may be regarded as a mapping
$$U{}_{n}{}^{(k)}:^{k+1}\times Y{}_{n}{}^{(k)}[\mathrm{\Lambda }],(x,y)U{}_{n}{}^{(k)}(x|y)$$
(237)
such that
$$๐ฏ{}_{n}{}^{(k)}=U{}_{n}{}^{(k)}(s_n,\mathrm{},s_{n+k}|\omega _n,\mathrm{},\omega _{n+k})$$
(238)
The following theorem is a generalization of theorems for the $`๐ฅ`$ transformation \[36, Theorem 5\] and the $``$ transformation \[44, Theorem 5\].
###### Theorem 6
1. The $`๐ฏ^{(k)}`$ transformation can be regarded as continous mapping $`U_n^{(k)}`$ on $`^{k+1}\times Y{}_{n}{}^{(k)}[\mathrm{\Lambda }]`$ where $`Y{}_{n}{}^{(k)}[\mathrm{\Lambda }]`$ is defined in Eq. (236).
2. According to Theorem 5, $`U_n^{(k)}`$ is a homogeneous function of first degree in the first $`(k+1)`$ variables and a homogeneous function of degree zero in the last $`(k+1)`$ variables. Hence, for all vectors $`\stackrel{}{x}^{k+1}`$ and $`\stackrel{}{y}Y{}_{n}{}^{(k)}[\mathrm{\Lambda }]`$ and for all complex constants $`s`$ and $`t0`$ the equations
$$\begin{array}{cc}& U{}_{n}{}^{(k)}(s\stackrel{}{x}|\stackrel{}{y})=sU{}_{n}{}^{(k)}(\stackrel{}{x}|\stackrel{}{y}),\hfill \\ & U{}_{n}{}^{(k)}(\stackrel{}{x}|t\stackrel{}{y})=U{}_{n}{}^{(k)}(\stackrel{}{x}|\stackrel{}{y})\hfill \end{array}$$
(239)
hold.
3. $`U_n^{(k)}`$ is linear in the first $`(k+1)`$ variables. Thus, for all vectors $`\stackrel{}{x}^{k+1}`$, $`\stackrel{}{x}^{}^{k+1}`$, und $`\stackrel{}{y}Y{}_{n}{}^{(k)}[\mathrm{\Lambda }]`$
$$U{}_{n}{}^{(k)}(\stackrel{}{x}+\stackrel{}{x}^{}|\stackrel{}{y})=U{}_{n}{}^{(k)}(\stackrel{}{x}|\stackrel{}{y})+U{}_{n}{}^{(k)}(\stackrel{}{x}^{}|\stackrel{}{y})$$
(240)
holds.
4. For all constant vectors $`\stackrel{}{c}=(c,c,\mathrm{},c)^{k+1}`$ and all vectors $`\stackrel{}{y}Y{}_{n}{}^{(k)}[\mathrm{\Lambda }]`$ we have
$$U{}_{n}{}^{(k)}(\stackrel{}{c}|\stackrel{}{y})=c.$$
(241)
PROOF. These are immediate consequences of the definitions. $`\mathrm{}`$
### 6.2 The Limiting Transformation
We note that if a limiting transformation $`\stackrel{}{๐ฏ}[\stackrel{}{\Lambda }]`$ exists, it is also of Levin-type, and thus, the above theorems apply to the limiting transformation as well.
Also, we have the following result for the kernel of the limiting transformation:
###### Theorem 7
Suppose that for a Levin-type sequence transformation $`๐ฏ^{(k)}`$ of order $`k`$ there exists a limiting transformation $`\stackrel{}{๐ฏ}^{(k)}`$ with characteristic polynomial $`\stackrel{}{\Pi }^k`$ given by
$$\stackrel{}{\Pi }{}_{}{}^{(k)}(z)=\underset{j=0}{\overset{k}{}}\stackrel{}{๐}{}_{j}{}^{(k)}z_{}^{j}=\underset{\mathrm{}=1}{\overset{M}{}}(z\zeta _{\mathrm{}})^m_{\mathrm{}}$$
(242)
where the zeroes $`\zeta _{\mathrm{}}0`$ have multiplicities $`m_{\mathrm{}}`$. Then the kernel of the limiting transformation consists of all sequences $`\{\{s_n\}\}`$ with elements of the form
$$\sigma _n=\sigma +\omega _n\underset{\mathrm{}=1}{\overset{M}{}}\zeta _{\mathrm{}}^nP_{\mathrm{}}(n)$$
(243)
where $`P_{\mathrm{}}^{m_{\mathrm{}}1}`$ are arbitrary polynomials and $`\{\{\omega _n\}\}\stackrel{}{๐}^{(k)}`$.
PROOF. This follows directly from the observation that for such sequences $`(\sigma _n\sigma )/\omega _n`$ is nothing but a finite linear combination of the solutions $`\phi {}_{n,\mathrm{},j_{\mathrm{}}}{}^{(k)}=n^j_{\mathrm{}}\zeta _{\mathrm{}}^n`$ with $`\mathrm{}=1,\mathrm{},M`$ and $`j_{\mathrm{}}=0,\mathrm{},m_{\mathrm{}}1`$ of the recursion relation
$$\underset{j=0}{\overset{k}{}}\stackrel{}{๐}{}_{j}{}^{(k)}v_{n+j}^{}=0,$$
(244)
and thus, it is annihilated by $`๐ซ[\stackrel{}{\Pi }{}_{}{}^{(k)}]`$. $`\mathrm{}`$
### 6.3 Application to Power Series
Here, we generalize some results of Weniger that regard the application of Levin-type sequence transformations to power series.
We use the definitions in Eq. (27).
Like Padรฉ approximants, Levin-type sequence transformations yield rational approximants when applied to the partial sums $`f_n(z)`$ of a power series $`f(z)`$ with terms $`a_j=c_jz^j`$. These approximations offer a practical way for the analytical continuation of power series to regions outside of their circle of convergence. Furthermore, the poles of the rational approximations model the singularities of $`f(z)`$. They may also be used to approximate further terms beyond the last one used in constructing the rational approximant.
When applying a Levin-type sequence transformation $`๐ฏ`$ to a power series, remainder estimates $`\omega _n=m_nz^{\gamma +n}`$ will be used. We note that $`t`$ variants correspond to $`m_n=c_n`$, $`\gamma =0`$, $`u`$ variants correspond to $`m_n=c_n(n+\beta )`$, $`\gamma =0`$, $`\stackrel{~}{t}`$ variants to $`m_n=c_{n+1}`$, $`\gamma =1`$. Thus, for these variants, $`m_n`$ is independent of $`z`$ (Case A). For $`v`$ variants, we have $`m_n=c_{n+1}c_n/(c_nc_{n+1}z)`$, and $`\gamma =1`$. In this case, $`1/m_n^{(1)}`$ is a linear function of $`z`$ (Case B).
Application of $`๐ฏ`$ yields after some simplification
$$๐ฏ{}_{n}{}^{(k)}(\{\{f_n(z)\}\},\{\{m_nz^{\gamma +n}\}\})=\frac{{\displaystyle \underset{\mathrm{}=0}{\overset{n+k}{}}}z^{\mathrm{}}{\displaystyle \underset{j=\mathrm{max}(0,k\mathrm{})}{\overset{k}{}}}{\displaystyle \frac{\lambda _{n,j}^{(k)}}{m_{n+j}}}c_{\mathrm{}(kj)}}{{\displaystyle \underset{j=0}{\overset{k}{}}}{\displaystyle \frac{\lambda _{n,j}^{(k)}}{m_{n+j}}}z^{kj}}=\frac{P{}_{n}{}^{(k)}[T](z)}{Q{}_{n}{}^{(k)}[T](z)},$$
(245)
where in Case A, we have $`P{}_{n}{}^{(k)}[T]^{n+k}`$, $`Q{}_{n}{}^{(k)}[T]^k`$, and in Case B, we have $`P{}_{n}{}^{(k)}[T]^{n+k+1}`$, $`Q{}_{n}{}^{(k)}[T]^{k+1}`$. One needs the $`k+1+\gamma `$ partial sums $`f_n(z),\mathrm{},f_{n+k+\gamma }(z)`$ to compute these rational approximants. This should be compared to the fact that for the computation of the Padรฉ approximant $`[n+k+\gamma /k+\gamma ]`$ one needs the $`2k+2\gamma +1`$ partial sums $`f_n(z),\mathrm{},f_{n+2k+2\gamma }(z)`$.
We show that Taylor expansion of these rational approximants reproduces all terms of power series that have been used to calculate the rational approximation.
###### Theorem 8
We have
$$๐ฏ{}_{n}{}^{(k)}(\{\{f_n(z)\}\},\{\{m_nz^{\gamma +n}\}\})f(z)=O(z^{n+k+1+\tau })$$
(246)
where $`\tau =0`$ for $`t`$ and $`u`$ variants corresponding to $`m_n=c_n`$, $`\gamma =0`$, or $`m_n=c_n(n+\beta )`$, $`\gamma =0`$, respectively, while $`\tau =1`$ holds for the $`v`$ variant corresponding to $`m_n=c_{n+1}c_n/(c_nc_{n+1}z)`$, $`\gamma =1`$, and for the $`\stackrel{~}{t}`$ variants corresponding to $`m_n=c_{n+1}`$, $`\gamma =1`$, one obtains $`\tau =1`$ if $`๐ฏ`$ is convex.
PROOF. Using the identity
$$๐ฏ{}_{n}{}^{(k)}(\{\{f_n(z)\}\},\{\{m_nz^{\gamma +n}\}\})=f(z)+๐ฏ{}_{n}{}^{(k)}(\{\{f_n(z)f(z)\}\},\{\{m_nz^{\gamma +n}\}\})=$$
(247)
that follows from Theorem 5, we obtain after some easy algebra
$$๐ฏ{}_{n}{}^{(k)}(\{\{f_n(z)\}\},\{\{m_nz^{\gamma +n}\}\})f(z)=z^{n+k+1}\frac{{\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}}z^{\mathrm{}}{\displaystyle \underset{j=0}{\overset{k}{}}}{\displaystyle \frac{\lambda _{n,j}^{(k)}}{m_{n+j}}}c_{\mathrm{}+n+j+1}}{{\displaystyle \underset{j=0}{\overset{k}{}}}{\displaystyle \frac{\lambda _{n,j}^{(k)}}{m_{n+j}}}z^{kj}}.$$
(248)
This shows that the right hand side is at least $`O(z^{n+k+1})`$ since the denominator is $`O(1)`$ due to $`\lambda {}_{n,k}{}^{(k)}0`$. For the $`\stackrel{~}{t}`$ variant, the term corresponding to $`\mathrm{}=0`$ in the numerator is $`_{j=0}^k\lambda {}_{n,j}{}^{(k)}=\mathrm{\Pi }{}_{n}{}^{(k)}(1)`$ that vanishes for convex $`๐ฏ`$. For the $`v`$ variant, that term is $`_{j=0}^k\lambda {}_{n,j}{}^{(k)}(c_{n+j}c_{n+j+1}z)/c_{n+j}`$ that simplifies to $`(z)_{j=0}^k\lambda {}_{n,j}{}^{(k)}c_{n+j+1}^{}/c_{n+j}`$ for convex $`๐ฏ`$. This finishes the proof. $`\mathrm{}`$
## 7 Convergence Acceleration Results <br>for Levin-Type Transformations
### 7.1 General Results
We note that Germain-Bonne developed a theory of the regularity and convergence acceleration properties of sequence transformations that was later extended by Weniger \[84, Section 12\], \[88, Section 6\] to sequence transformations that depend explicitly on $`n`$ and on an auxiliary sequence of remainder estimates. The essential results of this theory apply to convergence acceleration of linearly convergent sequences. Of course, this theory can be applied to Levin-type sequence transformations. However, for the latter transformations, many results can be obtained more easily and also, one may obtain results of a general nature that are also applicable to other convergence types like logarithmic convergence. Thus, we are not going to use the Germain-Bonne-Weniger theory in the present article.
Here, we present some general convergence acceleration results for Levin-type sequence transformations that have a limiting transformation. The results, however, do not completely determine which transformation provides the best extrapolation results for a given problem sequence since the results are asymptotic in nature, but in practice, one is interested in obtaining good extrapolation results from as few members of the problem sequence as possible. Thus, it may well be that transformations with the same asymptotic behavior of the results perform rather differently in practice.
Nevertheless, the results presented below provide a first indication which results one may expect for large classes of Levin-type sequence transformations.
First, we present some results that show that the limiting transformation essentially determines for which sequences Levin-type sequence transformations are accelerative. The speed of convergence will be analyzed later.
###### Theorem 9
Assume that the following asymptotic relations hold for large $`n`$:
$$\lambda {}_{n,j}{}^{(k)}\stackrel{}{๐}{}_{j}{}^{(k)},\stackrel{}{๐}{}_{k}{}^{(k)}0,$$
(249)
$$\frac{s_ns}{\omega _n}\underset{\nu =1}{\overset{A}{}}c_\nu \zeta _\nu ^n,c_\nu \zeta _\nu 0,\stackrel{}{\Pi }{}_{}{}^{(k)}(\zeta _\nu )=0,$$
(250)
$$\frac{\omega _{n+1}}{\omega _n}\rho 0,\stackrel{}{\Pi }{}_{}{}^{(k)}(1/\rho )0.$$
(251)
Then, $`\{\{๐ฏ{}_{n}{}^{(k)}\}\}`$ accelerates $`\{\{s_n\}\}`$ to $`s`$, i.e., we have
$$\underset{n\mathrm{}}{lim}\frac{๐ฏ{}_{n}{}^{(k)}s}{s_ns}=0.$$
(252)
PROOF. Rewriting
$$\frac{๐ฏ{}_{n}{}^{(k)}s}{s_ns}=\frac{\omega _n}{s_ns}\frac{{\displaystyle \underset{j=0}{\overset{k}{}}}\lambda {}_{n,j}{}^{(k)}{\displaystyle \frac{s_{n+j}s}{\omega _{n+j}}}}{{\displaystyle \underset{j=0}{\overset{k}{}}}\lambda {}_{n,j}{}^{(k)}{\displaystyle \frac{\omega _n}{\omega _{n+j}}}}$$
(253)
one may perform the limit for $`n\mathrm{}`$ upon using the assumptions according to
$$\frac{๐ฏ{}_{n}{}^{(k)}s}{s_ns}\frac{{\displaystyle \underset{j=0}{\overset{k}{}}}\stackrel{}{๐}{\displaystyle {}_{j}{}^{(k)}\underset{\nu }{}}c_\nu \zeta _\nu ^{n+j}}{{\displaystyle \underset{\nu }{}}c_\nu \zeta _\nu ^n{\displaystyle \underset{j=0}{\overset{k}{}}}\stackrel{}{๐}{}_{j}{}^{(k)}\rho _{}^{j}}=\frac{{\displaystyle \underset{\nu }{}}c_\nu \zeta _\nu ^n\stackrel{}{\Pi }{}_{}{}^{(k)}(\zeta _\nu )}{\stackrel{}{\Pi }{}_{}{}^{(k)}(1/\rho ){\displaystyle \underset{\nu }{}}c_\nu \zeta _\nu ^n}=0$$
(254)
since $`\omega _n/\omega _{n+j}\rho ^j`$. $`\mathrm{}`$
Thus, the zeroes $`\zeta _\nu `$ of the characteristic polynomial of the limiting transformation are of particular importance.
It should be noted that the above assumptions correspond to a more complicated convergence type than linear or logarithmic convergence if $`\left|\zeta _1\right|=\left|\zeta _2\right|\left|\zeta _3\right|\mathrm{}`$. This is the case, for instance, for the $`^{(k)}`$ transformation where the limiting transformation has the characteristic polynomial $`P^{(2k)}(\alpha )`$ with $`k`$-fold zeroes at $`\mathrm{exp}(\text{i}\alpha )`$ and $`\mathrm{exp}(\text{i}\alpha )`$. Another example is the $`^{(k)}`$ transformation where the limiting transformation has characteristic polynomials $`Q^{(2k)}(\alpha )`$ with zeroes at $`\mathrm{exp}(\pm \text{i}\alpha )/\mathrm{\Theta }_j`$, $`j=0,\mathrm{},k1`$.
Specializing to $`A=1`$ in Theorem 9, we obtain the following corollary:
###### Corollary 1
Assume that the following asymptotic relations hold for large $`n`$:
$$\lambda {}_{n,j}{}^{(k)}\stackrel{}{๐}{}_{j}{}^{(k)},\stackrel{}{๐}{}_{k}{}^{(k)}0,$$
(255)
$$\frac{s_ns}{\omega _n}cq^n,cq0,\stackrel{}{\Pi }{}_{}{}^{(k)}(q)=0,$$
(256)
$$\frac{\omega _{n+1}}{\omega _n}\rho 0,\stackrel{}{\Pi }{}_{}{}^{(k)}(1/\rho )0.$$
(257)
Then, $`\{\{๐ฏ{}_{n}{}^{(k)}\}\}`$ accelerates $`\{\{s_n\}\}`$ to $`s`$, i.e., we have
$$\underset{n\mathrm{}}{lim}\frac{๐ฏ{}_{n}{}^{(k)}s}{s_ns}=0.$$
(258)
Note that the assumptions of Corollary 1 imply
$$\frac{s_{n+1}s}{s_ns}=\frac{s_{n+1}s}{\omega _{n+1}}\frac{\omega _n}{s_ns}\frac{\omega _{n+1}}{\omega _n}\rho \frac{cq^{n+1}}{cq^n}=\rho q$$
(259)
and thus, Corollary 1 corresponds to linear convergence for $`0<\left|\rho q\right|<1`$ and to logarithmic convergence for $`\rho q=1`$.
Many important sequence transformations have convex limiting transformations, i.e., the characteristic polynomials satisfy $`\stackrel{}{\Pi }{}_{}{}^{(k)}(1)=0`$. In this case, they accelerate linear convergence. More exactly, we have the following corollary:
###### Corollary 2
Assume that the following asymptotic relations hold for large $`n`$:
$$\lambda {}_{n,j}{}^{(k)}\stackrel{}{๐}{}_{j}{}^{(k)},\stackrel{}{๐}{}_{k}{}^{(k)}0,$$
(260)
$$\frac{s_ns}{\omega _n}c,c0,\stackrel{}{\Pi }{}_{}{}^{(k)}(1)=0,$$
(261)
$$\frac{\omega _{n+1}}{\omega _n}\rho 0,\stackrel{}{\Pi }{}_{}{}^{(k)}(1/\rho )0.$$
(262)
Then, $`\{\{๐ฏ{}_{n}{}^{(k)}\}\}`$ accelerates $`\{\{s_n\}\}`$ to $`s`$, i.e., we have
$$\underset{n\mathrm{}}{lim}\frac{๐ฏ{}_{n}{}^{(k)}s}{s_ns}=0.$$
(263)
Hence, any Levin-type sequence transformation with a convex limiting transformation accelerates linearly convergent sequences with
$$\underset{n\mathrm{}}{lim}\frac{s_{n+1}s}{s_ns}=\rho ,0<\left|\rho \right|<1$$
(264)
such that $`\stackrel{}{\Pi }{}_{}{}^{(k)}(1/\rho )0`$ for suitably chosen remainder estimates $`\omega _n`$ satisfying $`(s_ns)/\omega _nc0`$.
PROOF. Specializing Corollary 1 to $`q=1`$, it suffices to prove the last assertion. Here, the proof follows from the observation that $`(s_{n+1}s)/(s_ns)\rho `$ and $`(s_ns)/\omega _nc`$ imply $`\omega _{n+1}/\omega _n\rho `$ for large $`n`$ in view of the assumptions. $`\mathrm{}`$
Note that Corollary 2 applies for instance to suitable variants of the Levin transformation, the $`{}_{p}{}^{}๐`$ transformation and, more generally, of the $`๐ฅ`$ transformation. In particular, it applies to $`t`$, $`\stackrel{~}{t}`$, $`u`$ and $`v`$ variants, since in the case of linear convergence, one has $`\mathrm{}s_n/\mathrm{}s_{n1}\rho `$ which entails $`(s_ns)/\omega _nc`$ for all these variants by simple algebra.
Now, some results for the speed of convergence are given. Matos presented convergence theorems for sequence transformations based on annihilation difference operators with characteristic polynomials with constants coefficients that are close in spirit to the theorems given below. However, it should be noted that the theorems presented here apply to large classes of Levin-type transformations that have a limiting transformation (the latter, of course, has a characteristic polynomial with constants coefficients).
###### Theorem 10
Suppose that for a Levin-type sequence transformation $`๐ฏ^{(k)}`$ of order $`k`$ there is a limiting transformation $`\stackrel{}{๐ฏ}^{(k)}`$ with characteristic polynomial $`\stackrel{}{\Pi }^k`$ given by Eq. (242) where the multiplicities $`m_{\mathrm{}}`$ of the zeroes $`\zeta _{\mathrm{}}0`$ satisfy $`m_1m_2\mathrm{}m_M`$. Let
$$\lambda {}_{n,j}{}^{(k)}\frac{n^{m_11}}{(n+j)^{m_11}}\stackrel{}{๐}{}_{j}{}^{(k)}\left(\underset{t=0}{\overset{\mathrm{}}{}}\frac{e_t^{(k)}}{(n+j)^t}\right),e{}_{0}{}^{(k)}=1$$
(265)
for $`n\mathrm{}`$.
Assume that $`\{\{s_n\}\}๐^๐`$ and $`\{\{\omega _n\}\}๐^๐`$. Assume further that for $`n\mathrm{}`$ the asymptotic expansion
$$\frac{s_ns}{\omega _n}\underset{\mathrm{}=1}{\overset{M}{}}\zeta _{\mathrm{}}^n\underset{r=0}{\overset{\mathrm{}}{}}c_{\mathrm{},r}n^r$$
(266)
holds, and put
$$r_{\mathrm{}}=\mathrm{min}\{r_0|f_{\mathrm{},r+m_1}0\}$$
(267)
where
$$f_{\mathrm{},v}=\underset{r=0}{\overset{v}{}}e{}_{vr}{}^{(k)}c_{\mathrm{},r}^{}$$
(268)
and
$$B_{\mathrm{}}=(1)^m_{\mathrm{}}\frac{d^m_{\mathrm{}}\stackrel{}{\Pi }^{(k)}}{dx^m_{\mathrm{}}}(\zeta _{\mathrm{}})$$
(269)
for $`\mathrm{}=1,\mathrm{},M`$.
Assume that the following limit exists and satisfies
$$0\underset{n\mathrm{}}{lim}\frac{\omega _{n+1}}{\omega _n}=\rho \{\zeta _{\mathrm{}}^1|\mathrm{}=1,\mathrm{},M\}.$$
(270)
Then we have
$$\frac{๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\})s}{\omega _n}\frac{{\displaystyle \underset{\mathrm{}=1}{\overset{M}{}}}f_{\mathrm{},r_{\mathrm{}}+m_1}\zeta _{\mathrm{}}^{n+m_{\mathrm{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{r_{\mathrm{}}+m_{\mathrm{}}}{r_{\mathrm{}}}}\right){\displaystyle \frac{B_{\mathrm{}}}{n^{r_{\mathrm{}}+m_{\mathrm{}}m_1}}}}{\stackrel{}{\Pi }{}_{}{}^{(k)}(1/\rho )}\frac{1}{n^{2m_1}}.$$
(271)
Thus, $`\{\{๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\})\}\}`$ accelerates $`\{\{s_n\}\}`$ to $`s`$ at least with order $`2m_1`$, i.e.,
$$\frac{๐ฏ{}_{n}{}^{(k)}s}{s_ns}=O(n^{2m_1\tau }),\tau 0,$$
(272)
if $`c_{\mathrm{},0}0`$ for all $`\mathrm{}`$.
PROOF. We rewrite $`๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\})=๐ฏ_n^{(k)}`$ as defined in Eq. (11) in the form
$$๐ฏ{}_{n}{}^{(k)}s=\omega _n\frac{{\displaystyle \underset{j=0}{\overset{k}{}}}\lambda {}_{n,j}{}^{(k)}{\displaystyle \frac{s_{n+j}s}{\omega _{n+j}}}}{{\displaystyle \underset{j=0}{\overset{k}{}}}\lambda {}_{n,j}{}^{(k)}{\displaystyle \frac{\omega _n}{\omega _{n+j}}}}\omega _n\frac{{\displaystyle \underset{j=0}{\overset{k}{}}}\stackrel{}{๐}{\displaystyle {}_{j}{}^{(k)}\underset{t=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{e_t^{(k)}}{(n+j)^t}}{\displaystyle \frac{(n+j)^{m_11}}{n^{m_11}}}{\displaystyle \frac{s_{n+j}s}{\omega _{n+j}}}}{{\displaystyle \underset{j=0}{\overset{k}{}}}\stackrel{}{๐}{}_{j}{}^{(k)}{\displaystyle \frac{1}{\rho ^j}}}$$
(273)
for large $`n`$ where we used Eq. (265) in the numerator, and in the denominator the relation $`\omega _n/\omega _{n+j}\rho ^j`$ that follows by repeated application of Eq. (270). Insertion of (266) now yields
$$๐ฏ{}_{n}{}^{(k)}s\frac{\omega _n}{n^{m_11}\stackrel{}{\Pi }{}_{}{}^{(k)}(1/\rho )}\underset{\mathrm{}=1}{\overset{M}{}}\underset{r=0}{\overset{\mathrm{}}{}}f_{\mathrm{},r+m_1}\underset{j=0}{\overset{k}{}}\stackrel{}{๐}{}_{j}{}^{(k)}\frac{\zeta _{\mathrm{}}^{n+j}}{(n+j)^{r+1}}$$
(274)
where Eq. (268) was used. Also the fact was used that $`๐ซ[\stackrel{}{\Pi }{}_{}{}^{(k)}]`$ annihilates any linear combination of the solutions $`\phi {}_{n,\mathrm{},j_{\mathrm{}}}{}^{(k)}=n^j_{\mathrm{}}\zeta _{\mathrm{}}^n`$ with $`\mathrm{}=1,\mathrm{},M`$ and $`j_{\mathrm{}}=0,\mathrm{},m_11`$ of the recursion relation (244) since each $`\zeta _{\mathrm{}}`$ is a zero with multiplicity exceeding $`m_11`$. Invoking Lemma 1 given in Appendix C one obtains
$$๐ฏ{}_{n}{}^{(k)}s\frac{\omega _n}{n^{m_11}\stackrel{}{\Pi }{}_{}{}^{(k)}(1/\rho )}\underset{\mathrm{}=1}{\overset{M}{}}\underset{r=0}{\overset{\mathrm{}}{}}f_{\mathrm{},r+m_1}\zeta _{\mathrm{}}^{n+m_{\mathrm{}}}\left(\genfrac{}{}{0pt}{}{r+m_{\mathrm{}}}{r}\right)\frac{(1)^m_{\mathrm{}}}{n^{r+m_{\mathrm{}}+1}}\frac{d^m_{\mathrm{}}\stackrel{}{\Pi }^{(k)}}{dx^m_{\mathrm{}}}(\zeta _{\mathrm{}})$$
(275)
The proof of Eq. (271) is completed taking leading terms in the sums over $`r`$. Since $`s_ns\omega _nZ^n_\mathrm{}I(\zeta _{\mathrm{}}/Z)^nc_{\mathrm{},0}`$ where $`Z=\mathrm{max}\{|\zeta _{\mathrm{}}||\mathrm{}=1,\mathrm{},M\}`$, and $`I=\{\mathrm{}=1,\mathrm{},M|Z=|\zeta _{\mathrm{}}|\}`$, Eq. (272) is obtained where $`\tau =\mathrm{min}\{r_{\mathrm{}}+m_{\mathrm{}}m_1|\mathrm{}I\}`$. $`\mathrm{}`$
If $`\omega _{n+1}/\omega _n\rho `$, where $`\stackrel{}{\Pi }{}_{}{}^{(k)}(1/\rho )=0`$, i.e., if (C-3) of Theorem 10 does not hold, then the denominators vanish asymptotically. In this case, one has to investigate whether the numerators or the denominators vanish faster.
###### Theorem 11
Assume that (C-1) and (C-2) of Theorem 10 hold.
Assume that for $`n\mathrm{}`$ the asymptotic relation
$$\frac{\omega _{n+1}}{\omega _n}\rho \mathrm{exp}(ฯต_n),\rho 0$$
(276)
holds where
$$\frac{1}{\lambda !}\frac{d^\lambda \stackrel{}{\Pi }^{(k)}}{dx^\lambda }(1/\rho )=\{\begin{array}{cc}0\hfill & \text{for }\lambda =0,\mathrm{},\mu 1\hfill \\ C0\hfill & \text{for }\lambda =\mu \hfill \end{array}$$
(277)
and
$$ฯต_n0,\frac{ฯต_{n+1}}{ฯต_n}1$$
(278)
for large $`n`$. Define $`\delta _n`$ via $`\mathrm{exp}(ฯต_n)=1+\delta _n\rho `$.
Then we have for large $`n`$
$$\frac{๐ฏ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\})s}{\omega _n}\frac{{\displaystyle \underset{\mathrm{}=1}{\overset{M}{}}}f_{\mathrm{},r_{\mathrm{}}+m_1}\zeta _{\mathrm{}}^{n+m_{\mathrm{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{r_{\mathrm{}}+m_{\mathrm{}}}{r_{\mathrm{}}}}\right){\displaystyle \frac{B_{\mathrm{}}}{n^{r_{\mathrm{}}+m_{\mathrm{}}m_1}}}}{C(\delta _n)^\mu }\frac{1}{n^{2m_1}}.$$
(279)
PROOF. The proof proceeds as the proof of Theorem 10 but in the denominator we use
$$\underset{j=0}{\overset{k}{}}\lambda {}_{n,j}{}^{(k)}\frac{\omega _n}{\omega _{n+j}}C(\delta _n)^\mu $$
(280)
that follows from Lemma 2 given in Appendix C. $`\mathrm{}`$
Thus, the effect of the sequence transformation in this case essentially depends on the question whether $`(\delta _n)^\mu n^{2m_1}`$ goes to 0 for large $`n`$ or not. In many important cases like the Levin transformation and the $`{}_{p}{}^{}๐`$ transformations, we have $`M=1`$ and $`m_1=k`$. We note that Theorem 11 becomes especially important in the case of logarithmic convergence since for instance for $`M=1`$ one observes that $`(s_{n+1}s)/(s_ns)1`$ and $`(s_ns)/\omega _n\zeta _1^nc_{1,0}0`$ imply $`\omega _{n+1}/\omega _n1/\zeta _1`$ for large $`n`$ such that the denominators vanish asymptotically. In this case, we have $`\mu =m_1`$ whence $`(\delta _n)^\mu n^{2m_1}=O(n^{m_1})`$ if $`\delta _n=O(1/n)`$. This reduction of the speed of convergence of the acceleration process from $`O(n^{2k})`$ to $`O(n^k)`$ in the case of logarithmic convergence is a generic behavior that is reflected in a number of theorems regarding convergence acceleration properties of Levin-type sequence transformations. Examples are Sidiโs theorem for the Levin transformation given below (Theorem 13), and for the $`{}_{p}{}^{}๐`$ transformation the Corollaries 3 and 4 given below. Compare also \[84, Theorems 13.5, 13.9, 13.11, 13.12, 14.2\].
The following theorem was given by Matos where the proof may be found. To formulate it, we define that a sequence $`\{\{u_n\}\}`$ has *property M* if it satisfies
$$\frac{u_{n+1}}{u_n}1+\frac{\alpha }{n}+r_n\text{with }r_n=o(1/n),\mathrm{}^{\mathrm{}}r_n=o(\mathrm{}^{\mathrm{}}(1/n))\text{for }n\mathrm{}.$$
(281)
###### Theorem 12
\[59, Theorem 13\] Let $`\{\{s_n\}\}`$ be a sequence such that
$$s_ns=\omega _n(a_1g{}_{1}{}^{(1)}(n)+\mathrm{}+a_kg{}_{1}{}^{(k)}(n)+\rho _n)$$
(282)
with $`g{}_{1}{}^{(j+1)}(n)=o(g{}_{1}{}^{(j)}(n))`$, $`\rho _n=o(g{}_{1}{}^{(k)}(n))`$ for $`n\mathrm{}`$. Let us consider an operator $`L`$ of the form (193) for which we know a basis of solutions $`\{\{u{}_{n}{}^{(j)}\}\}`$, $`j=1,\mathrm{},k`$, and each one can be written as
$$u{}_{n}{}^{(j)}\underset{m=1}{\overset{\mathrm{}}{}}\alpha {}_{m}{}^{(j)}g{}_{m}{}^{(j)}(n),g{}_{m+1}{}^{(j)}(n)=o(g{}_{m}{}^{(j)}(n))$$
(283)
as $`n\mathrm{}`$ for all $`m`$ and $`j=1,\mathrm{},k`$. Suppose that
$$\begin{array}{cc}\hfill a)& g{}_{2}{}^{(j+1)}(n)=o(g{}_{2}{}^{(j)}(n))\text{ for }n\mathrm{},j=1,\mathrm{},k1\hfill \\ \hfill b)& g{}_{2}{}^{(1)}(n)=o(g{}_{1}{}^{(k)}(n)),\text{ and }\rho _nKg{}_{2}{}^{(1)}(n)\text{for }n\mathrm{}\hfill \\ \hfill c)& \{\{g{}_{m}{}^{(j)}(n)\}\}\text{has property M for }\hfill \\ & m,j=1,\mathrm{},k.\hfill \end{array}$$
(284)
Then
1. If $`\{\{\omega _n\}\}`$ satisfies $`lim_n\mathrm{}\omega _n/\omega _{n+1}=\lambda 1`$, the sequence transformation $`๐ฏ_n^{(k+1)}`$ corresponding to the operator $`L`$ accelerates the convergence of $`\{\{s_n\}\}`$. Moreover, the acceleration can be measured by
$$\frac{๐ฏ{}_{n}{}^{(k+1)}s}{s_ns}Cn^k\frac{g{}_{2}{}^{(1)}(n)}{g{}_{1}{}^{(1)}(n)},n\mathrm{}.$$
(285)
2. If $`\{\{1/\omega _n\}\}`$ has property M, then the speed of convergence of $`๐ฏ_n^{(k+1)}`$ can be measured by
$$\frac{๐ฏ{}_{n}{}^{(k+1)}s}{s_ns}C\frac{g{}_{2}{}^{(1)}(n)}{g{}_{1}{}^{(1)}(n)},n\mathrm{}.$$
(286)
### 7.2 Results for Special Cases
In the case that peculiar properties of a Levin-type sequence transformation are used, more stringent theorems can often be proved as regards convergence acceleration using this particular transformation.
In the case of the Levin transformation, Sidi proved the following theorem:
###### Theorem 13
\[14, Theorem 2.32\] If $`s_n=s+\omega _nf_n`$ where $`f_n_{j=0}^{\mathrm{}}\beta _j/n^j`$ with $`\beta _00`$ and $`\omega _n_{j=0}^{\mathrm{}}\delta _j/n^{j+a}`$ with $`a>0`$, $`\delta _00`$ for $`n\mathrm{}`$ then, if $`\beta _k0`$
$${}_{n}{}^{(k)}s\frac{\delta _0\beta _k}{\left(\genfrac{}{}{0pt}{}{a}{k}\right)}n^{ak}(n\mathrm{}).$$
(287)
For the $`W`$ algorithm and the $`d^{(1)}`$ transformation that may be regarded as direct generalizations of the Levin transformation, Sidi has obtained a large number of results. The interested reader is referred to the literature (See and references therein).
Convergence results for the Levin transformation, the Drummond transformation and the Weniger transformations may be found in Section 13 of Wenigerโs report .
Results for the $`๐ฅ`$ transformation and in particular, for the $`{}_{p}{}^{}๐`$ transformation are given in . Here, we recall the following theorems:
###### Theorem 14
Assume that the following holds:
The sequence $`\{\{s_n\}\}`$ has the (anti)limit $`s`$.
For every $`n`$, the elements of the sequence $`\{\{\omega _n\}\}`$ are strictly alternating in sign and do not vanish.
For all $`n`$ and $`k`$, the elements of the sequence $`\{\{\delta {}_{n}{}^{(k)}\}\}=\{\{\mathrm{\Delta }r{}_{n}{}^{(k)}\}\}`$ are of the same sign and do not vanish.
For all $`n_0`$ the ratio $`(s_ns)/\omega _n`$ can be expressed as a series of the form
$$\frac{s_ns}{\omega _n}=c_0+\underset{j=1}{\overset{\mathrm{}}{}}c_j\underset{n>n_1>n_2>\mathrm{}>n_j}{}\delta {}_{n_1}{}^{(0)}\delta {}_{n_2}{}^{(1)}\mathrm{}\delta _{n_j}^{(j1)}$$
(288)
with $`c_00`$.
Then the following holds for $`s{}_{n}{}^{(k)}=๐ฅ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\},\{\{\delta {}_{n}{}^{(k)}\}\})`$:
The error $`s{}_{n}{}^{(k)}s`$ satisfies
$$s{}_{n}{}^{(k)}s=\frac{b_n^{(k)}}{{}_{n}{}^{(k1)}{}_{n}{}^{(k2)}\mathrm{}{}_{n}{}^{(0)}[1/\omega _n]}$$
(289)
with
$$b{}_{n}{}^{(k)}=c_k+\underset{j=k+1}{\overset{\mathrm{}}{}}c_j\underset{n>n_{k+1}>n_{k+2}>\mathrm{}>n_j}{}\delta {}_{n_{k+1}}{}^{(k)}\delta {}_{n_{k+2}}{}^{(k+1)}\mathrm{}\delta {}_{n_j}{}^{(j1)}.$$
(290)
The error $`s{}_{n}{}^{(k)}s`$ is bounded in magnitude according to
$$|s{}_{n}{}^{(k)}s||\omega _nb{}_{n}{}^{(k)}\delta {}_{n}{}^{(0)}\delta {}_{n}{}^{(1)}\mathrm{}\delta {}_{n}{}^{(k1)}|.$$
(291)
For large $`n`$ the estimate
$$\frac{s{}_{n}{}^{(k)}s}{s_ns}=O(\delta {}_{n}{}^{(0)}\delta {}_{n}{}^{(1)}\mathrm{}\delta {}_{n}{}^{(k1)})$$
(292)
holds if $`b{}_{n}{}^{(k)}=O(1)`$ and $`(s_ns)/\omega _n=O(1)`$ as $`n\mathrm{}`$.
###### Theorem 15
Define $`s{}_{n}{}^{(k)}=๐ฅ{}_{n}{}^{(k)}(\{\{s_n\}\},\{\{\omega _n\}\},\{\{\delta {}_{n}{}^{(k)}\}\})`$ and $`\omega {}_{n}{}^{(k)}=1/D_n^{(k)}`$ where the $`D_n^{(k)}`$ are defined as in Eq. (94). Put $`e{}_{n}{}^{(k)}=1\omega {}_{n+1}{}^{(k)}/\omega _n^{(k)}`$ and $`b{}_{n}{}^{(k)}=(s{}_{n}{}^{(k)}s)/\omega _n^{(k)}`$. Assume that (A-0) of Theorem 14 holds and that the following conditions are satisfied:
Assume that
$$\underset{n\mathrm{}}{lim}\frac{b_n^{(k)}}{b_n^{(0)}}=B_k$$
(293)
exists and is finite.
Assume that
$$\mathrm{\Omega }_k=\underset{n\mathrm{}}{lim}\frac{\omega _{n+1}^{(k)}}{\omega _n^{(k)}}0$$
(294)
and
$$F_k=\underset{n\mathrm{}}{lim}\frac{\delta _{n+1}^{(k)}}{\delta _n^{(k)}}0,$$
(295)
exist for all $`k_0`$. Hence the limits $`\mathrm{\Phi }_k=lim_n\mathrm{}\mathrm{\Phi }_n^{(k)}`$ (compare Eq. (97)) exist for all $`k_0`$.
Then, the following holds:
If $`\mathrm{\Omega }_0\{\mathrm{\Phi }_0=1,\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_{k1}\}`$, then
$$\underset{n\mathrm{}}{lim}\frac{s{}_{n}{}^{(k)}s}{s_ns}\{\underset{l=0}{\overset{k1}{}}\delta {}_{n}{}^{(l)}\}^1=B_k\frac{[\mathrm{\Omega }_0]^k}{{\displaystyle \underset{l=0}{\overset{k1}{}}}(\mathrm{\Phi }_l\mathrm{\Omega }_0)}$$
(296)
and, hence,
$$\frac{s{}_{n}{}^{(k)}s}{s_ns}=O(\delta {}_{n}{}^{(0)}\delta {}_{n}{}^{(1)}\mathrm{}\delta {}_{n}{}^{(k1)})$$
(297)
holds in the limit $`n\mathrm{}`$.
If $`\mathrm{\Omega }_l=1`$ for $`l\{0,1,2,\mathrm{},k\}`$ then
$$\underset{n\mathrm{}}{lim}\frac{s{}_{n}{}^{(k)}s}{s_ns}\left\{\underset{l=0}{\overset{k1}{}}\frac{\delta _n^{(l)}}{e_n^{(l)}}\right\}^1=B_k$$
(298)
and, hence,
$$\frac{s{}_{n}{}^{(k)}s}{s_ns}=O\left(\underset{l=0}{\overset{k1}{}}\frac{\delta _n^{(l)}}{e_n^{(l)}}\right)$$
(299)
holds in the limit $`n\mathrm{}`$.
This theorem has the following two corollaries for the $`{}_{p}{}^{}๐`$ transformation :
###### Corollary 3
Assume that the following holds:
(C-1) Let $`\beta >0,p1`$ and $`\delta {}_{n}{}^{(k)}=\mathrm{}[(n+\beta +(p1)k)^1]`$. Thus, we deal with the $`{}_{p}{}^{}๐`$ transformation and, hence, the equations $`F_k=lim_n\mathrm{}\delta {}_{n+1}{}^{(k)}/\delta {}_{n}{}^{(k)}=1`$ and $`\mathrm{\Phi }_k=1`$ hold for all $`k`$.
(C-2) Assumptions (A-2) of Theorem 14 and (B-1) of Theorem 15 are satisfied for the particular choice (C-1) for $`\delta _n^{(k)}`$.
(C-3) The limit $`\mathrm{\Omega }_0=lim_n\mathrm{}\omega _{n+1}/\omega _n`$ exists, and it satisfies $`\mathrm{\Omega }_0\{0,1\}`$. Hence, all the limits $`\mathrm{\Omega }_k=lim_n\mathrm{}\omega {}_{n+1}{}^{(k)}/\omega _n^{(k)}`$ exist for $`k`$ exist and satisfy $`\mathrm{\Omega }_k=\mathrm{\Omega }_0`$.
Then the transformation $`s{}_{n}{}^{(k)}={}_{p}{}^{}๐{}_{n}{}^{(k)}(\beta ,\{\{s_n\}\},\{\{\omega _n\}\})`$ satisfies
$$\underset{n\mathrm{}}{lim}\frac{s{}_{n}{}^{(k)}s}{s_ns}\{\underset{l=0}{\overset{k1}{}}\delta {}_{n}{}^{(l)}\}^1=B_k\left\{\frac{\mathrm{\Omega }_0}{1\mathrm{\Omega }_0}\right\}^k$$
(300)
and, hence,
$$\frac{s{}_{n}{}^{(k)}s}{s_ns}=O\left((n+\beta )^{2k}\right)$$
(301)
holds in the limit $`n\mathrm{}`$.
Note that Corollary 3 can be applied in the case of linear convergence because then $`0<\left|\mathrm{\Omega }_0\right|<1`$ holds.
Corollary 3 allows to conclude that in the case of linear convergence, the $`{}_{p}{}^{}๐`$ transformations should be superior to Wynnโs epsilon algorithm . Consider for instance the case that
$$s_ns+\lambda ^nn^\theta \underset{n=0}{\overset{\mathrm{}}{}}c_j/n^j,c_00,n\mathrm{}$$
(302)
is an asymptotic expansion of the sequence elements $`s_n`$. Assuming $`\lambda 1`$ and $`\theta \{0,1,\mathrm{},k1\}`$ it follows that (\[102, p. 127\]; \[84, p. 333, Eq. (13.4-7)\])
$$\frac{ฯต{}_{2k}{}^{(n)}s}{s_ns}=O\left(n^{2k}\right),n\mathrm{}.$$
(303)
This is the same order of convergence acceleration as in Eq. (301). But it should be noted that for the computation of $`ฯต_{2k}^{(n)}`$ the $`2k+1`$ sequence elements $`\{s_n,\mathrm{},s_{n+2k}\}`$ are required. But for the computation of $`{}_{p}{}^{}๐`$$`{}_{}{}^{(k)}{}_{n}{}^{}`$ only the $`k+1`$ sequence elements $`\{s_n,\mathrm{},s_{n+k}\}`$ are required in the case of the $`t`$ and $`u`$ variants, and additionally $`s_{n+k+1}`$ in the case of the $`\stackrel{~}{t}`$ variant. Again, this is similar to Levin-type accelerators (\[84, p. 333\]).
The following corollary 4 applies to the case of logarithmic convergence:
###### Corollary 4
Assume that the following holds:
(D-1) Let $`\beta >0,p1`$ and $`\delta {}_{n}{}^{(k)}=\mathrm{\Delta }[(n+\beta +(p1)k)^1]`$. Thus, we deal with the $`{}_{p}{}^{}๐`$ transformation and, hence, the equations $`F_k=lim_n\mathrm{}\delta {}_{n+1}{}^{(k)}/\delta {}_{n}{}^{(k)}=1`$ and $`\mathrm{\Phi }_k=1`$ hold for all $`k`$
(D-2) Assumptions (A-2) of Theorem 14 and (B-1) of Theorem 15 are satisfied for the particular choice (C-1) for $`\delta _n^{(k)}`$.
(D-3) Some constants $`a_l^{(j)}`$, $`j=1,2`$, exist such that
$$e{}_{n}{}^{(l)}=1\omega {}_{n+1}{}^{(l)}/\omega {}_{n}{}^{(l)}=\frac{a_l^{(1)}}{n+\beta }+\frac{a_l^{(2)}}{(n+\beta )^2}+O((n+\beta )^3)$$
(304)
holds for $`l=0`$. This implies that this equation, and hence, $`\mathrm{\Omega }_l=1`$ holds for $`l\{0,1,2,\mathrm{},k\}`$. Assume further that $`a{}_{l}{}^{(1)}0`$ for $`l\{0,1,2,\mathrm{},k1\}`$.
Then the transformation $`s{}_{n}{}^{(k)}={}_{p}{}^{}๐{}_{n}{}^{(k)}(\beta ,\{\{s_n\}\},\{\{\omega _n\}\})`$ satisfies
$$\underset{n\mathrm{}}{lim}\frac{s{}_{n}{}^{(k)}s}{s_ns}\left\{\underset{l=0}{\overset{k1}{}}\frac{\delta _n^{(l)}}{e_n^{(l)}}\right\}^1=B_k$$
(305)
and, hence,
$$\frac{s{}_{n}{}^{(k)}s}{s_ns}=O\left((n+\beta )^k\right)$$
(306)
holds in the limit $`n\mathrm{}`$.
For convergence acceleration results regarding the $``$ and $``$ transformations, see and .
## 8 Stability Results for Levin-Type Transformations
### 8.1 General Results
We remind the reader of the definition of the *stability indices* $`๐ช{}_{n}{}^{(k)}(๐ฏ)1`$ as given in Eq. (76). We consider the sequence $`\{\{\omega _n\}\}๐^๐`$ as given. We call the transformation $`๐ฏ`$ *stable along the path* $`๐ซ=\{(n_{\mathrm{}},k_{\mathrm{}})|[n_{\mathrm{}}>n_\mathrm{}1\text{ and }k_{\mathrm{}}k_\mathrm{}1]\text{ or }[n_{\mathrm{}}n_\mathrm{}1\text{ and }k_{\mathrm{}}>k_\mathrm{}1]\}`$ in the $`๐ฏ`$ table if the limit of its stability index along the path $`๐ซ`$ exists and is bounded, i.e., if
$$\underset{\mathrm{}\mathrm{}}{lim}๐ช{}_{n_{\mathrm{}}}{}^{(k_{\mathrm{}})}(๐ฏ)=\underset{\mathrm{}\mathrm{}}{lim}\underset{j=0}{\overset{k_{\mathrm{}}}{}}|\gamma {}_{n_{\mathrm{}},j}{}^{(k_{\mathrm{}})}(\stackrel{}{\omega }_n)|<\mathrm{}$$
(307)
where the $`\gamma {}_{n,j}{}^{(k)}(\stackrel{}{\omega }_n)`$ are defined in Eq. (75). The transformation $`๐ฏ`$ is called *S-stable*, if it is stable along all paths $`๐ซ^{(k)}=\{(n,k)|n=0,1,\mathrm{}\}`$ for fixed $`k`$, i.e. along all columns in the $`๐ฏ`$ table.
The case of stability along diagonal paths is much more difficult to treat analytically unless Theorem 17 applies. Up to now it seems that such diagonal stability issues have only been analysed by Sidi for the case of the $`d^{(1)}`$ transformation (see and references therein). We will treat only S-stability in the sequel.
The higher the stability index $`๐ช(๐ฏ)`$ is, the smaller is the numerical stability of the transformation $`๐ฏ`$: If $`ฯต_j`$ is the numerical error of $`s_j`$,
$$ฯต_j=s_j\text{fl}\left(s_j\right)$$
(308)
then the difference between the true value $`๐ฏ_n^{(k)}`$ and the numerically computed approximation $`\text{fl}(๐ฏ{}_{n}{}^{(k)})`$ may be bounded according to
$$\left|๐ฏ{}_{n}{}^{(k)}\text{fl}(๐ฏ{}_{n}{}^{(k)})\right|๐ช{}_{n}{}^{(k)}(๐ฏ)(\underset{j\{0,1,\mathrm{},k\}}{\mathrm{max}}|ฯต_{n+j}|).$$
(309)
Compare also .
###### Theorem 16
If the Levin-type sequence transformation $`๐ฏ^{(k)}`$ has a limiting transformation $`\stackrel{}{๐ฏ}^{(k)}`$ with characteristic polynomial $`\stackrel{}{\Pi }{}_{}{}^{(k)}^{(k)}`$ for all $`k`$, and if $`\{\{\omega _n\}\}๐^๐`$ satisfies $`\omega _{n+1}/\omega _n\rho 0`$ for large $`n`$ with $`\stackrel{}{\Pi }{}_{}{}^{(k)}(1/\rho )0`$ for all $`k`$ then the transformation $`๐ฏ`$ is S-stable. If additionally, the coefficients $`\stackrel{}{๐}_j^{(k)}`$ of the characteristic polynomial alternate in sign, i.e., if $`\stackrel{}{๐}{}_{j}{}^{(k)}=(1)^j|\stackrel{}{๐}{}_{j}{}^{(k)}|/\tau _k`$ with $`|\tau _k|=1`$, then the limits $`\stackrel{}{\Gamma }{}_{}{}^{(k)}(๐ฏ)=lim_n\mathrm{}\mathrm{\Gamma }{}_{n}{}^{(k)}(๐ฏ)`$ obey
$$\stackrel{}{๐ช}{}_{}{}^{(k)}(๐ฏ)=\tau _k\frac{\stackrel{}{\Pi }{}_{}{}^{(k)}(1/|\rho |)}{|\stackrel{}{\Pi }{}_{}{}^{(k)}(1/\rho )|}$$
(310)
PROOF. We have for fixed $`k`$
$$\gamma {}_{n,j}{}^{(k)}(\stackrel{}{\omega }_n)=\lambda {}_{n,j}{}^{(k)}\frac{\omega _n}{\omega _{n+j}}\left[\underset{j^{}=0}{\overset{k}{}}\lambda {}_{n,j^{}}{}^{(k)}\frac{\omega _n}{\omega _{n+j^{}}}\right]^1\stackrel{}{๐}{}_{j}{}^{(k)}\rho _{}^{j}\left[\underset{j^{}=0}{\overset{k}{}}\stackrel{}{๐}{}_{j^{}}{}^{(k)}\rho _{}^{j^{}}\right]^1=\frac{\stackrel{}{๐}{}_{j}{}^{(k)}\rho _{}^{j}}{\stackrel{}{\Pi }{}_{}{}^{(k)}(1/\rho )}$$
(311)
whence
$$\underset{n\mathrm{}}{lim}๐ช{}_{n}{}^{(k)}(๐ฏ)=\underset{n\mathrm{}}{lim}\underset{j=0}{\overset{k}{}}|\gamma {}_{n,j}{}^{(k)}(\stackrel{}{\omega }_n)|=\frac{_{j=0}^k|\stackrel{}{๐}{}_{j}{}^{(k)}||\rho |^j}{|\stackrel{}{\Pi }{}_{}{}^{(k)}(1/\rho )|}<\mathrm{}.$$
(312)
If the $`\stackrel{}{๐}_j^{(k)}`$ alternate in sign, we obtain for these limits
$$\stackrel{}{๐ช}{}_{}{}^{(k)}(๐ฏ)=\tau _k\frac{_{j=0}^k\stackrel{}{๐}{}_{j}{}^{(k)}(|\rho |)_{}^{j}}{|\stackrel{}{\Pi }{}_{}{}^{(k)}(1/\rho )|}.$$
(313)
This implies Eq. (310). $`\mathrm{}`$
###### Corollary 5
Assume that the Levin-type sequence transformation $`๐ฏ^{(k)}`$ has a limiting transformation $`\stackrel{}{๐ฏ}^{(k)}`$ with characteristic polynomial $`\stackrel{}{\Pi }{}_{}{}^{(k)}^{(k)}`$ and the coefficients $`\stackrel{}{๐}_j^{(k)}`$ of the characteristic polynomial alternate in sign, i.e., if $`\stackrel{}{๐}{}_{j}{}^{(k)}=(1)^j|\stackrel{}{๐}{}_{j}{}^{(k)}|/\tau _k`$ with $`|\tau _k|=1`$ for all $`k`$. The sequence $`\{\{\omega _n\}\}๐^๐`$ is assumed to be alternating and to satisfy $`\omega _{n+1}/\omega _n\rho <0`$ for large $`n`$. Then the transformation $`๐ฏ`$ is S-stable. Additionally the limits are $`\stackrel{}{๐ช}{}_{}{}^{(k)}(๐ฏ)=1`$.
PROOF. Since
$$\underset{j^{}=0}{\overset{k}{}}\frac{\lambda _{n,j^{}}^{(k)}}{\tau _k^1}\frac{\omega _n}{\omega _{n+j^{}}}\frac{\stackrel{}{\Pi }{}_{}{}^{(k)}(1/\rho )}{\tau _k^1}=\underset{j^{}=0}{\overset{k}{}}\frac{\stackrel{}{๐}_j^{}^{(k)}}{\tau _k^1}(1)^j^{}|\rho |^j^{}=\underset{j^{}=0}{\overset{k}{}}|\stackrel{}{๐}{}_{j^{}}{}^{(k)}||\rho |^j^{}|\stackrel{}{๐}{}_{k}{}^{(k)}|/|\rho |^k>0$$
(314)
$`1/\rho `$ cannot be a zero of $`\stackrel{}{\Pi }^{(k)}`$. Then, Theorem 16 entail that $`๐ฏ`$ is S-stable. Furthermore, Eq. (310) is applicable and yields $`\stackrel{}{๐ช}{}_{}{}^{(k)}(๐ฏ)=1`$. $`\mathrm{}`$
This result can be improved if all the coefficients $`\lambda _{n,j}^{(k)}`$ are alternating:
###### Theorem 17
Assume that the Levin-type sequence transformation $`๐ฏ^{(k)}`$ has a characteristic polynomials $`\mathrm{\Pi }{}_{n}{}^{(k)}^{(k)}`$ with alternating coefficients $`\lambda _{n,j}^{(k)}`$ i.e., $`\lambda {}_{n,j}{}^{(k)}=(1)^j|\lambda {}_{n,j}{}^{(k)}|/\tau _k`$ with $`|\tau _k|=1`$ for all $`n_0`$ and $`k`$. The sequence $`\{\{\omega _n\}\}๐^๐`$ is assumed to be alternating and to satisfy $`\omega _{n+1}/\omega _n<0`$ for all $`n_0`$. Then we have $`๐ช{}_{n}{}^{(k)}(๐ฏ)=1`$. Hence, the transformation $`๐ฏ`$ is stable along all paths for such remainder estimates.
PROOF. We have for fixed $`n`$ and $`k`$
$$\gamma {}_{n,j}{}^{(k)}(\stackrel{}{\omega }_n)=\frac{\lambda {}_{n,j}{}^{(k)}\frac{\omega _n}{\omega _{n+j}}}{_{j^{}=0}^k\lambda {}_{n,j^{}}{}^{(k)}\frac{\omega _n}{\omega _{n+j^{}}}}=\frac{\lambda {}_{n,j}{}^{(k)}\tau _{k}^{}(1)^j|\omega _n/\omega _{n+j}|}{_{j^{}=0}^k\lambda {}_{n,j^{}}{}^{(k)}\tau _{k}^{}(1)^j^{}\left|\frac{\omega _n}{\omega _{n+j^{}}}\right|}=\frac{|\lambda {}_{n,j}{}^{(k)}||\omega _n/\omega _{n+j}|}{_{j^{}=0}^k|\lambda {}_{n,j^{}}{}^{(k)}||\frac{\omega _n}{\omega _{n+j^{}}}|}0.$$
(315)
Note that the denominators cannot vanish and are bounded from below by $`|\lambda {}_{n,k}{}^{(k)}\omega _{n}^{}/\omega _{n+k}|>0`$. Hence, we have $`\gamma {}_{n,j}{}^{(k)}(\stackrel{}{\omega }_n)=|\gamma {}_{n,j}{}^{(k)}(\stackrel{}{\omega }_n)|`$ and consequently, $`๐ช{}_{n}{}^{(k)}(๐ฏ)=1`$ since $`_{j=0}^k\gamma {}_{n,j}{}^{(k)}(\stackrel{}{\omega }_n)=1`$ according to Eq. (75). $`\mathrm{}`$
### 8.2 Results for Special Cases
Here, we collect some special results on the stability of various Levin-type sequence transformations that have been reported in and generalize some results of Sidi on the S-stability of the $`d^{(1)}`$ transformation.
###### Theorem 18
If the sequence $`\omega _{n+1}/\omega _n`$ possesses a limit according to
$$\underset{n\mathrm{}}{lim}\omega _{n+1}/\omega _n=\rho 0$$
(316)
and if $`\rho \{1,\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_k,\mathrm{}\}`$ such that the limiting transformation exists, the $`๐ฅ`$ transformation is S-stable with the same limiting stability indices as the transformation $`\stackrel{}{๐ฅ}`$, i.e., we have
$$\underset{n\mathrm{}}{lim}๐ช{}_{n}{}^{(k)}=\frac{_{j=0}^k|\lambda {}_{j}{}^{(k)}\rho _{}^{kj}|}{_{j^{}=0}^{k1}|\mathrm{\Phi }_j^{}\rho |}<\mathrm{}.$$
(317)
If all $`\mathrm{\Phi }_k`$ are positive then
$$\underset{n\mathrm{}}{lim}๐ช{}_{n}{}^{(k)}=\underset{j=0}{\overset{k1}{}}\frac{\mathrm{\Phi }_j+|\rho |}{|\mathrm{\Phi }_j\rho |}<\mathrm{}$$
(318)
holds.
As corollaries, we get the following results
###### Corollary 6
If the sequence $`\omega _{n+1}/\omega _n`$ possesses a limit according to
$$\underset{n\mathrm{}}{lim}\omega _{n+1}/\omega _n=\rho \{0,1\}$$
(319)
the $`{}_{p}{}^{}๐`$ transformation for $`p>1`$ and $`\beta >0`$ is S-stable and we have
$$\underset{n\mathrm{}}{lim}๐ช{}_{n}{}^{(k)}=\frac{_{j=0}^k\left(\genfrac{}{}{0pt}{}{k}{j}\right)|\rho ^{kj}|}{|1\rho |^k}=\frac{(1+|\rho |)^k}{|1\rho |^k}<\mathrm{}.$$
(320)
###### Corollary 7
If the sequence $`\omega _{n+1}/\omega _n`$ possesses a limit according to
$$\underset{n\mathrm{}}{lim}\omega _{n+1}/\omega _n=\rho \{0,1\}$$
(321)
the Weniger $`๐ฎ`$ transformation \[84, Sec. 8\] for $`\beta >0`$ is S-stable and we have
$$\underset{n\mathrm{}}{lim}๐ช{}_{n}{}^{(k)}(๐ฎ)=\frac{_{j=0}^k\left(\genfrac{}{}{0pt}{}{k}{j}\right)|\rho ^{kj}|}{|1\rho |^k}=\frac{(1+|\rho |)^k}{|1\rho |^k}<\mathrm{}.$$
(322)
###### Corollary 8
If the sequence $`\omega _{n+1}/\omega _n`$ possesses a limit according to
$$\underset{n\mathrm{}}{lim}\omega _{n+1}/\omega _n=\rho \{0,1\}$$
(323)
the Levin $``$ transformation is S-stable and we have
$$\underset{n\mathrm{}}{lim}๐ช{}_{n}{}^{(k)}()=\frac{_{j=0}^k\left(\genfrac{}{}{0pt}{}{k}{j}\right)|\rho ^{kj}|}{|1\rho |^k}=\frac{(1+|\rho |)^k}{|1\rho |^k}<\mathrm{}.$$
(324)
###### Corollary 9
Assume that the elements of the sequence $`\{t_n\}_n`$ satisfy $`t_n0`$ for all $`n`$ and $`t_nt_n^{}`$ for all $`nn^{}`$. If the sequence $`t_{n+1}/t_n`$ possesses a limit
$$\underset{n\mathrm{}}{lim}t_{n+1}/t_n=\tau \text{ with }0<\tau <1$$
(325)
and if the sequence $`\omega _{n+1}/\omega _n`$ possesses a limit according to
$$\underset{n\mathrm{}}{lim}\omega _{n+1}/\omega _n=\rho \{0,1,\tau ^1,\mathrm{},\tau ^k,\mathrm{}\}$$
(326)
then the generalized Richardson extrapolation process $``$ introduced by Sidi that is identical to the $`๐ฅ`$ transformation with $`\delta {}_{n}{}^{(k)}=t_nt_{n+k+1}`$ as shown in , i.e., the $`W`$ algorithm is S-stable and we have
$$\underset{n\mathrm{}}{lim}๐ช{}_{n}{}^{(k)}()=\frac{_{j=0}^k|\stackrel{~}{\lambda }{}_{j}{}^{(k)}\rho _{}^{kj}|}{_{j^{}=0}^{k1}|\tau ^j^{}\rho |}=\underset{j^{}=0}{\overset{k1}{}}\frac{1+\tau ^j^{}|\rho |}{|1\tau ^j^{}\rho |}<\mathrm{}.$$
(327)
Here,
$$\stackrel{~}{\lambda }{}_{j}{}^{(k)}=(1)^{kj}\underset{\genfrac{}{}{0pt}{}{j_0+j_1+\mathrm{}+j_{k1}=j,}{j_0\{0,1\},\mathrm{},j_{k1}\{0,1\}}}{}\underset{m=0}{\overset{k1}{}}(\tau )^{mj_m}$$
(328)
such that
$$\underset{j=0}{\overset{k}{}}\stackrel{~}{\lambda }{}_{j}{}^{(k)}\rho _{}^{kj}=\underset{j=0}{\overset{k1}{}}(\tau ^j\rho )=\tau ^{k(k1)/2}\underset{j=0}{\overset{k1}{}}(1\tau ^j\rho ).$$
(329)
Note that the preceeding corollary is essentially the same as a result of Sidi \[78, Theorem 2.2\] that now appears as a special case of the more general Theorem 18 that applies to a much wider class of sequence transformations. As noted above, Sidi has also derived conditions under which the $`d^{(1)}`$ transformation is stable along the paths $`๐ซ_n=\{(n,k)|k=0,1,\mathrm{}\}`$ for fixed $`n`$. For details and more references see . Analogous work for the $`๐ฅ`$ transformation is in progress.
An efficient algorithm for the computation of the stability index of the $`๐ฅ`$ transformation can be given in the case $`\delta {}_{n}{}^{(k)}>0`$. Since the $`๐ฅ`$ transformation is invariant under $`\delta {}_{n}{}^{(k)}\alpha ^{(k)}\delta _n^{(k)}`$ for any $`\alpha ^{(k)}0`$ according to \[36, Theorem 4\], $`\delta {}_{n}{}^{(k)}>0`$ can always be achieved if for given $`k`$, all $`\delta _n^{(k)}`$ have the same sign. This is the case, for instance, for the $`{}_{p}{}^{}๐`$ transformation .
###### Theorem 19
Define
$$F{}_{n}{}^{(0)}=(1)^n|D{}_{n}{}^{(0)}|,F{}_{n}{}^{(k+1)}=(F{}_{n+1}{}^{(k)}F{}_{n}{}^{(k)})/\delta _n^{(k)}$$
(330)
and $`\widehat{F}{}_{n}{}^{(0)}=F_n^{(0)}`$, $`\widehat{F}{}_{n}{}^{(k)}=(\delta {}_{n}{}^{(0)}\mathrm{}\delta {}_{n}{}^{(k1)})F_n^{(k)}`$. If all $`\delta {}_{n}{}^{(k)}>0`$ then
1. $`F{}_{n}{}^{(k)}=(1)^{n+k}|F{}_{n}{}^{(k)}|`$,
2. $`\lambda {}_{n,j}{}^{(k)}=(1)^{j+k}|\lambda {}_{n,j}{}^{(k)}|`$, and
3. $$๐ช{}_{n}{}^{(k)}=\frac{|\widehat{F}{}_{n}{}^{(k)}|}{|\widehat{D}{}_{n}{}^{(k)}|}=\frac{|F{}_{n}{}^{(k)}|}{|D{}_{n}{}^{(k)}|}.$$
(331)
This generalizes Sidiโs method for the computation of stability indices to a larger class of sequence transformations.
## 9 Application of Levin-type Sequence Transformations
### 9.1 Practical Guidelines
Here, we address shortly the following questions:
One can only hope for good convergence acceleration, extrapolation, or summation results if a) the $`s_n`$ have some asymptotic structure for large $`n`$ and are not erratic or random, b) a sufficiently large number of decimal digits is available. Many problems can be successfully tackled if 13-15 digits are available but some require a much larger number of digits in order to overcome some inevitable rounding errors, especially for the acceleration of logarithmically convergent sequences. The asymptotic information that is required for a successful extrapolation is often hidden in the last digits of the problem data.
The recommended mode of application is that one computes the highest possible order $`k`$ of the transformation from the data. In the case of triangular recursive schemes like that of the $`๐ฅ`$ transformation and the Levin transformation, this means that one computes as transformed sequence $`\{\{๐ฏ{}_{0}{}^{(n)}\}\}`$. For L-shaped recursive schemes as in the case of the $``$, $``$, and $`๐ฆ`$ transformations, one usually computes as transformed sequence $`\{\{๐ฏ{}_{n2[[n/2]]}{}^{([[n/2]])}\}\}`$. The error $`ฯต`$ of the current estimate can usually be approximated a posteriori using sums of magnitudes of differences of a few entries of the $`๐ฏ`$ table, e.g.,
$$ฯต|๐ฏ{}_{1}{}^{(n)}๐ฏ{}_{0}{}^{(n)}|+|๐ฏ{}_{0}{}^{(n1)}๐ฏ{}_{0}{}^{(n)}|$$
(332)
for transformations with triangular recursive schemes. Such a simple approach works surprisingly well in practice. The loss of decimal digits can be estimated computing stability indices. An example is given below.
The occurrence of zeroes in the $`D`$ table for specific combinations of $`n`$ and $`k`$ is usually no problem since the recurrences for numerators and denominators still work in this case. Thus, no special devices are required to jump over such singular points in the $`๐ฏ`$ table.
This depends on the type of convergence of the problem sequence. For linearly convergent sequences, $`t`$, $`\stackrel{~}{t}`$, $`u`$ and $`v`$ variants of the Levin transformation, the $`{}_{p}{}^{}๐`$ transformation, especially the $`{}_{2}{}^{}๐`$ transformation are usually a good choice as long as one is not too close to a singularity or to a logarithmically convergent problem. Especially well-behaved is usually the application to alternating series since then, the stability is very good as discussed above. For the summation of alternating divergent sequences and series, usually the $`t`$ and the $`\stackrel{~}{t}`$ variants of the Levin transformation, the $`{}_{2}{}^{}๐`$ and the Weniger $`๐ฎ`$ and $``$ transformations provide often surprisingly accurate results. In the case of logarithmic convergence, $`t`$ and $`\stackrel{~}{t}`$ variants become useless, and the order of acceleration is dropping from $`2k`$ to $`k`$ when the transformation is used columnwise. If a Kummer-related series is available (cp. Section 2.2.1), then $`K`$ and $`lu`$ variants leading to linear sequence transformations can be efficient . Similarly, linear variants can be based on some good asymptotic estimates $`{}_{}{}^{asy}\omega _{n}^{}`$, that have to be obtained via a separate analysis . In the case of logarithmcic convergence, it pays to consider special devices like using subsequences $`\{\{s_{\xi _n}\}\}`$ where the $`\xi _n`$ grow exponentially like $`\xi _n=[[\sigma \xi _{n1}]]+1`$ like in the $`d`$ transformations. This choice can be also used in combination with the $``$ transformation. Alternatively, one can use some other transformations like the condensation transformation or interpolation to generate a linearly convergent sequence , before applying an usually nonlinear sequence transformation. A somewhat different approach is possible if one can obtain a few terms $`a_n`$ with large $`n`$ easily .
When extrapolating power series or, more generally, sequences depending on certain parameters, quite often extrapolation becomes difficult near the singularities of the limit function. In the case of linear convergence, one can often transform to a problem with a larger distance to the singularity: If Eq. (28) holds, then the subsequence $`\{\{s_{\tau n}\}\}`$ satisfies
$$\underset{n\mathrm{}}{lim}(s_{\tau (n+1)}s)/(s_{\tau n}s)=\rho ^\tau $$
(333)
This is a method of Sidi that has can, however, be applied to large classes of sequence transformations .
Here, one should try to rewrite the problem sequence as a sum of sequences with more simple convergence behavior. Then, nonlinear sequence transformations are used to extrapolate each of these simpler series, and to sum the extrapolation results to obtain an estimate for the original problem. This is for instance often possible for (generalized) Fourier series where it leads to complex series that may asymptotically be regarded as power series. For details, the reader is referred to the literature . If this approach is not possible one is forced to use more complicated sequence transformations like the $`d^{(m)}`$ transformations or the (generalized) $``$ transformation. These more complicated sequence transformations, however, do require more numerical effort to achieve a desired accuracy.
### 9.2 Numerical Examples
In Table 3, we present results of the application of certain variants of the $``$ transformation and the $`W`$ algorithm to the series
$$S(z,a)=1+\underset{j=1}{\overset{\mathrm{}}{}}z^j\underset{\mathrm{}=0}{\overset{j1}{}}\frac{1}{\mathrm{ln}(a+\mathrm{})}$$
(334)
with partial sums
$$S_n(z,a)=1+\underset{j=1}{\overset{n}{}}z^j\underset{\mathrm{}=0}{\overset{j1}{}}\frac{1}{\mathrm{ln}(a+\mathrm{})}$$
(335)
for $`z=1.2`$ and $`a=1.01`$. Since the terms $`a_j`$ satisfy $`a_{j+1}/a_j=z/\mathrm{ln}(a+j)`$, the ratio test reveals that $`S(z,a)`$ converges for all $`z`$ and, hence, represents an analytic function. Nevertheless, only for $`ja+\mathrm{exp}(|z|)`$, the ratio of the terms becomes less than unity in absolute value. Hence, for larger $`z`$ the series converges rather slowly.
It should be noted that for cases $`C_n`$ and $`G_n`$, the $``$ transformation is identical to the Weniger transformation $`๐ฎ`$, i.e., to the $`{}_{3}{}^{}๐`$ transformation, and for cases $`C_n`$ and $`H_n`$ the $`W`$ algorithm is identical to the Levin transformation. In the upper part of the table, we use u-type remainder estimates while in the lower part, we use $`\stackrel{~}{t}`$ variants. It is seen that the choices $`x_n=1+\mathrm{ln}(a+n)`$ for the $``$ transformation and $`t_n=1/(1+\mathrm{ln}(a+n))`$ for the $`W`$ algorithm perform for both variants nearly identical (columns $`A_n`$, $`B_n`$, $`E_n`$ and $`F_n`$) and are superior to the choices $`x_n=1+n+a`$ and $`t_n=1/(1+n+a)`$, respectively, that correspond to the Weniger and the Levin transformation as noted above. For the latter two transformations, the Weniger $`{}_{}{}^{\stackrel{~}{t}}๐ฎ`$ transformation is slightly superior the $`{}_{}{}^{\stackrel{~}{t}}`$ transformation for this particular example (columns $`G_n`$ vs. $`H_n`$) while the situation is reversed for the u-type variants displayed in colums $`C_n`$ and $`D_n`$.
The next example is taken from , namely the โinflated Riemann $`\zeta `$ functionโ, i.e., the series
$$\zeta (ฯต,1,q)=\underset{j=0}{\overset{\mathrm{}}{}}\frac{q^j}{(j+1)^ฯต}$$
(336)
that is a special case of the Lerch zeta function $`\zeta (s,b,z)`$ (cf. \[30, p. 142, Eq. (6.9.7)\] and \[20, Sec. 1.11\]). The partial sums are defined as
$$s_n=\underset{j=0}{\overset{n}{}}\frac{q^j}{(j+1)^ฯต}.$$
(337)
The series converges linearly for $`0<|q|<1`$ for any complex $`ฯต`$. In fact, we have in this case $`\rho =lim_n\mathrm{}(s_{n+1}s)/(s_ns)=q`$. We choose $`q=0.95`$ and $`ฯต=0.1+10i`$. Note that for this value of $`ฯต`$, there is a singularity of $`\zeta (ฯต,1,q)`$ at $`q=1`$ where the defining series diverges since $`\mathrm{}(ฯต)<1`$.
The results of applying $`u`$ variants of the $`{}_{p}{}^{}๐`$ transformation with $`p=1,2,3`$ and of the Levin transformation to the sequence of partial sums is displayed in Table 4. For each of these four variants of the $`๐ฅ`$ transformation, we give the relative error and the stability index. The true value of the series (that is used to compute the errors) was computed using a more accurate method described below. It is seen that the $`{}_{2}{}^{}๐`$ transformation achieves the best results. The attainable accuracy for this transformation is limited to about 9 decimal digits by the fact that the stability index displayed in the column $`D_n`$ of Table 4 grows relatively fast. Note that for $`n=46`$, the number of digits (as given by the negative decadic logarithm of the relative error) and the decadic logarithm of the stability index sum up to approximately 32 which corresponds to the maximal number of decimal digits that could be achieved in the run. Since the stability index increases with $`n`$, indicating decreasing stability, it is clear that for higher values of $`n`$ the accuracy will be lower.
The magnitude of the stability index is largely controlled by the value of $`\rho `$, compare Corollary 6. If one can treat a related sequence with a smaller value of $`\rho `$, the stability index will be smaller and thus, the stability of the extrapolation will be greater.
Such a related sequence is given by putting $`\stackrel{ห}{s}_{\mathrm{}}=s_\xi _{\mathrm{}}`$ for $`\mathrm{}_0`$, where the sequence $`\xi _{\mathrm{}}`$ is a monotonously increasing sequence of nonnegative integers. In the case of linear convergent sequences, the choice $`\xi _{\mathrm{}}=\tau \mathrm{}`$ with $`\tau `$ can be used as in the case of the $`d^{(1)}`$ transformation. It is easily seen that the new sequence also converges linearly with $`\rho =lim_n\mathrm{}(\stackrel{ห}{s}_{n+1}s)/(\stackrel{ห}{s}_ns)=q^\tau `$. For $`\tau >1`$, both the effectiveness and the stability of the various transformations are increased as shown in Table 5 for the case $`\tau =10`$. Note that this value was chosen to display basic features relevant to the stability analysis, and is not necessarily the optimal value. As in Table 4, the relative errors and the stability indices of some variants of the $`๐ฅ`$ transformation are displayed. These are nothing but the $`{}_{p}{}^{}๐`$ transformation for $`p=1,2,3`$ and the Levin transformation as applied to the sequence $`\{\{\stackrel{ห}{s}_n\}\}`$ with remainder estimates $`\omega _n=(n\tau +\beta )(s_{n\tau }s_{n\tau 1})`$ for $`\beta =1`$. Since constant factors in the remainder estimates are irrelevant since the $`๐ฅ`$ transformation is invariant under any scaling $`\omega _n\alpha \omega _n`$ for $`\alpha 0`$, the same results would have been obtained for $`\omega _n=(n+\beta /\tau )(s_{n\tau }s_{n\tau 1})`$.
If the Levin transformation is applied to the series with partial sums $`\stackrel{ห}{s}_n=s_{\tau n}`$, and if the remainder estimates $`\omega _n=(n+\beta /\tau )(s_{\tau n}s_{(\tau n)1})`$ are used, then one obtains nothing but the $`d^{(1)}`$ transformation with $`\xi _{\mathrm{}}=\tau \mathrm{}`$ for $`\tau `$.
It is seen from Table 5 that again the best accuracy is obtained for the $`{}_{2}{}^{}๐`$ transformation. The $`d^{(1)}`$ transformation is worse, but better than the $`{}_{p}{}^{}๐`$ transformations for $`p=1`$ and $`p=3`$. Note that the stability indices are now much smaller and do not limit the achievable accuracy for any of the transformations up to $`n=30`$. The true value of the series was computed numerically by applying the $`{}_{2}{}^{}๐`$ transformation to the further sequence $`\{\{s_{40n}\}\}`$ and using 64 decimal digits in the calculation. In this way, a sufficiently accurate approximation was obtained that was used to compute the relative errors in Tables 4 and 5. A comparison value was computed using the representation \[20, p. 29, Eq. (8)\]
$$\zeta (s,1,q)=\frac{\mathrm{\Gamma }(1s)}{z}(\mathrm{log}1/q)^{s1}+z^1\underset{j=0}{\overset{\mathrm{}}{}}\zeta (sj)\frac{(\mathrm{log}q)^j}{j!}$$
(338)
that holds for $`|\mathrm{log}q|<2\pi `$ and $`s`$. Here, $`\zeta (z)`$ denotes the Riemann zeta function. Both values agreed to all relevant decimal digits.
In Table 6, we display stability indices corresponding to the acceleration of $`\stackrel{ห}{s}_n`$ with the $`{}_{2}{}^{}๐`$ transformation columnwise, as obtainable by using the sequence elements up to $`\stackrel{ห}{s}_{50}=s_{500}`$. In the row labelled *Cor. 6*, we display the limits of the $`๐ช_n^{(k)}`$ for large $`n`$, i.e., the quantities
$$\underset{n\mathrm{}}{lim}๐ช{}_{n}{}^{(k)}=\left(\frac{1+q^\tau }{1q^\tau }\right)^k$$
(339)
that are the limits according to Corollary 6. It is seen that the values for finite $`n`$ are still relatively far off the limits. In order to check numerically the validity of the corollary, we extrapolated the values of all $`๐ช_n^{(k)}`$ for fixed $`k`$ with $`n`$ up to the maximal $`n`$ for which there is an entry in the corresponding column of Table 6 using the $`u`$ variant of the $`{}_{1}{}^{}๐`$ transformation. The results of the extrapolation are displayed in the row labelled *Extr* in Table 6 and coincide nearly perfectly with the values expected according to Corollary 6.
As a final example, we consider the evaluation of the $`F_m(z)`$ functions that are used in quantum chemistry calculations via the series representation
$$F_m(z)=\underset{j=0}{\overset{\mathrm{}}{}}(z)^j/j!(2m+2j+1),$$
(340)
with partial sums
$$s_n=\underset{j=0}{\overset{n}{}}(z)^j/j!(2m+2j+1).$$
(341)
In this case, for larger $`z`$, the convergence is rather slow although the convergence finally is hyperlinear. As a $`K`$ variant, one may use
$${}_{}{}^{K}\omega _{n}^{}=\left(\underset{j=0}{\overset{n}{}}(z)^j/(j+1)!(1e^z)/z\right).$$
(342)
since $`(1e^z)/z`$ is a Kummer related series. The results for several variants in Table 7 show that the $`K`$ variant is superior to $`u`$ and $`t`$ variants in this case.
Many further numerical examples are given in the literature .
## Appendix A Stieltjes Series and Functions
A Stieltjes series is a formal expansion
$$f(z)=\underset{j=0}{\overset{\mathrm{}}{}}(1)^j\mu _jz^j$$
(343)
with partial sums
$$f_n(z)=\underset{j=0}{\overset{n}{}}(1)^j\mu _jz^j.$$
(344)
The coefficients $`\mu _n`$ are the moments of an uniquely given positive measure $`\psi (t)`$ that has infinitely many different values on $`0t<\mathrm{}`$ \[4, p. 159\]:
$$\mu _n=_0^{\mathrm{}}t^n๐\psi (t),n_0.$$
(345)
Formally, the Stieltjes series can be identified with a Stieltjes integral
$$f(z)=_0^{\mathrm{}}\frac{d\psi (t)}{1+zt},\left|\mathrm{arg}(z)\right|<\pi .$$
(346)
If such an integral exists for a function $`f`$ then the function is called a Stieltjes function. For every Stieltjes function there exist a unique asymptotical Stieltjes series (343), uniformly in every sector $`\left|\mathrm{arg}(z)\right|<\theta `$ for all $`\theta <\pi `$. For any Stieltjes series, however, several different corresponding Stieltjes functions may exist. To ensure uniqueness, additional criteria are necessary \[88, Sec. 4.3\].
In the context of convergence acceleration and summation of divergent series, it is important that for given $`z`$ the tails $`f(z)f_n(z)`$ of a Stieltjes series are bounded in absolute value by the by the next term of the series,
$$|f(z)f_n(z)|\mu _{n+1}z^{n+1}z0.$$
(347)
Hence, for Stieltjes series the remainder estimates may be chosen as
$$\omega _n=(1)^{n+1}\mu _{n+1}z^{n+1}$$
(348)
This corresponds to $`\omega _n=\mathrm{\Delta }f_n(z)`$, i.e., to a $`\stackrel{~}{t}`$ variant.
## Appendix B Derivation of the Recursive Scheme (148)
We show that for the divided difference operator $`\mathrm{}{}_{n}{}^{(k)}=\mathrm{}{}_{n}{}^{(k)}[\{\{x_n\}\}]`$ the identity
$$\mathrm{}{}_{n}{}^{(k+1)}((x)_{\mathrm{}+1}g(x))=\frac{(x_{n+k+1}+\mathrm{})\mathrm{}{}_{n+1}{}^{(k)}((x)_{\mathrm{}}g(x))(x_n+\mathrm{})\mathrm{}{}_{n}{}^{(k)}((x)_{\mathrm{}}g(x))}{x_{n+k+1}x_n}$$
(349)
holds. The proof is based on the Leibniz formula for divided differences (see, e.g., \[69, p. 50\]) that yields upon use of $`(x)_{\mathrm{}+1}=(x+\mathrm{})(x)_{\mathrm{}}`$ and $`\mathrm{}{}_{n}{}^{(k)}(x)=x_n\delta _{k,0}+\delta _{k,1}`$
$$\begin{array}{cc}\hfill \mathrm{}{}_{n}{}^{(k+1)}((x)_{\mathrm{}+1}g(x))& =\mathrm{}\mathrm{}{}_{n}{}^{(k+1)}((x)_{\mathrm{}}g(x))+\underset{j=0}{\overset{k+1}{}}\mathrm{}{}_{n}{}^{(j)}(x)\mathrm{}{}_{n+j}{}^{(k+1j)}((x)_{\mathrm{}}g(x))\hfill \\ & =(x_n+\mathrm{})\mathrm{}{}_{n}{}^{(k+1)}((x)_{\mathrm{}}g(x))+\mathrm{}{}_{n+1}{}^{(k)}((x)_{\mathrm{}}g(x))\hfill \end{array}$$
(350)
Using the recursion relation of the divided differences, one obtains
$$\mathrm{}{}_{n}{}^{(k+1)}((x)_{\mathrm{}+1}g(x))=(x_n+\mathrm{})\frac{\mathrm{}{}_{n+1}{}^{(k)}((x)_{\mathrm{}}g(x))\mathrm{}{}_{n}{}^{(k)}((x)_{\mathrm{}}g(x))}{x_{n+k+1}x_n}+\mathrm{}{}_{n+1}{}^{(k)}((x)_{\mathrm{}}g(x)).$$
(351)
Simple algebra then yields Eq. (349).
Comparison with Eq. (140) shows that using the interpolation conditions $`g_n=g(x_n)=s_n/\omega _n`$ and $`\mathrm{}=k1`$ in Eq. (349) yields the recursion for the numerators in Eq. (148), while the recursion for the denominators in Eq. (148) follows for $`\mathrm{}=k1`$ and using the interpolation conditions $`g_n=g(x_n)=1/\omega _n`$. In each case, the initial conditions follow directly from Eq. (140) in combination with the definition of the divided difference operator: For $`k=0`$, we use $`(a)_1=1/(a1)`$ and obtain $`\mathrm{}{}_{n}{}^{(k)}(x_n)_{k1}^{}g_n=(x_n)_1g_n=g_n/(x_n1)`$.
## Appendix C Two Lemmata
###### Lemma 1
Define
$$A=\underset{j=0}{\overset{k}{}}\stackrel{}{๐}{}_{j}{}^{(k)}\frac{\zeta ^{n+j}}{(n+j)^{r+1}}$$
(352)
where $`\zeta `$ is a zero of of multiplicity $`m`$ of $`\stackrel{}{\Pi }{}_{}{}^{(k)}(z)=_{j=0}^k\stackrel{}{๐}{}_{j}{}^{(k)}z_{}^{j}`$. Then
$$A\zeta ^{n+m}\left(\genfrac{}{}{0pt}{}{r+m}{r}\right)\frac{(1)^m}{n^{r+m+1}}\frac{d^m\stackrel{}{\Pi }^{(k)}}{dx^m}(\zeta ),(n\mathrm{}).$$
(353)
PROOF. Use
$$\frac{1}{a^{r+1}}=\frac{1}{r!}_0^{\mathrm{}}\mathrm{exp}(at)t^r๐t,a>0$$
(354)
to obtain
$$A=\frac{1}{r!}_0^{\mathrm{}}\underset{j=0}{\overset{k}{}}\stackrel{}{๐}{}_{j}{}^{(k)}\zeta _{}^{n+j}\mathrm{exp}((n+j)t)t^rdt=\frac{\zeta ^n}{r!}_0^{\mathrm{}}\mathrm{exp}(nt)\stackrel{}{\Pi }{}_{}{}^{(k)}(\zeta \mathrm{exp}(t))t^rdt.$$
(355)
Taylor expansion of the polynomial yields due to the zero at $`\zeta `$
$$\stackrel{}{\Pi }{}_{}{}^{(k)}(\zeta \mathrm{exp}(t))=\frac{(\zeta )^m}{m!}\frac{d^m\stackrel{}{\Pi }{}_{}{}^{(k)}(x)}{dx^m}|_{x=\zeta }t^m(1+O(t)).$$
(356)
Invoking Watsonโs lemma \[6, p. 263ff\] completes the proof. $`\mathrm{}`$
###### Lemma 2
Assume that assumption (C-3โ) of Theorem 11 holds. Further assume $`\lambda {}_{n,j}{}^{(k)}\stackrel{}{๐}_j^{(k)}`$ for $`n\mathrm{}`$. Then, Eq. (280) holds.
PROOF. We have
$$\frac{\omega _{n+j}}{\omega _n}\rho ^j\mathrm{exp}\left(ฯต_n\underset{t=0}{\overset{j1}{}}\frac{ฯต_{n+t}}{ฯต_n}\right)\rho ^j\mathrm{exp}(jฯต_n)$$
(357)
for large $`n`$. Hence,
$$\underset{j=0}{\overset{k}{}}\lambda {}_{n,j}{}^{(k)}\frac{\omega _n}{\omega _{n+j}}\underset{j=0}{\overset{k}{}}\stackrel{}{๐}{}_{j}{}^{(k)}(\rho \mathrm{exp}(ฯต_n))_{}^{j}=\stackrel{}{\Pi }{}_{}{}^{(k)}(1/\rho +\delta _n)$$
(358)
Since the characteristic polynomial $`\stackrel{}{\Pi }{}_{}{}^{(k)}(z)`$ has a zero of order $`\mu `$ at $`z=1/\rho `$ according to the assumptions, Eq. (280) follows using Taylor expansion. $`\mathrm{}`$
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# One-body density matrix and momentum distribution in ๐ -๐ and ๐ -๐ shell nuclei
## I INTRODUCTION
The momentum distribution (MD) is of interest in many research subjects of modern physics, including those referring to helium, electronic, nuclear, and quark systems . In the last two decades, there has been significant effort for the determination of the MD in nuclear matter and finite nucleon systems . MD is related to the cross sections of various kinds of nuclear reactions. Specifically, the interaction of particles with nuclei at high energies, such as (p,2p), (e,ep), and (e,e) reactions, the nuclear photo-effect, meson absorption by nuclei, the inclusive proton production in proton-nucleus collisions, and even phenomena at low energies such as giant multipole resonances, give significant information about the nucleon MD. The experimental evidence obtained from inclusive and exclusive electron scattering on nuclei established the existence of a high-momentum component for momenta $`k>2\mathrm{fm}^1`$ . It has been shown that, in principle, mean field theories can not describe correctly MD and density distribution simultaneously and the main features of MD depend little on the effective mean field considered . The reason is that MD is sensitive to short-range and tensor nucleon-nucleon correlations which are not included in the mean field theories. Thus, theoretical approaches, which take into account short range correlations (SRC) due to the character of the nucleon-nucleon forces at small distances, are necessary to be developed.
Zabolitzky and Ey , employing the coupled-cluster (or $`\mathrm{exp}(S)`$) method for the microscopic evaluation of nuclear MD for the ground states of <sup>4</sup>He and <sup>16</sup>O and using various realistic NN-potentials, showed that the contribution of correlations dominates for momenta beyond $`2\mathrm{fm}^1`$. A realistic interaction and a many-body approach have been used by Benhar et al for the evaluation of MD of <sup>12</sup>C, <sup>16</sup>O and <sup>40</sup>Ca. Their results have yielded a much larger content of the high momentum component with respect to the results obtained within the Hartree-Fock approach or within methods which take into account the effect of correlations phenomenologically.
Bohigas and Stringari and Dal Ri et al evaluated the effect of SRCโs on the one- and two- body densities by developing a low order approximation (LOA) in the framework of Jastrow formalism. They showed that one-body quantities provide an adequate test for the presence of SRCโs in nuclei, which indicates that the independent-particle wave functions cannot reproduce simultaneously the form factor and the MD of a correlated system and also the effect of SRCโs strongly modify the MD by introducing an important contribution in the region $`k>2`$ fm<sup>-1</sup>. Stoitsov et al generalised the model of Jastrow correlations within the LOA of Ref. , to heavier nuclei as <sup>16</sup>O, <sup>36</sup>Ar, <sup>40</sup>Ca. Their analytical expressions for the MD show the high momentum tail. They found that there is an A dependence of MD for small values of $`k`$, while for large values of $`k`$ the slope of $`\mathrm{log}n(k)`$ versus $`k`$ is roughly the same for the above three nuclei as well as for <sup>4</sup>He. The same behaviour of the MD of protons and neutrons for the $`A=3,4`$ nuclei has been found earlier by Schiavilla et al performing variational calculations with realistic interactions. MD for the nuclei <sup>4</sup>He, <sup>16</sup>O and <sup>40</sup>Ca was also calculated by Traini and Orlandini within a phenomenological model in which dynamical short-range and tensor correlations effects were included. They showed that SRC increase the high momentum component considerably while the tensor correlations do not affect the MD appreciably . In heavy nuclei, the local density approximation was used for the study of the effect of SRCโs in MD and the predictions were in agreement with the results of microscopic calculations in nuclear matter and in light nuclei.
The influence of SRCโs on the MD of nucleons in nuclei has also been evaluated by Mรผther et al within the Green-function approach assuming a realistic meson-exchange potential for the nucleon-nucleon interaction. Their analysis on <sup>16</sup>O demonstrates that a non-negligible contribution to the MD should be found in partial waves which are unoccupied in the simple shell model. Another approach is to consider the average occupancy of the relevant shell-model orbitals . It has been found that the depletion of such orbitals can be of the order $`15\%`$ or more for single-particles (SP) states below the Fermi energy .
In the various approaches, the MD of the closed shell nuclei <sup>4</sup>He, <sup>16</sup>O and <sup>40</sup>Ca as well as of <sup>208</sup>Pb and nuclear matter is usually studied. There is no systematic study of the one body density matrix (OBDM) and MD which include both the case of closed and open shell nuclei. This would be helpful in the calculations of the overlap integrals and reactions in that region of nuclei if one wants to go beyond the mean field theories . For that reason, in the present work, we attempt to find some general expressions for the OBDM $`\rho (๐ซ,๐ซ^{})`$ and MD $`n(๐ค)`$ which could be used both for closed and open shell nuclei. This work is a continuation of our previous study on the form factors and densities of the $`s`$-$`p`$ and $`s`$-$`d`$ shell nuclei. The expression of $`\rho (๐ซ,๐ซ^{})`$ was found, first, using the factor cluster expansion of Clark and co-workers and Jastrow correlation function which introduces SRC for closed shell nuclei and then was extrapolated to the case of $`N=Z`$ open shell nuclei. $`n(๐ค)`$ was found by Fourier transform of $`\rho (๐ซ,๐ซ^{})`$. These expressions are functionals of the harmonic oscillator (HO) orbitals and depend on the HO parameter $`b`$ and the correlation parameter $`\beta `$. The values of the parameters $`b`$ and $`\beta `$, which we have used for the closed shell nuclei <sup>4</sup>He, <sup>16</sup>O and <sup>40</sup>Ca, are the ones which have been determined in Ref. by fit of the theoretical $`F_{ch}(q)`$, derived with the same cluster expansion, to the experimental one. For the open shell nuclei <sup>12</sup>C, <sup>24</sup>Mg, <sup>28</sup>Si and <sup>32</sup>S we provide new values for these parameters, which have been found to give a better fit to the experimental form factors than in our previous analysis . It is found that the high-momentum tail of the MD of all the nuclei we have considered appears for $`k>2\mathrm{fm}^1`$ and also there is an A dependence of the values of $`n(k)`$ for $`2\mathrm{fm}^1<k<5\mathrm{fm}^1`$. This $`A`$ dependence of MD was first investigated considering <sup>24</sup>Mg, <sup>28</sup>Si and <sup>32</sup>S as $`1d`$ shell nuclei. Next we treated the above nuclei as $`1d`$-$`2s`$ shell nuclei having the occupation probability of the $`2s`$ state as an extra free parameter in the fit of the form factors. The $`A`$ dependence is quite small in the second case.
The paper is organised as follows. In Sec. II the general expressions of the correlated OBDM and MD are derived using a Jastrow correlation function. In Sec. III the analytical expressions of the above quantities for the $`s`$-$`p`$ and $`s`$-$`d`$ shell nuclei, in the case of the HO orbitals, are given. Numerical results are reported and discussed in Sec. IV, while the summary of the present work is given in Sec. V.
## II CORRELATED ONE-BODY DENSITY MATRIX AND MOMENTUM DISTRIBUTION
A nucleus with $`A`$ nucleons is described by the wave function $`\mathrm{\Psi }(๐ซ_1,๐ซ_2,\mathrm{},๐ซ_A)`$ which depends on $`3A`$ coordinates as well as on spins and isospins. The evaluation of the single particle characteristics of the system needs the one-body density matrix
$$\rho (๐ซ,๐ซ^{})=\mathrm{\Psi }^{}(๐ซ,๐ซ_2,\mathrm{},๐ซ_A)\mathrm{\Psi }(๐ซ^{},๐ซ_2,\mathrm{},๐ซ_A)d๐ซ_2\mathrm{}d๐ซ_A,$$
(1)
where the integration is carried out over the radius vectors $`๐ซ_2,\mathrm{},๐ซ_A`$ and summation over spin and isospin variables is implied. $`\rho (๐ซ,๐ซ^{})`$ can also be represented by the form
$$\rho (๐ซ,๐ซ^{})=\frac{\mathrm{\Psi }|๐_{\mathrm{๐ซ๐ซ}^{}}|\mathrm{\Psi }^{}}{\mathrm{\Psi }|\mathrm{\Psi }}=N\mathrm{\Psi }|๐_{\mathrm{๐ซ๐ซ}^{}}|\mathrm{\Psi }^{}=N๐_{\mathrm{๐ซ๐ซ}^{}},$$
(2)
where $`\mathrm{\Psi }^{}=\mathrm{\Psi }(๐ซ_1^{},๐ซ_2^{},\mathrm{},๐ซ_A^{})`$ and $`N`$ is the normalization factor. The one-body โdensity operatorโ $`๐_{\mathrm{๐ซ๐ซ}^{}}`$, has the form
$$๐_{\mathrm{๐ซ๐ซ}^{}}=\underset{i=1}{\overset{A}{}}\delta (๐ซ_i๐ซ)\delta (๐ซ_i^{}๐ซ^{})\underset{ji}{\overset{A}{}}\delta (๐ซ_j๐ซ_j^{}).$$
(3)
In the case where the nuclear wave function $`\mathrm{\Psi }`$ can be expressed as a Slater determinant depending on the SP wave functions $`\varphi _i(๐ซ)`$ we have
$$\rho _{SD}(๐ซ,๐ซ^{})=\underset{i=1}{\overset{A}{}}\varphi _i^{}(๐ซ)\varphi _i(๐ซ^{})$$
(4)
The diagonal elements of the OBDM give the density distribution
$$\rho (๐ซ,๐ซ)=\rho (๐ซ),$$
(5)
while the MD is given by the Fourier transform of $`\rho (๐ซ,๐ซ^{})`$,
$$n(๐ค)=\frac{1}{(2\pi )^3}\mathrm{exp}[i๐ค(๐ซ๐ซ^{})]\rho (๐ซ,๐ซ^{})d๐ซd๐ซ^{}.$$
(6)
In the case of a Slater determinant, MD takes the form
$$n_{SD}(๐ค)=\underset{i=1}{\overset{A}{}}\stackrel{~}{\varphi }_i^{}(๐ค)\stackrel{~}{\varphi }_i(๐ค),$$
(7)
where
$$\stackrel{~}{\varphi }_i(๐ค)=\frac{1}{(2\pi )^{3/2}}\varphi _i(๐ซ)\mathrm{exp}[i\mathrm{๐ค๐ซ}]d๐ซ.$$
(8)
The second moment of the MD is related to the expectation value of the kinetic energy, $`๐`$, by the expression
$$๐=\frac{\mathrm{}^2}{2m}n(๐ค)k^2d๐ค.$$
(9)
### A One-body density matrix
If we denote the model operator, which introduces SRC, by $``$, an eigenstate $`\mathrm{\Phi }`$ of the model system corresponds to an eigenstate
$$\mathrm{\Psi }=\mathrm{\Phi }$$
(10)
of the true system.
Several restrictions can be made on the model operator $``$, as for example, that it depends on (the spins, isospins and) relative coordinates and momenta of the particles in the system, that be a scalar with respect to rotations etc. . Further, it is required that $``$ be translationally invariant and symmetrical in its arguments $`1\mathrm{}i\mathrm{}A`$ and possesses the cluster property. That is if any subset $`i_1\mathrm{}i_p`$ of the particles is removed far from the rest $`i_{p+1}\mathrm{}i_A`$, $``$ decomposes into a product of two factors, $`(1\mathrm{}A)=(i_1\mathrm{}i_p)(i_{p+1}\mathrm{}i_A)`$ . In the present work $``$ is taken to be of the Jastrow type ,
$$=\underset{i<j}{\overset{A}{}}f(r_{ij}),$$
(11)
where $`f(r_{ij})`$ is the state-independent correlation function of the form
$$f(r_{ij})=1\mathrm{exp}[\beta (๐ซ_i๐ซ_j)^2].$$
(12)
The correlation function $`f(r_{ij})`$ goes to 1 for large values of $`r_{ij}=๐ซ_i๐ซ_j`$ and it goes to 0 for $`r_{ij}0`$. It is obvious that the effect of SRC, introduced by the function $`f(r_{ij})`$, becomes large when the SRC parameter $`\beta `$ becomes small and vice versa.
In order to evaluate the correlated one-body density matrix $`\rho _{cor}(๐ซ,๐ซ^{})`$, we consider, first, the generalized integral
$$I(\alpha )=\mathrm{\Psi }|\mathrm{exp}[\alpha I(0)๐_{\mathrm{๐ซ๐ซ}^{}}]|\mathrm{\Psi }^{},$$
(13)
corresponding to the one-body โdensity operatorโ $`๐_{\mathrm{๐ซ๐ซ}^{}}`$ (given by (3)), from which we have
$$๐_{\mathrm{๐ซ๐ซ}^{}}=\left[\frac{\mathrm{ln}I(\alpha )}{\alpha }\right]_{\alpha =0}.$$
(14)
For the cluster analysis of equation (14), we consider the sub-product integrals , for the sub-systems of the $`A`$-nucleons system
$`I_i(\alpha )`$ $`=`$ $`i^{}(r_1)\mathrm{exp}[\alpha I_i(0)๐จ_{\mathrm{๐ซ๐ซ}^{}}(1)](r_1^{})i^{},`$ (15)
$`I_{ij}(\alpha )`$ $`=`$ $`ij^{}(r_{12})\mathrm{exp}[\alpha I_{ij}(0)๐จ_{\mathrm{๐ซ๐ซ}^{}}(2)](r_{12}^{})i^{}j^{}_\mathrm{a},`$ (16)
$`I_{ijk}(\alpha )`$ $`=`$ $`ijk^{}(r_{12})^{}(r_{13})^{}(r_{23})\mathrm{exp}[\alpha I_{ijk}(0)๐จ_{\mathrm{๐ซ๐ซ}^{}}(3)](r_{12}^{})(r_{13}^{})(r_{23}^{})i^{}j^{}k^{}_\mathrm{a},`$ (17)
. (18)
. (19)
. (20)
$`I_{12\mathrm{}A}`$ $`=`$ $`I(\alpha ),`$ (21)
where the operators $`๐จ_{\mathrm{๐ซ๐ซ}^{}}(1)`$, $`๐จ_{\mathrm{๐ซ๐ซ}^{}}(2)`$, $`\mathrm{}`$ have the form
$`๐จ_{\mathrm{๐ซ๐ซ}^{}}(1)`$ $`=`$ $`\delta (๐ซ_1๐ซ)\delta (๐ซ_1^{}๐ซ^{}),`$ (22)
$`๐จ_{\mathrm{๐ซ๐ซ}^{}}(2)`$ $`=`$ $`\delta (๐ซ_1๐ซ)\delta (๐ซ_1^{}๐ซ^{})\delta (๐ซ_2๐ซ_2^{})+\delta (๐ซ_2๐ซ)\delta (๐ซ_2^{}๐ซ^{})\delta (๐ซ_1๐ซ_1^{}),`$ (23)
and so on.The factor cluster decomposition of the above integrals, following the factor cluster expansion of Ristig,Ter Low, and Clark , gives
$$๐_{\mathrm{๐ซ๐ซ}^{}}=๐_{\mathrm{๐ซ๐ซ}^{}}_1+๐_{\mathrm{๐ซ๐ซ}^{}}_2+\mathrm{}+๐_{\mathrm{๐ซ๐ซ}^{}}_A,$$
(24)
where
$`๐_{\mathrm{๐ซ๐ซ}^{}}_1`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{A}{}}}\left[{\displaystyle \frac{\mathrm{ln}I_i(\alpha )}{\alpha }}\right]_{\alpha =0},`$ (25)
$`๐_{\mathrm{๐ซ๐ซ}^{}}_2`$ $`=`$ $`{\displaystyle \underset{i<j}{\overset{A}{}}}{\displaystyle \frac{}{\alpha }}\left[\mathrm{ln}I_{ij}(\alpha )\mathrm{ln}I_i(\alpha )\mathrm{ln}I_j(\alpha )\right]_{\alpha =0},`$ (26)
and so on. $`(r_1)`$ is chosen to be the identity operator.
Three- and many-body terms will be neglected in the present analysis. Thus, in the two-body approximation, $`\rho _{cor}(๐ซ,๐ซ^{})`$, defined by Eq. (2), is written
$$\rho _{cor}(๐ซ,๐ซ^{})N[๐_{\mathrm{๐ซ๐ซ}^{}}_1+๐_{\mathrm{๐ซ๐ซ}^{}}_{22}๐_{\mathrm{๐ซ๐ซ}^{}}_{21}],$$
(27)
where
$$๐_{\mathrm{๐ซ๐ซ}^{}}_1=\underset{i=1}{\overset{A}{}}i๐จ_{\mathrm{๐ซ๐ซ}^{}}(1)i^{},$$
(28)
$$๐_{\mathrm{๐ซ๐ซ}^{}}_{22}=\underset{i<j}{\overset{A}{}}ij^{}(r_{12})๐จ_{\mathrm{๐ซ๐ซ}^{}}(2)(r_{12}^{})i^{}j^{}_\mathrm{a},$$
(29)
$$๐_{\mathrm{๐ซ๐ซ}^{}}_{21}=\underset{i<j}{\overset{A}{}}ij๐จ_{\mathrm{๐ซ๐ซ}^{}}(2)i^{}j^{}_\mathrm{a}$$
(30)
If the two-body operator $`(r_{12}^{})`$ is taken to be the correlation function given by Eq. (12), then
$$^{}(r_{12})๐จ_{\mathrm{๐ซ๐ซ}^{}}(2)(r_{12}^{})=๐จ_{\mathrm{๐ซ๐ซ}^{}}(2)\left[1\mathrm{g}_1(๐ซ,๐ซ_2)\mathrm{g}_2(๐ซ^{},๐ซ_2)+\mathrm{g}_3(๐ซ,๐ซ^{},๐ซ_\mathrm{๐})\right],$$
(31)
where
$`\mathrm{g}_1(๐ซ,๐ซ_2)`$ $`=`$ $`\mathrm{exp}[\beta (r^2+r_2^2)]\mathrm{exp}[2\beta \mathrm{๐ซ๐ซ}_2],\mathrm{g}_2(๐ซ^{},๐ซ_2)=\mathrm{g}_1(๐ซ^{},๐ซ_2),`$ (32)
$`\mathrm{g}_3(๐ซ,๐ซ^{},๐ซ_2)`$ $`=`$ $`\mathrm{exp}[\beta (r^2+r_{}^{}{}_{}{}^{2})]\mathrm{exp}[2\beta r_2^2]\mathrm{exp}[2\beta (๐ซ+๐ซ^{})๐ซ_2]`$ (33)
and the term $`๐_{\mathrm{๐ซ๐ซ}^{}}_{22}`$ is written
$$๐_{\mathrm{๐ซ๐ซ}^{}}_{22}=๐_{\mathrm{๐ซ๐ซ}^{}}_{21}O_{22}(๐ซ,๐ซ^{},\mathrm{g}_1)O_{22}(๐ซ,๐ซ^{},\mathrm{g}_2)+O_{22}(๐ซ,๐ซ^{},\mathrm{g}_3),$$
(34)
where
$$O_{22}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})=\underset{i<j}{\overset{A}{}}ij๐จ_{\mathrm{๐ซ๐ซ}^{}}(2)\mathrm{g}_{\mathrm{}}(๐ซ,๐ซ^{},๐ซ_2)i^{}j^{}_\mathrm{a},\mathrm{}=1,2,3.$$
(35)
Thus, $`\rho _{cor}(๐ซ,๐ซ^{})`$ takes the form
$$\rho _{cor}(๐ซ,๐ซ^{})N[๐_{\mathrm{๐ซ๐ซ}^{}}_1O_{22}(๐ซ,๐ซ^{},\mathrm{g}_1)O_{22}(๐ซ,๐ซ^{},\mathrm{g}_2)+O_{22}(๐ซ,๐ซ^{},\mathrm{g}_3)].$$
(36)
This is also expressed in the following form
$`\rho _{cor}(๐ซ,๐ซ^{})`$ $``$ $`N[\rho _{SD}(๐ซ,๐ซ^{})+{\displaystyle }[\mathrm{g}_1(๐ซ,๐ซ_2)\mathrm{g}_2(๐ซ^{},๐ซ_2)+\mathrm{g}_3(๐ซ,๐ซ^{},๐ซ_2)]`$ (39)
$`\times [\rho _{SD}(๐ซ,๐ซ^{})\rho _{SD}(๐ซ_2,๐ซ_2)\rho _{SD}(๐ซ,๐ซ_2)\rho _{SD}(๐ซ_2,๐ซ^{})]\mathrm{d}๐ซ_2{\displaystyle \frac{}{}}],`$
where $`\rho _{SD}(๐ซ,๐ซ^{})`$ is the uncorrelated OBDM associated with the Slater determinant.
It should be noted that a similar expression for $`\rho _{cor}(๐ซ,๐ซ^{})`$, given by Eq. (39), was derived by Gaudin et al. in the framework of LOA. Their expansion contains one- and two-body terms and a part of the three-body term which was chosen so that the normalization of the wave function was preserved. Expression (39) of the present work has only one- and two-body terms and the normalization of the wave function is preserved by the normalization factor $`N`$.
In the above expression of $`\rho _{cor}(๐ซ,๐ซ^{})`$, the one-body contribution to the OBDM is well known and is given by the equation
$$๐_{\mathrm{๐ซ๐ซ}^{}}_1=\rho _{SD}(๐ซ,๐ซ^{})=\frac{1}{\pi }\underset{nl}{}\eta _{nl}(2l+1)\varphi _{nl}^{}(r)\varphi _{nl}(r^{})P_l(\mathrm{cos}\omega _{rr^{}})$$
(40)
where $`\eta _{nl}`$ are the occupation probabilities of the states $`nl`$ (0 or 1 in the case of closed shell nuclei) and $`\varphi _{nl}(r)`$ is the radial part of the SP wave function and $`\omega _{rr^{}}`$ the angle between the vectors $`๐ซ`$ and $`๐ซ^{}`$.
The term $`O_{22}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})`$, performing the spin-isospin summation and the angular integration, takes the general form
$`O_{22}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})`$ $`=`$ $`4{\displaystyle \underset{n_il_i,n_jl_j}{}}\eta _{n_il_i}\eta _{n_jl_j}(2l_i+1)(2l_j+1)`$ (42)
$`\times [4A_{n_il_in_jl_j}^{n_il_in_jl_j,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}}){\displaystyle \underset{k=0}{\overset{l_i+l_j}{}}}l_i0l_j0k0^2A_{n_il_in_jl_j}^{n_jl_jn_il_i,k}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})],\mathrm{}=1,2,3,`$
where
$`A_{n_1l_1n_2l_2}^{n_3l_3n_4l_4,k}(๐ซ,๐ซ^{},\mathrm{g}_1)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}\varphi _{n_1l_1}^{}(r)\varphi _{n_3l_3}(r^{})\mathrm{exp}[\beta r^2]P_{l_3}(\mathrm{cos}\omega _{rr^{}})`$ (44)
$`\times {\displaystyle _0^{\mathrm{}}}\varphi _{n_2l_2}^{}(r_2)\varphi _{n_4l_4}(r_2)\mathrm{exp}[\beta r_2^2]i_k(2\beta rr_2)r_2^2\mathrm{d}r_2,`$
and the matrix element $`A_{n_1l_1n_2l_2}^{n_3l_3n_4l_4,k}(๐ซ,๐ซ^{},\mathrm{g}_2)`$ can be found from (44) replacing $`๐ซ๐ซ^{}`$ and $`n_1l_1n_3l_3`$ while the matrix element corresponding to the factor $`\mathrm{g}_3`$ is
$`A_{n_1l_1n_2l_2}^{n_3l_3n_4l_4,k}(๐ซ,๐ซ^{},\mathrm{g}_3)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}\varphi _{n_1l_1}^{}(r)\varphi _{n_3l_3}(r^{})\mathrm{exp}[\beta (r^2+r^2)]\mathrm{\Omega }_{l_1l_3}^k(\omega _{rr^{}})\times `$ (46)
$`{\displaystyle _0^{\mathrm{}}}\varphi _{n_2l_2}^{}(r_2)\varphi _{n_4l_4}(r_2)\mathrm{exp}[2\beta r_2^2]i_k(2\beta |๐ซ+๐ซ^{}|r_2)r_2^2dr_2,`$
In Eqs. (44) and (46) the modified spherical Bessel function, $`i_k(z)`$, comes from the expansion of the exponential function $`\mathrm{exp}[2\beta ๐ฑ_1๐ฑ_2]`$ of the factors $`\mathrm{g}_{\mathrm{}}`$ in spherical harmonics, that is
$`\mathrm{exp}[2\beta ๐ฑ_1๐ฑ_2]=4\pi {\displaystyle \underset{km_k}{}}i_k(2\beta x_1x_2)Y_{km_k}^{}(\mathrm{\Omega }_1)Y_{km_k}(\mathrm{\Omega }_2),`$
while the factor $`\mathrm{\Omega }_{l_1l_3}^k(\omega _{rr^{}})`$ which depends on the directions of $`๐ซ`$ and $`๐ซ^{}`$ is,
$$\mathrm{\Omega }_{l_1l_3}^k(\omega _{rr^{}})=\underset{m_1,m_3}{}\frac{l_1m_1l_3m_3k0}{l_10l_30k0}\left(\frac{(l_1m_1)!}{(l_1+m_1)!}\frac{(l_3m_3)!}{(l_3+m_3)!}\right)^{\frac{1}{2}}P_{l_1}^{m_1}(\mathrm{cos}\omega _r)P_{l_3}^{m_3}(\mathrm{cos}\omega _r^{}),$$
(47)
where $`\omega _r`$ and $`\omega _r^{}`$ are the angles between the vectors $`๐ซ,๐ซ+๐ซ^{}`$ and $`๐ซ^{},๐ซ+๐ซ^{}`$, respectively. The final expression of $`\mathrm{\Omega }_{l_1l_3}^k(\omega _{rr^{}})`$ depends on $`\omega _{rr^{}}`$, the angle between the vectors $`๐ซ`$ and $`๐ซ^{}`$.
The expression of the term $`O_{22}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})`$ depends on the SP wave functions and so it is suitable to be used for analytical calculations with the HO orbitals and in principle for numerical calculations with more realistic SP orbitals. Expressions (40) and (42) were derived for the closed shell nuclei with $`N=Z`$, where $`\eta _{nl}`$ is 0 or 1. For the open shell nuclei (with $`N=Z`$) we use the same expressions, where now $`0\eta _{nl}1`$. In this way the mass dependence of the correlation parameter $`\beta `$ and the OBDM or MD can be studied.
Finally, using the known values of the Clebsch-Gordan coefficients, Eq. (42), for the case of $`s`$-$`p`$ and $`s`$-$`d`$ shell nuclei, takes the form
$`O_{22}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})=`$ (48)
$`4[{\displaystyle \frac{}{}}3A_{0000}^{0000,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})\eta _{1s}^2+[33A_{0101}^{0101,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})6A_{0101}^{0101,2}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})]\eta _{1p}^2+3A_{1010}^{1010,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})\eta _{2s}^2`$ (49)
(50)
$`+\left[95A_{0202}^{0202,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}}){\displaystyle \frac{50}{7}}A_{0202}^{0202,2}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}}){\displaystyle \frac{90}{7}}A_{0202}^{0202,4}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})\right]\eta _{1d}^2`$ (51)
(52)
$`+\left[12A_{0001}^{0001,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})+12A_{0100}^{0100,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})3A_{0001}^{0100,1}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})3A_{0100}^{0001,1}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})\right]\eta _{1s}\eta _{1p}`$ (53)
(54)
$`+\left[20A_{0002}^{0002,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})+20A_{0200}^{0200,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})5A_{0002}^{0200,2}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})5A_{0200}^{0002,2}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})\right]\eta _{1s}\eta _{1d}`$ (55)
(56)
$`+\left[4A_{0010}^{0010,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})+4A_{1000}^{1000,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})A_{0010}^{1000,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})A_{1000}^{0010,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})\right]\eta _{1s}\eta _{2s}`$ (57)
(58)
$`+[60A_{0102}^{0102,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})+60A_{0201}^{0201,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})6A_{0102}^{0201,1}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})6A_{0201}^{0102,1}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})`$ (59)
(60)
$`9A_{0102}^{0201,3}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})9A_{0201}^{0102,3}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})]\eta _{1p}\eta _{1d}`$ (61)
(62)
$`+\left[12A_{0110}^{0110,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})+12A_{1001}^{1001,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})3A_{0110}^{1001,1}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})3A_{1001}^{0110,1}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})\right]\eta _{1p}\eta _{2s}`$ (63)
(64)
$`+[20A_{0210}^{0210,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})+20A_{1002}^{1002,0}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})5A_{0210}^{1002,2}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})5A_{1002}^{0210,2}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})]\eta _{_{1d}}\eta _{2s}{\displaystyle \frac{}{}}].`$ (65)
It should be noted that Eqs. (42) and (65) are also valid for the cluster expansion of the density distribution and the form factor as it has been found in ref. and also in the cluster expansion of the MD. The only difference is the expressions of the matrix elements $`A`$.
### B Momentum distribution
The MD for the above mentioned nuclei can be found either by following the same cluster expansion or by taking the Fourier transform of $`\rho (๐ซ,๐ซ^{})`$ given by (36). In both cases the correlated momentum distribution takes the form
$$n_{cor}(๐ค)N\left[\stackrel{~}{๐}_๐ค_12\stackrel{~}{O}_{22}(๐ค,\mathrm{g}_1)+\stackrel{~}{O}_{22}(๐ค,\mathrm{g}_3)\right],$$
(66)
where
$$\stackrel{~}{๐}_๐ค_1=n_{SD}(๐ค)=\underset{i=1}{\overset{A}{}}\stackrel{~}{\varphi }_i^{}(๐ค)\stackrel{~}{\varphi }_i(๐ค).$$
(67)
The term $`\stackrel{~}{O}_{22}(๐ค,\mathrm{g}_{\mathrm{}})`$, as in the case of OBDM, is given again by the right-hand side of Eqs. (42) and (65) replacing the matrix elements $`A_{n_1l_1n_2l_2}^{n_3l_3n_4l_4,k}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})`$, defined by Eqs. (44) and (46), by the Fourier transform of them, that is by the matrix elements
$$\stackrel{~}{A}_{n_1l_1n_2l_2}^{n_3l_3n_4l_4,k}(๐ค,\mathrm{g}_{\mathrm{}})=\frac{1}{(2\pi )^3}A_{n_1l_1n_2l_2}^{n_3l_3n_4l_4,k}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})\mathrm{exp}[i๐ค(๐ซ๐ซ^{})]d๐ซd๐ซ^{},\mathrm{}=1,3.$$
(68)
As in the case of the OBDM, expression (66) is suitable for the study of the MD for the $`s`$-$`p`$ and $`s`$-$`d`$ shell nuclei and also for the study of the mass dependence of the kinetic energy of these nuclei. The mean value of the kinetic energy has the form
$$๐=N[๐_12T_{22}(\mathrm{g}_1)+T_{22}(\mathrm{g}_3)],$$
(69)
where
$$๐_1=\frac{\mathrm{}^2}{2m}k^2n_{SD}(๐ค)d๐ค,T_{22}(\mathrm{g}_{\mathrm{}})=\frac{\mathrm{}^2}{2m}k^2\stackrel{~}{O}_{22}(๐ค,\mathrm{g}_{\mathrm{}})d๐ค,\mathrm{}=1,3.$$
(70)
## III ANALYTICAL EXPRESSIONS
In the case of the HO wave functions, with radial part in coordinate and momentum space,
$`\varphi _{nl}(r)`$ $`=`$ $`N_{nl}b^{3/2}r_b^lL_n^{l+\frac{1}{2}}\left(r_b^2\right)\mathrm{e}^{r_b^2/2},r_b=r/b,`$ (71)
$`\stackrel{~}{\varphi }_{nl}(k)`$ $`=`$ $`i^l(1)^{n+l}N_{nl}b^{3/2}k_b^lL_n^{l+\frac{1}{2}}\left(k_b^2\right)\mathrm{e}^{k_b^2/2},k_b=kb,`$ (72)
where
$`N_{nl}=\left({\displaystyle \frac{2n!}{\mathrm{\Gamma }(n+l+\frac{3}{2})}}\right)^{1/2},`$
analytical expressions of the one-body terms, $`๐_{\mathrm{๐ซ๐ซ}^{}}_1`$ and $`\stackrel{~}{๐}_๐ค_1`$ as well as of the matrix elements $`A_{n_1l_1n_2l_2}^{n_3l_3n_4l_4,k}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})`$ and $`\stackrel{~}{A}_{n_1l_1n_2l_2}^{n_3l_3n_4l_4,k}(๐ค,\mathrm{g}_{\mathrm{}})`$, which have been defined in Sec. II, can be found. From these expressions, the analytical expressions of the terms $`O_{22}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})`$ and $`\stackrel{~}{O}_{22}(๐ค,\mathrm{g}_{\mathrm{}})`$, defined by Eq. (65), can also be found.
The expressions of the one-body terms, $`๐_{\mathrm{๐ซ๐ซ}^{}}_1`$ and $`\stackrel{~}{๐}_๐ค_1`$, have the forms
$`๐_{\mathrm{๐ซ๐ซ}^{}}_1=\rho _{SD}(๐ซ,๐ซ^{})`$ $`=`$ $`{\displaystyle \frac{2}{\pi ^{3/2}b^3}}[{\displaystyle \frac{}{}}2\eta _{1s}+3\eta _{2s}2\eta _{2s}(r_b^2+r_{b}^{}{}_{}{}^{2})+4\eta _{1p}r_br_b^{}\mathrm{cos}\omega _{rr^{}}`$ (75)
$`+{\displaystyle \frac{4}{3}}[\eta _{2s}+\eta _{1d}(3\mathrm{cos}^2\omega _{rr^{}}1)r_b^2r_{b}^{}{}_{}{}^{2}]{\displaystyle \frac{}{}}]\mathrm{exp}[(r_b^2+r_{b}^{}{}_{}{}^{2})/2]`$
$$\stackrel{~}{๐}_๐ค_1=n_{SD}(๐ค)=\frac{2b^3}{\pi ^{3/2}}\mathrm{exp}[k_b^2]\underset{k=0}{\overset{2}{}}C_{2k}k_b^{2k},$$
(76)
where the coefficients $`C_{2k}`$ are
$$C_0=2\eta _{1s}+3\eta _{2s},C_2=4(\eta _{1p}\eta _{2s}),C_4=\frac{4}{3}(2\eta _{1d}+\eta _{2s}).$$
(77)
The analytical expressions of the matrix element $`A_{n_1l_1n_2l_2}^{n_3l_3n_4l_4,k}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})`$ $`(\mathrm{}=1,3)`$ have the form
$`A_{n_1l_1n_2l_2}^{n_3l_3n_4l_4,k}(๐ซ,๐ซ^{},\mathrm{g}_1)`$ $`=`$ $`B_0b^3y^kr_b^{k+l_1}r_{b}^{}{}_{}{}^{l_3}L_{n_1}^{l_1+\frac{1}{2}}(r_b^2)L_{n_3}^{l_3+\frac{1}{2}}(r_{b}^{}{}_{}{}^{2})\mathrm{exp}\left[{\displaystyle \frac{1+3y}{2(1+y)}}r_b^2{\displaystyle \frac{1}{2}}r_{b}^{}{}_{}{}^{2}\right]`$ (79)
$`\times P_{l_3}(\mathrm{cos}\omega _{rr^{}}){\displaystyle \underset{w_2=0}{\overset{n_2}{}}}{\displaystyle \underset{w_4=0}{\overset{n_4}{}}}B_{w_2w_4}(y)L_{\frac{1}{2}(l_2+l_4k)+w_2+w_4}^{k+\frac{1}{2}}\left({\displaystyle \frac{y^2}{1+y}}r_b^2\right),`$
and
$`A_{n_1l_1n_2l_2}^{n_3l_3n_4l_4,k}(๐ซ,๐ซ^{},\mathrm{g}_3)`$ $`=`$ $`B_0b^3y^k๐ซ_b+๐ซ_b^{}^kr_b^{l_1}r_{b}^{}{}_{}{}^{l_3}L_{n_1}^{l_1+\frac{1}{2}}(r_b^2)L_{n_3}^{l_3+\frac{1}{2}}(r_{b}^{}{}_{}{}^{2})\mathrm{\Omega }_{l_1l_3}^k(\omega _{rr^{}})`$ (82)
$`\times \mathrm{exp}\left[{\displaystyle \frac{1+2y}{2}}(r_b^2+r_{b}^{}{}_{}{}^{2})\right]\mathrm{exp}\left[{\displaystyle \frac{y^2}{1+2y}}(๐ซ_b+๐ซ_b^{})^2\right]`$
$`\times {\displaystyle \underset{w_2=0}{\overset{n_2}{}}}{\displaystyle \underset{w_4=0}{\overset{n_4}{}}}B_{w_2w_4}(2y)L_{\frac{1}{2}(l_2+l_4k)+w_2+w_4}^{k+\frac{1}{2}}({\displaystyle \frac{y^2}{1+2y}}(๐ซ_b+๐ซ_b^{})^2)`$
where $`y=\beta b^2`$ and
$$B_0=\frac{1}{16\sqrt{\pi }}\left(\underset{i=1}{\overset{4}{}}N_{n_il_i}\right),$$
(83)
$$B_{w_2w_4}(z)=[\frac{1}{2}(l_2+l_4k)+w_2+w_4]!\underset{i=2,4}{}\frac{(1)^{w_i}}{w_i!}(\begin{array}{c}n_i+l_i+\frac{1}{2}\\ n_iw_i\end{array})(1+z)^{\frac{1}{2}l_iw_i\frac{1}{2}(3+k)},$$
(84)
while the one corresponding to the factor $`\mathrm{g}_2`$ can be found from (79) replacing $`r_br_b^{}`$ and $`n_1l_1n_3l_3`$.
The substitution of $`A_{n_1l_1n_2l_2}^{n_3l_3n_4l_4,k}(๐ซ,๐ซ^{},\mathrm{g}_{\mathrm{}})`$ to the expression of $`O_{22}(๐ซ,๐ซ^{},g_{\mathrm{}})`$ which is given by Eq. (65) leads to the analytical expression of the two-body term of the OBDM, which is of the form
$`O_{22}(๐ซ_b,๐ซ_b^{})`$ $`=`$ $`f_1(r_b,r_b^{},\mathrm{cos}\omega _{rr^{}})\mathrm{exp}\left[{\displaystyle \frac{1+3y}{2(1+y)}}r_b^2{\displaystyle \frac{1}{2}}r_{b}^{}{}_{}{}^{2}\right]`$ (87)
$`+f_1(r_b^{},r_b,\mathrm{cos}\omega _{rr^{}})\mathrm{exp}\left[{\displaystyle \frac{1+3y}{2(1+y)}}r_{b}^{}{}_{}{}^{2}{\displaystyle \frac{1}{2}}r_b^2\right]`$
$`+f_3(r_b,r_b^{},\mathrm{cos}\omega _{rr^{}})\mathrm{exp}\left[{\displaystyle \frac{1+2y}{2}}(r_b^2+r_{b}^{}{}_{}{}^{2})\right]\mathrm{exp}\left[{\displaystyle \frac{y^2}{1+2y}}(๐ซ_b+๐ซ_b^{})^2\right],`$
where $`f_{\mathrm{}}(r_b,r_b^{},\mathrm{cos}\omega _{rr^{}}),(\mathrm{}=1,3)`$ are polynomials of $`r_b,r_b^{}`$ and $`\mathrm{cos}\omega _{rr^{}}`$ which depend also on $`y=\beta b^2`$ and the occupation probabilities of the various states.
The corresponding analytical expressions of the matrix elements $`\stackrel{~}{A}_{n_1l_1n_2l_2}^{n_3l_3n_4l_4,k}(๐ค,\mathrm{g}_{\mathrm{}})`$, $`(\mathrm{}=1,3)`$ which contribute to the two-body term of the MD were found substituting $`\varphi _{nl}(r)`$ with that of the HO wave function into Eq. (68). The expression of $`\stackrel{~}{A}_{n_1l_1n_2l_2}^{n_3l_3n_4l_4,k}(๐ค,\mathrm{g}_1)`$, which can be found easily, has the form
$`\stackrel{~}{A}_{n_1l_1n_2l_2}^{n_3l_3n_4l_4,k}(๐ค,\mathrm{g}_1)`$ $`=`$ $`B_0b^3(1)^{n_3}k_b^{2l_3}L_{n_3}^{l_3+\frac{1}{2}}\left(k_b^2\right)\mathrm{exp}\left[{\displaystyle \frac{1+2y}{1+3y}}k_b^2\right]`$ (93)
$`\times {\displaystyle \underset{w_1=0}{\overset{n_1}{}}}{\displaystyle \underset{w_2=0}{\overset{n_2}{}}}{\displaystyle \underset{w_4=0}{\overset{n_4}{}}}{\displaystyle \underset{t=0}{\overset{\frac{1}{2}(l_2+l_4k)+w_2+w_4}{}}}\stackrel{~}{B}_{w_2w_4t}(y)[{\displaystyle \frac{1}{2}}(l_1l_3+k)+w_1+t)]!`$
$`\times {\displaystyle \frac{(1)^{w_1}}{w_1!}}\left(\begin{array}{c}n_1+l_1+\frac{1}{2}\\ n_1w_1\end{array}\right)2^{\frac{1}{2}(l_1l_3)+w_1}(1+y)^{\frac{1}{2}(l_1+l_3l_2l_4)+w_1w_2w_4}`$
$`\times (1+3y)^{\frac{1}{2}(3+l_1+l_3+k)w_1t}L_{\frac{1}{2}(l_1l_3+k)+w_1+t}^{l_3+\frac{1}{2}}\left({\displaystyle \frac{1+y}{2(1+3y)}}k_b^2\right),`$
where
$`\stackrel{~}{B}_{w_2w_4t}(y)`$ $`=`$ $`(\sqrt{2}y)^{k+2t}\left[{\displaystyle \frac{1}{2}}(1+l_2+l_4+k)+w_2+w_4\right]!\left[{\displaystyle \frac{kl_2l_4}{2}}w_2w_4\right]_t`$ (97)
$`\times {\displaystyle \frac{1}{(k+t+\frac{1}{2})!}}{\displaystyle \underset{i=2,4}{}}{\displaystyle \frac{(1)^{w_i+t}}{w_i!t!}}(\begin{array}{c}n_i+l_i+\frac{1}{2}\\ n_iw_i\end{array}).`$
The expression of $`\stackrel{~}{A}_{n_1l_1n_2l_2}^{n_3l_3n_4l_4,k}(๐ค,\mathrm{g}_3)`$ is more complicated. It has the the general form
$`\stackrel{~}{A}_{n_1l_1n_2l_2}^{n_3l_3n_4l_4,k}(๐ค,\mathrm{g}_3)`$ $`=`$ $`{\displaystyle \frac{2}{\sqrt{\pi }}}B_0b^3\mathrm{exp}\left[{\displaystyle \frac{k_b^2}{1+2y}}\right]{\displaystyle \underset{w_2=0}{\overset{n_2}{}}}{\displaystyle \underset{w_4=0}{\overset{n_4}{}}}{\displaystyle \underset{t=0}{\overset{\frac{1}{2}(l_2+l_4k)+w_2+w_4}{}}}\stackrel{~}{B}_{w_2w_4t}(y)`$ (99)
$`\times (1+2y)^{\frac{1}{2}(3+l_2+l_4+k)w_2w_4t}I_{n_1l_1}^{n_3l_3,k}(k_b).`$
The general expression of the quantity $`I_{n_1l_1}^{n_3l_3,k}(k_b)`$ is quite complicated. For that reason we calculated it for various cases which are needed for the $`s`$-$`p`$ and $`s`$-$`d`$ shell nuclei. The various cases and the corresponding expressions of $`I_{n_1l_1}^{n_3l_3,k}(k_b)`$ are given bellow.
#### 1 Case 1: $`n_1l_1=n_3l_3`$ and $`k=0`$
$`I_{n_1l_1}^{n_1l_1,0}(k_b)`$ $`=`$ $`{\displaystyle \underset{\rho =0}{\overset{[l_1/2]}{}}}{\displaystyle \underset{\sigma =0}{\overset{n_1}{}}}{\displaystyle \underset{\tau =0}{\overset{\rho +\sigma }{}}}{\displaystyle \underset{\nu =0}{\overset{2(\rho +\sigma \tau )}{}}}{\displaystyle \underset{\alpha =0}{\overset{l_12\rho }{}}}{\displaystyle \underset{\mu =0}{\overset{n_1\sigma }{}}}{\displaystyle \underset{w_1=0}{\overset{n_1\sigma \mu }{}}}{\displaystyle \underset{w_3=0}{\overset{\mu }{}}}{\displaystyle \frac{(1)^{\rho +\alpha +\tau +w_1+w_3}}{n_1!\rho !\sigma !w_1!w_3!}}L_{\alpha +w_1+\tau +\nu }^{\frac{1}{2}}\left({\displaystyle \frac{k_b^2}{1+2y}}\right)`$ (115)
$`\times (\begin{array}{c}l_12\rho \\ \alpha \end{array})(\begin{array}{c}\rho +\sigma \\ \tau \end{array})(\begin{array}{c}2(\rho +\sigma \tau )\\ \nu \end{array})(\begin{array}{c}n_1+\sigma \mu \\ n_1\sigma \mu w_1\end{array})(\begin{array}{c}\mu +l_1\frac{1}{2}\\ \mu w_3\end{array})`$
$`\times {\displaystyle \frac{(\alpha +w_1+\tau +\nu )!}{2\tau +1}}{\displaystyle \frac{(2l_12\rho )!(n_1+l_1+\frac{1}{2})!(l_1+w_3+t+2\sigma \alpha \tau \nu +\frac{1}{2})!}{(l_1\rho )!(l_12\rho )!(\sigma +l_1+\frac{1}{2})!}}`$
$`\times 2^{l_1+w_1+w_3+t+2\tau }(1+2y)^{l_1w_1+w_3+t+2(\sigma \alpha \tau \nu )}(1+4y)^{\frac{3}{2}l_1w_3t2\sigma +\alpha +\tau +\nu }.`$
#### 2 Case 2: $`n_1l_1=n_3l_3=01`$ and $`k=2`$
$`I_{01}^{01,2}(k_b)`$ $`=`$ $`2^{1+t}\left({\displaystyle \frac{5}{2}}+t\right)!(1+2y)^{2+t}(1+4y)^{\frac{7}{2}t}.`$ (116)
#### 3 Case 3: $`n_1l_1=n_3l_3=02`$ and $`k=2,4`$
$`I_{02}^{02,k}(k_b)`$ $`=`$ $`2^{\frac{k}{2}+t}(1+2y)^{1+t}(1+4y)^{\frac{7}{2}t}`$ (118)
$`\times \left[\left({\displaystyle \frac{5+k}{2}}+t\right)!(1+2y)^{1+\frac{k}{2}}(1+4y)^{\frac{k}{2}}\delta _{2k}{\displaystyle \frac{7}{3}}\left({\displaystyle \frac{5}{2}}+t\right)!L_1^{\frac{1}{2}}\left({\displaystyle \frac{k_b^2}{1+2y}}\right)\right].`$
#### 4 Case 4: $`n_1l_1n_3l_3`$ and $`l_1=0`$ or/and $`l_3=0`$
$`I_{n_1l_1}^{n_3l_3,k}(k_b)`$ $`=`$ $`{\displaystyle \underset{w_1=0}{\overset{n_1}{}}}{\displaystyle \underset{w_3=0}{\overset{n_3}{}}}{\displaystyle \underset{\rho _1=0}{\overset{[l_1/2]}{}}}{\displaystyle \underset{\rho _3=0}{\overset{[l_3/2]}{}}}{\displaystyle \underset{\tau _1=0}{\overset{l_12\rho _1}{}}}{\displaystyle \underset{\tau _3=0}{\overset{l_32\rho _3}{}}}{\displaystyle \underset{\sigma _1=0}{\overset{\rho _1+w_1}{}}}{\displaystyle \underset{\sigma _3=0}{\overset{\rho _3+w_3}{}}}{\displaystyle \underset{\nu =0}{\overset{\rho _1+\rho _3+w_1+w_3\sigma _1\sigma _3}{}}}{\displaystyle \frac{1+(1)^{\tau _1+\tau _3+\sigma _1+\sigma _3}}{2(\tau _1+\tau _3+\sigma _1+\sigma _3+1)}}`$ (131)
$`\times [{\displaystyle \underset{i=1,3}{}}(1)^{\rho _i+w_i+\tau _3+\sigma _3}{\displaystyle \frac{2^{\frac{1}{2}kl_i+\sigma _i+t}(2l_i2\rho _i)!}{w_i!(l_i\rho _i)!(l_i2\rho _i)!}}(\begin{array}{c}n_i+l_i+\frac{1}{2}\\ n_iw_i\end{array})(\begin{array}{c}l_i2\rho _i\\ \tau _i\end{array})(\begin{array}{c}\rho _i+w_i\\ \sigma _i\end{array})`$
$`\times (1+2y)^{\frac{1}{2}(k+l_i)+w_i\sigma _i\tau _i2\nu +t}(1+4y)^{\frac{1}{2}(3+k+l_i\sigma _i\tau _i)w_it+\nu }{\displaystyle \frac{}{}}]`$
$`\times \left[{\displaystyle \frac{1+k+l_1+l_3\sigma _1\sigma _3\tau _1\tau _3}{2}}+t+w_1+w_3\nu \right]!\left[{\displaystyle \frac{\sigma _1+\sigma _3+\tau _1+\tau _3}{2}}+\nu \right]!`$
$`\times (\begin{array}{c}\rho _1+\rho _3+w_1+w_3\sigma _1\sigma _3\\ \nu \end{array})L_{\frac{1}{2}(\sigma _1+\sigma _3+\tau _1+\tau _3+2\nu )}^{\frac{1}{2}}\left({\displaystyle \frac{k_b^2}{1+2y}}\right).`$
#### 5 Case 5: $`n_1l_1=01,n_3l_3=02`$ or $`n_1l_1=02,n_3l_3=01`$, and $`k=1,3`$
$`I_{01}^{02,k}(k_b)`$ $`=`$ $`I_{02}^{01,k}(k_b)=2^{\frac{k}{2}+t}\left(1+{\displaystyle \frac{k}{2}}+t\right)!(1+2y)^{\frac{1}{2}(3+k)+t}(1+4y)^{3\frac{k}{2}t}`$ (133)
$`\times \left[\left(2+{\displaystyle \frac{k}{2}}+t\right)+\left({\displaystyle \frac{4}{k^2+k6}}{\displaystyle \frac{2}{3}}\right){\displaystyle \frac{1+4y}{(1+2y)^2}}L_1^{\frac{1}{2}}\left({\displaystyle \frac{k_b^2}{1+2y}}\right)\right].`$
The substitution of $`\stackrel{~}{A}_{n_1l_1n_2l_2}^{n_3l_3n_4l_4,k}(๐ค,\mathrm{g}_{\mathrm{}})`$ to the expression of $`\stackrel{~}{O}_{22}(๐ค,\mathrm{g}_{\mathrm{}})`$ which is given by Eq. (65) leads to the analytical expression of the two-body term of the MD, which is of the form
$`\stackrel{~}{O}_{22}(k)`$ $`=`$ $`\stackrel{~}{f}_1(k_b^2)\mathrm{exp}\left[{\displaystyle \frac{1+2y}{1+3y}}k_b^2\right]+\stackrel{~}{f}_3(k_b^2)\mathrm{exp}\left[{\displaystyle \frac{1}{1+2y}}k_b^2\right],`$ (134)
where $`\stackrel{~}{f}_{\mathrm{}}(k_b^2),(\mathrm{}=1,3)`$ are polynomials of $`k_b^2`$ which depend also on $`y=\beta b^2`$ and the occupation probabilities of the various states. Similar expressions have been found for the mean value of the kinetic energy.
It should be noted that, although the above expressions of the matrix elements $`A`$ and $`\stackrel{~}{A}`$ seem to be quite complicated, they can easily be used for analytical calculations with programs such as Macsyma or Mathematica. As the above expressions have been found for the ($`N=Z`$) $`s`$-$`p`$ and $`s`$-$`d`$ shell nuclei they can be used for the systematic study of the OBDM and MD in this region of nuclei.
## IV RESULTS AND DISCUSSION
The calculations of the MD for the various $`s`$-$`p`$ and $`s`$-$`d`$ shell nuclei, with $`N=Z`$, have been carried out on the basis of Eq. (66) and the analytical expressions of the one- and two-body terms which were given in Sec. III. Two cases have been examined, named case 1 and case 2 corresponding to the analytical calculations with HO orbitals without and with SRC, respectively.
The parameters $`b`$ and $`\beta `$ of the model in case 1 and for <sup>4</sup>He, <sup>16</sup>O, <sup>36</sup>Ar and <sup>40</sup>Ca in case 2 were the ones which have been determined in our previous work by fit of the theoretical $`F_{ch}(q)`$, derived with the same cluster expansion, to the experimental one. These values of the parameters are given in Table I. The values of the correlation parameter $`\beta `$ of the open shell nuclei which have been reported in Ref. were quite large. That is the correlations for these nuclei were quite small. The MD of the open shell nuclei, which we found with these values of the parameters, had a high momentum tail at values of $`k`$ larger than expected. As that seems to us quite unreasonable we tried to redetermine more carefully the parameters of the model by fit of the theoretical $`F_{ch}(q)`$ to the experimental one in order to obtain a better fit.
The new values of $`b`$ and $`\beta `$ for case 2 and for <sup>12</sup>C, <sup>24</sup>Mg, <sup>28</sup>Si and <sup>32</sup>S are shown in Table I. The theoretical $`F_{ch}(q)`$ for these nuclei, which are shown in Fig. 1, are closer to the experimental data than they were in Ref. . From the values of $`\chi ^2`$, which have been found in cases 1 and 2 and also from Fig. 1 it can been seen that the inclusion of SRCโs improves the fit of the form factor of the above mentioned nuclei. Also, all the diffraction minima, even the third one which seems to exist in the experimental data of <sup>24</sup>Mg, <sup>28</sup>Si and <sup>32</sup>S are reproduced in the correct place.
Although the values of the parameters $`b`$ and $`\beta `$, for the open shell nuclei, are different from those reported in Ref. , their behaviour, still, indicates that there should be a shell effect in the case of closed shell nuclei. This behaviour has an effect on the MD of nuclei as it is seen from Fig. 2, where the MD, of the various $`s`$-$`p`$ and $`s`$-$`d`$ shell nuclei calculated with the values of $`b`$ and $`\beta `$ of Table I for case 2, have been plotted. It is seen that the inclusion of SRCโs increases considerably the high momentum component of $`n(k)`$, for all nuclei we have considered. Also, while the general structure of the high momentum component of the MD for $`A=4,12,16,24,28,32,36,40`$, is almost the same, in agreement with other studies , there is an $`A`$ dependence of $`n(k)`$ both at small values of $`k`$ and in the region $`2\mathrm{fm}^1<k<5\mathrm{fm}^1`$. The $`A`$ dependence of the high momentum component of $`n(k)`$ is larger in the open shell nuclei than in the closed shell nuclei. It is seen that the high momentum component is almost the same for the closed shell nuclei <sup>4</sup>He, <sup>16</sup>O and <sup>40</sup>Ca as expected from other studies .
In the previous analysis, the nuclei <sup>24</sup>Mg, <sup>28</sup>Si and <sup>32</sup>S were treated as $`1d`$ shell nuclei, that is, the occupation probability of the $`2s`$ state was taken to be zero. The formalism of the present work has the advantage that the occupation probabilities of the various states can be treated as free parameters in the fitting procedure of $`F_{ch}(q)`$. Thus, the analysis can be made with more free parameters. For that reason we considered case $`2^{}`$ in which the occupation probability $`\eta _{2s}`$ of the nuclei <sup>24</sup>Mg, <sup>28</sup>Si and <sup>32</sup>S was taken to be a free parameter together with the parameters $`b`$ and $`\beta `$. We found that the $`\chi ^2`$ values become better, compared to those of case 2 and the $`A`$ dependence of the parameter $`\beta `$ is not so large as it was before. The new values of $`b`$ and $`\beta `$ are shown in Table I and the theoretical $`F_{ch}(q)`$ in Fig 1. The values of the occupation probability $`\eta _{2s}`$ of the above-mentioned three nuclei are 0.19982, 0.17988 and 0.50921, respectively, while the corresponding values of $`\eta _{1d}`$, which can be found from the values of $`\eta _{2s}`$ through the relation
$`\eta _{1d}=[(Z8)2\eta _{2s}]/10,`$
are 0.36004, 0.56402 and 0.69816, respectively. The MD of these three nuclei together with the closed shell nuclei <sup>4</sup>He, <sup>16</sup>O and <sup>40</sup>Ca found in case 2 are shown in Fig. 3. It is seen that the $`A`$ dependence of the high momentum component is now not so large as it was in case 2. As $`F_{ch}(q)`$ calculated in case $`2^{}`$ is closer to the experimental data than in case 2, we might say that this result is in the correct direction, that is the high momentum component of the MD of nuclei is almost the same. We would like to mention that experimental data for $`n(k)`$ are not directly measured but are obtained by means of $`y`$-scaling analysis and only for <sup>4</sup>He and <sup>12</sup>C in $`s`$-$`p`$ and $`s`$-$`d`$ shell region. We expect that the above conclusion could be corroborated if new experimental data are obtained in the future for MD for several nuclei and we carry out a simultaneous fit both to MD and to form factors.
Finally, in table I we give the one and the two-body terms of the mean kinetic energy, $`๐`$, of the various $`s`$-$`p`$ and $`s`$-$`d`$ shell nuclei calculated on the basis of Eq. (69), as well as the rms charge radii, $`r_{ch}^2^{1/2}`$ which are compared with the experimental values. It is seen that the introduction of SRCโs (in case 2) increases the mean kinetic energy relative to case 1 ($`(๐_{case2}๐_{case1})/๐_{case2}`$) about $`50\%`$ in <sup>4</sup>He and $`23\%`$ in <sup>24</sup>Mg. This relative increase follows the fluctuation of the parameter $`\beta `$. Also the values of the kinetic energy in percents, $`100๐_{SRC}/๐_{Total}`$, as well as the ratio $`<๐_{Total}>/๐_{HO}`$ follow the fluctuation of the parameter $`\beta `$. In closed shell nuclei there is an increase of the above values by the increasing of mass number.
## V SUMMARY
In the present work, general expressions for the correlated OBDM and MD have been found using the factor cluster expansion of Clark and co-workers. These expressions can be used for analytical calculations, with HO orbitals and in principle for numerical calculations with more realistic orbitals.
The analytical expressions of the OBDM, MD and mean kinetic energy for the $`s`$-$`p`$ and $`s`$-$`d`$ shell nuclei, which have been found, are functions of the HO parameter $`b`$, the correlation parameter $`\beta `$ and the occupation probabilities of the various states. These expressions are suitable for the systematic study of the above quantities for the $`N=Z`$, $`s`$-$`p`$ and $`s`$-$`d`$ shell nuclei and also for the study of the dependence of these quantities on the various parameters.
It is found that, while the general structure of the MD at high momenta is almost the same for all the nuclei we have considered, in agreement with other studies, there is an $`A`$ dependence on $`n(k)`$ both at small values of $`k`$ and the high momentum component. The $`A`$ dependence of the high momentum component becomes quite small if the occupation probability of the $`2s`$-state for the nuclei <sup>24</sup>Mg, <sup>28</sup>Si and <sup>32</sup>S is treated as a free parameter in the fitting procedure of the charge form factor.
##
The authors would like to thank Professor M.E. Grypeos and Dr. C.P. Panos for useful comments on the manuscript. One of the author (Ch.C.M) would like to thank Dr. P. Porfyriadis for technical assistants.
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# Level Set Modeling of Transient Electromigration Grooving This research was supported by the Israeli Ministry of Science and Technology grant #9672-1-96 - 9672-3-98.
## 1 Introduction
This paper is a continuation of our work on numerical modeling of the formation and propagation of groove-like defects at GBs in thin film polycrystalline interconnects used in microelectronics (ME).
In modern ME industry, the reliability of ME integrated circuits has become no less important than their performance. Some of the most vulnerable elements of ME circuits, susceptible to several types of failures, are the interconnects. These are thin film metallic conductors which connect the active elements.
The defects (due to the small cross-section, high current density, mechanical stresses and presence of GBs acting as fast diffusion pathways) lead to a loss (in relatively short times) of electrical and mechanical integrity, i.e. to line opens or shorts. For example, in the presence of a large GB flux ($`J_{gb}10^4\mu m^2/s`$) a groove can extend several micrometers in a few hours . Thus, GB grooving is one of the main failure mechanisms in advanced integrated circuits.
In the absence of an external potential field and mechanical stresses, the GB atomic flux $`J_{gb}=0`$, and the corresponding groove profile evolves via surface diffusion under well-known conditions of scale and temperature (the so-called Mullins problem ). In this case, mass transport by surface diffusion is driven only by the surface Laplacian of curvature. Essentially, matter flows from low-curvature regions to high-curvature regions.
In , we presented and discussed the numerical approach (e.g. the Level Set (LS) method) used to model GB grooving phenomena. We also tested the LS method on two simple, already solved, grooving problems: Mullins problem, and that of GB grooving by surface diffusion in a periodic array of stationary GBs . In both cases, the results obtained by means of the LS method are in good agreement with the theoretical predictions. In this paper, we consider the second geometry only, as being more realistic (see Fig. 1). Due to axial symmetry at $`x=0,x=L`$ (where $`L`$ is the grain size), we do not attempt to calculate groove branches at $`x<0,x>L`$.
Electric fields/currents in metallic conductors provide an additional driving force for surface/GB mass fluxes . In the presence of an electric field, collisions between the conduction electrons and the metal ions lead to drift of the ions. This process is known as electromigration (EM).
GB grooving with a GB flux in real thin film interconnects is a complex problem. An adequate numerical modeling technique should be capable to manage such issues as GB grooving with an arbitrary EM flux, and various ratios of GB to surface diffusivities; the latter was predicted to critically affect groove kinetics and shape, and thus account for various EM failure regimes (see and the references therein). In cited works, analytical and semi-analytical approaches for analysing steady-state grooving regimes were employed. However, to our knowledge, no effort has been made to directly numerically simulate the transient stage during GB grooving. This is the ultimate goal of this paper.
We do not consider mechanical stresses in GBs which, as a matter of fact, are invariably induced by the field (the approximations under which it is reasonable to neglect the stress are discussed in ). Also, under typical operational conditions of ME interconnects, lattice transport may be neglected compared to surface/GB transport .
Our paper proceeds as follows. In Section 2, we give details of the physical formulation. In Section 3, we discuss some improvements in the numerical algorithm and also the aspects which are due to incorporation of the electric field/GB flux in the model (for other algorithmic details, the interested reader should refer to ). Our numerical results and discussion are presented in Section 4.
## 2 Physical model
### 2.1 Driving forces for the diffusion
In the absence of an electric current, the surface diffusion is driven by a variation in chemical potential, $`\mu _s`$, which causes atoms to migrate from high potential to low potential regions. It was shown that
$$\mu _s=K_s\gamma _s\mathrm{\Omega }$$
(2.1)
where $`K_s`$ is the surface curvature, $`\gamma _s`$ is the surface tension, and $`\mathrm{\Omega }`$ is the atomic volume. Gradients of chemical potential are therefore associated with gradients of curvature.
Let $`\tau `$ be the tangential direction to the surface profile in 2D. Let $`x,y`$ be Cartesian coordinates along horizontal and vertical boundaries of the computational box (Fig. 1). If $`๐ง=(n_x,n_y)`$ is the unit vector normal to the surface, then the following relations hold:
$$\tau =(n_y,n_x),\frac{K_s}{\tau }=K_s\tau =\frac{K_s}{x}n_y\frac{K_s}{y}n_xK_\tau ^s.$$
(2.2)
The corresponding surface flux (volume crossing unit length per unit time) is then given by
$$J_s^{K_s}=\frac{D_s\delta _s}{kT}\frac{\mu _s}{\tau }=BK_\tau ^s$$
(2.3)
where the superscript indicates that the flux is due to the curvature gradient,
$$B=\frac{D_s\delta _s\gamma _s\mathrm{\Omega }}{kT}$$
(2.4)
is known as Mullins constant, and $`D_s,\delta _s,k`$, and $`T`$ denote surface diffusion coefficient, thickness of the surface diffusion layer, Bolzmannโs constant and absolute temperature, respectively. Note that $`J_s^{K_s}`$ is proportional to the first directional derivative of the curvature.
If an electric field is present, the flux $`J_s`$ of matter at the curved surface of the conductor is driven simultaneously by curvature gradients, and by the component $`E`$ of the local electric field along the surface. In what follows, we distinguish between two models which handle the electric field:
Let $`C`$ and $`O`$ denote the conductor (interconnect) material domain and the outer (surrounding) material domain above the surface profile, respectively (see Fig. 1). In this model, the vector of the electric field intensity is parallel to the GBs, $`๐_\mathrm{๐}=(0,E_{0_y})`$, and $`E_{0_y}`$ is a step function,
$$E_{0_y}(x_i,y_j,t)=\{\begin{array}{cc}E_0^{in}=const.\hfill & \text{ if grid point }(x_i,y_j)C\hfill \\ E_0^{out}=const.\hfill & \text{ if grid point }(x_i,y_j)O.\hfill \end{array}$$
(2.5)
We assume that the surrounding material is less conductive than the interconnect material, and therefore $`|E_0^{out}|<|E_0^{in}|`$. In our numerical experiments we chose the ratio $`|E_0^{out}|/|E_0^{in}|=0.1`$. In the finite difference approach, the discontinuous distribution of electric field intensity is smoothed out across the surface profile (see the details in Section 3). The component of the local electric field along the surface, $`E`$, is then approximated by the projection of $`๐_\mathrm{๐}`$ on the surface, $`E=๐_\mathrm{๐}\tau `$. This approximates the true value given by solving Laplaceโs equation for the potential, subject to the boundary conditions of constant fields of magnitudes $`E_0^{in}`$ and $`E_0^{out}`$ in the conductor and surrounding material domains, respectively.
The corresponding electrically induced surface flux of matter is given by
$$J_s^E=\frac{D_s\delta _sZ_s}{kT}E=B_eE$$
(2.6)
where the superscript indicates that the flux is due to the electric field, and
$$B_e=\frac{D_s\delta _sZ_s}{kT}$$
(2.7)
where $`Z_s=z_s^{}e`$ is the effective charge of the ions undergoing electromigration in the surface layer and $`e`$ is the unit electronic charge; the sign of $`z_s`$ is usually positive (i.e., matter flux in the direction of the electron flow).
Assume that (at a given time step of overall marching algorithm) $`U(x,y)`$ is the electric potential within the (rectangular) computational box. $`U^{}(U^+)`$ and $`U^+(U^{})`$ are its values on the upper and lower boundaries of the box, and $`U_n`$ is the normal derivative on the boundary. $`U^{}`$ and $`U^+`$ are assumed to be time-independent and uniform along the boundaries; $`U^+U^{}`$ is the external voltage applied to the interconnect. The distribution $`U(x,y)`$ is governed by a static elliptic PDE
$$\frac{}{x}\left(k\frac{U}{x}\right)+\frac{}{y}\left(k\frac{U}{y}\right)=0,$$
(2.8)
with boundary conditions $`U_n=U_x=0`$ on the vertical boundaries of the box (which in our case coincide with GBs). Equation (2.8) is derived from the well-posed three dimensional potential problem for the two-layer interconnect. The assumptions and complete derivation for the case of small aspect ratio are presented in . We also give some details in the Appendix. In eq. (2.8), $`k=k(x,y)`$ is the specific electrical conductivity (at a given time step) of the material which fills the computational box. To solve (2.8), a finite difference scheme was developed and analysed in . The distribution of the specific conductivity in the physical system under consideration is discontinuous: the conductivity inside the conductor material (domain $`C`$, Fig. 1) differs by a finite value from that of the surrounding material (domain $`O`$). We assume
$$k=\{\begin{array}{cc}k_{in}=const.>0\hfill & \text{ if grid point }(x_i,y_j)C\hfill \\ k_{out}=const.>0\hfill & \text{ if grid point }(x_i,y_j)O,\hfill \end{array}$$
(2.9)
i.e. $`k=k(y)`$ is a step function. In our numerical experiments we chose the ratio $`k_{out}/k_{in}=0.1`$. Since the surface of the conductor evolves in time and space, then to find the time-dependent solution $`U(x,y,t)`$ we need to solve the static equation (2.8) every time step with $`k`$ given by (2.9). In order to be able to compute accurately the electric field intensity (which is the derivative of $`U`$), the discontinuous distribution of the specific conductivity is smoothed out across the surface profile. The finite difference discretization of (2.8) in the computational domain leads to a set of linear algebraic equations with a sparse banded matrix. This set is solved with an effective multigrid iterative procedure . The solution of the previous time step is used as an initial approximation for the current step which allows fast convergence.
After the potential is established everywhere in the computational domain, the corresponding electrically induced surface flux $`J_s^E`$ is given by (2.6), where
$$E=\tau U$$
(2.10)
To summarize the above discussion, the total flux of matter along the surface is
$$J_s=J_s^K+J_s^E.$$
(2.11)
Physically, equation (2.11) says that atoms will diffuse in the direction of the electron flow if the field dominates, but toward the position with the large curvature if the surface energy dominates. This competition between the electric field and the surface energy is essential for the groove dynamics.
The electric field results also in the diffusion of matter along GBs. The diffusion flux along the GB, $`J_{gb}`$, is given by
$$J_{gb}=\frac{D_{gb}\delta _{gb}Z_{gb}}{kT}E$$
(2.12)
where $`D_{gb},\delta _{gb},Z_{gb}=z_{gb}^{}e>0`$ are the GB diffusion coefficient, thickness and effective ionic charge, respectively, and $`E`$ is the component of the electric field along the GB.
### 2.2 Boundary conditions at groove roots
The evolution of the surface is constrained by two conditions imposed locally at groove roots $`a`$ and $`b`$ (Fig. 1):
The boundary condition is dictated by the local equilibrium between the surface tension, $`\gamma _s`$, and the GB tension, $`\gamma _{gb}`$. In the symmetric case of a GB ($`x=0`$) normal to an original ($`y=const.`$) flat surface, the angle of inclination of the right branch of the surface at the groove root with respect to the $`x`$ axis is
$$\theta _0=sin^1(\gamma _{gb}/2\gamma _s)=const.$$
(2.13)
The rapid establishment of the equilibrium angles between the GBs and the surface by atomic migration in the vicinity of the intersections develops some curvature gradients at the adjacent surface, and thus induces surface diffusion fluxes, $`J_s^K`$, along the groove walls. The directions of the fluxes depend on the sign of the respective surface curvature gradients at the groove groots.
The boundary conditions read
$$J_{gb}=2J_s^E,$$
(2.14)
since both branches of the groove act as sinks or sources of matter.
## 3 The numerical procedure
The Level Set method is used to โcaptureโ the evolution of the conductor surface. The method was introduced by Osher and Sethian and was further developed during the last several years. The method enables to capture drastic changes in the shape of the curves (surfaces or interfaces) and even topology changes.
The basic idea of the method consists of embedding the curve $`y(x,t)`$ into a higher dimensional space. As a matter of fact, we consider the evolution of a two-dimensional field $`\varphi (x,y,t)`$ such that its zero level set, $`\varphi (x,y,t)=0`$, coincides with the curve of interest, $`y(x,t)`$, at any time $`t`$. The level set function $`\varphi (x,y,t)`$ can be interpreted as a signed distance from the curve $`y(x,t)`$, which moves in the direction normal to itself.
The evolution of $`\varphi (x,y,t)`$ is described by an Hamilton-Jacobi type equation. A remarkable trait of the method is that the function $`\varphi (x,y,t)`$ remains smooth, while the level surface $`\varphi =0`$ may change topology, break, merge, and form sharp corners as $`\varphi `$ evolves. Thus, it is possible to perform numerical simulations on a discrete grid in the spatial domain, and substitute a finite difference approximations for the spatial and temporal derivatives in time and space. Another nice feature of the method is that the explicit location of the interface needs not to be known in the computational process; all the necessary information is extracted from the level set function.
The evolution equation has the form
$$\varphi _t+F|\varphi |=0\text{given}\varphi (x,t=0).$$
(3.1)
The normal velocity, $`F`$, is considered to be a function of spatial derivatives of $`\varphi (x,y,t)`$. In many applications $`F`$ is a function of the curvature, $`K_s`$, and its spatial derivatives. The curvature $`K_s`$ may be computed via the level set function $`\varphi `$ as follows:
$$K_s=๐ง,๐ง=\frac{\varphi }{|\varphi |}=(\frac{\varphi _x}{\left(\varphi _x^2+\varphi _y^2\right)^{1/2}},\frac{\varphi _y}{\left(\varphi _x^2+\varphi _y^2\right)^{1/2}}).$$
(3.2)
Here $`๐ง`$ is a โnormal vectorโ, and it coincides with the (previously introduced) unit normal to the surface, $`y(x,t)`$, on the zero level set $`\varphi =0`$. Formulas (3.2) can be combined as follows
$$K_s=\frac{\varphi }{|\varphi |}=\frac{\varphi _{xx}\varphi _y^22\varphi _x\varphi _y\varphi _{xy}+\varphi _{yy}\varphi _x^2}{\left(\varphi _x^2+\varphi _y^2\right)^{3/2}},$$
(3.3)
and the sign of $`K_s`$ is chosen such that a sphere has a positive mean curvature equal to its radius. In the case of surface diffusion in 2D,
$$F=\frac{J_s}{\tau }=\frac{J_s^K}{\tau }+\frac{J_s^E}{\tau }$$
(3.4)
where $`J_s`$ is given by (2.11).
The difficulties in the numerical solution of (3.1) in our case are due to the fact that, as could be noted from (2.3), (3.3), (3.4), the first term in $`F`$ contains space derivatives of order 4 of the level set function. Therefore, the evolution equation (3.1) is highly sensitive to errors. Besides, this fourth derivative term leads to schemes with very small time steps.
In , we presented the computational algorithm which solves the problem of GB grooving by surface diffusion in the absence of electromigration. This could be viewed as the limiting case of the problem which is under consideration in this paper, corresponding to the situation where electrically induced surface and GBs fluxes $`J_s^E`$ and $`J_{gb}`$ vanish. The normal velocity function (3.4) in the latter case contains only the first term. The basic features of the algorithm are:
* the use of a uniform static grid in both space directions
* the use of a standard second order-accurate finite difference scheme in space
* the approximation of spatial derivatives (in normal direction) on the boundaries of the computational box by second order one-sided differences
* time marching is done by a second-order Total Variation Diminishing (TVD) RungeโKutta procedure
* the use of second-order Essentially Non-Oscillatory (ENO) scheme to approximate the gradient function in (3.1)
* the use of โreinitializationโ every time step to keep the level set function $`\varphi `$ a signed distance function
The solution of (3.1) (in Mullins case of an infinite bicrystal with a single GB) subject to the conditions of a constant angle of surface inclination and zero surface flux $`J_s^K`$ at the groove root $`a`$, is then a self-similar surface profile, whose linear dimensions are proportional to $`(Bt)^{1/4}`$, $`B`$ given by (2.4) (Fig. 2). If the dimensions of the crystal are finite, grooves develop at each GB; grooving stops when, at sufficiently long times, identical circular arcs develop connecting adjacent GBs (Fig. 3). The parameters chosen for the runs are typical for copper interconnects at temperatures relevant to experiments ($`T=600K`$) : $`\mathrm{\Omega }=1.18\times 10^{29}m^3,D_s=3.3\times 10^{14}m^2/s,\gamma _s=1.7J/m^2,\delta _s=3.5\times 10^{10}m,kT=8.28\times 10^{21}J`$. It is worth noting that our numerical treatment of GB grooving is not constrained by the assumption of small equilibrium angles (โsmall slope approximationโ), in contrast to the analytical approaches of the pioneer works .
In , special attention was given to the treatment of constant-angle and zero-flux conditions at the groove roots within the framework of the Level Set method. Two methods were developed, the first based on interface reconstruction every time step, with subsequent correction of the angles followed with reinitialization, and the second based on the extension of the $`\varphi `$-field beyond the GBs using the expansion in Taylor series up to second order. Both methods were successfully used in calculations, but we observed that sometimes both procedures resulted in a loss of accuracy. In this paper, we propose a new, robust and highly accurate procedure to keep the equilibrium angle (2.13) constant at the intersections of the surface with the GBs.
Consider equations (2.2), (3.2) and the zero level line of $`\varphi `$ (conductorโs surface) passing through the groove root point $`a`$ in Fig. 1 (a similar analysis could be performed for groove root point $`b`$). Since the tangential vector to the zero level line at $`a`$ (as well as to other level lines at $`x=0`$ if $`\varphi `$ is kept a signed distance function), is $`\tau =(\mathrm{cos}\theta _0,\mathrm{sin}\theta _0)`$, then (2.2), (3.2) imply
$$n_x=\mathrm{sin}\theta _0=\frac{\varphi _x}{\left(\varphi _x^2+\varphi _y^2\right)^{1/2}},n_y=\mathrm{cos}\theta _0=\frac{\varphi _y}{\left(\varphi _x^2+\varphi _y^2\right)^{1/2}}.$$
(3.5)
Dividing the first equation in (3.5) by the second one gives
$$\varphi _x=\varphi _y\mathrm{tan}\theta _0.$$
(3.6)
As pointed out above, we approximate the spatial derivatives (in normal direction) of $`\varphi `$ at the boundaries of the computational domain by second order one-sided differences. Therefore, equation (3.6) written on the left boundary $`x=0`$ takes the form
$$\frac{3\varphi _{0,j}+4\varphi _{1,j}\varphi _{2,j}}{2\mathrm{\Delta }x}=\frac{\varphi _{0,j+1}\varphi _{0,j1}}{2\mathrm{\Delta }y}\mathrm{tan}\theta _0,j=0,\mathrm{}m1$$
(3.7)
where $`m`$ is the number of grid points in the vertical direction, and $`\mathrm{\Delta }x`$ and $`\mathrm{\Delta }y`$ are grid spacings in the horizontal and vertical directions, respectively. Rearranging the terms in (3.7) gives the following set of nonhomogeneous linear algebraic equations with a tridiagonal matrix for unknowns $`\varphi _{0,j},j=0,\mathrm{}m1`$:
$$\frac{\mathrm{tan}\theta _0}{2\mathrm{\Delta }y}\varphi _{0,j1}\frac{3}{2\mathrm{\Delta }x}\varphi _{0,j}+\frac{\mathrm{tan}\theta _0}{2\mathrm{\Delta }y}\varphi _{0,j+1}=\frac{4\varphi _{1,j}\varphi _{2,j}}{2\mathrm{\Delta }x}.$$
(3.8)
The solution to this (and to the similar set at $`x=L`$) is easily and accurately found at the beginning of every stage of a RungeโKutta time marching, thus providing the field $`\varphi `$ which incorporates the correct equilibrium angles at the groove roots.
The described procedure allows us to have a straight horizontal line
$$y(x,0)=const.$$
(3.9)
where $`const.>0`$ gives the initial height of the material domain, as initial condition for LS simulations. Note that initial condition (3.9) does not match the boundary condition (2.13). This implies that a singularity exists at $`x=0,x=L`$ at $`t=0`$. This singularity does not present a barrier in solving the system numerically when we select an appropriate numerical scheme. Physically, the equilibrium angle is formed instantaneously compared with the time needed for the evolution of the surface. We are not concerned with the details of this instance. We could use an initial surface which is consistent with the boundary conditions (as done in where Mullins profile served as an initial condition, and we followed the evolution of this profile in time). However, the choice of initial condition (3.9) is more physical.
To close this section, we present the details of the calculation of the normal velocity function (3.4):
1. Calculate curvature induced flux $`J_s^{K_s}`$ from (2.3). It is nonzero even at the first time step, since the equilibrium angles (2.13) are formed instantly
2. Calculate the first term in (3.4) by applying the formula
$$\frac{J_s^K}{\tau }=\left[J_s^K\tau \right]\tau =$$
(3.10)
$$B\left[\frac{K_{xx}\varphi _y^2+2K_{xy}\varphi _x\varphi _yK_{yy}\varphi _x^2}{\varphi _x^2+\varphi _y^2}+K\left(K_y(n_x+n_y)K_x(n_yn_x)\right)\right]+KJ_s^{K_s}$$
3. Solve the electrical problem, find electrically induced surface flux $`J_s^E`$ from (2.6). As pointed out above, the discontinuous distributions of electrical quantities are smoothed across the surface profile \- by a hyperbolic tangent law:
$$r=\frac{r_{out}+r_{in}}{2}+\frac{r_{out}r_{in}}{2}\mathrm{tanh}\beta \varphi (x,y)$$
(3.11)
where $`\beta `$ is a large constant adjusting parameter, and $`r`$ is either the electric intensity $`E(x,y)`$ (first electrical model), or the specific conductivity $`k(x,y)`$ (second electrical model)
4. Given values of the electrical intensity $`E`$ along the GBs, calculate electrically induced GB fluxes, $`J_{gb}`$, from (2.12). Applying boundary condition (2.14), calculate corrected values of $`J_s^E`$ along grid lines $`x=0,x=L`$
5. Calculate the second term in (3.4) by applying the formula (2.2) where $`J_s^E`$ replaces $`K_s`$.
## 4 Numerical results and discussion
Several comments should be made before we present the results of the numerical simulations.
* Due to the large number of material parameters involved we concentrate on the influence of the one which was predicted to greatly affect the grooving process, i.e. the ratio of the GB to surface diffusivity, $`r_d=D_{gb}/D_s`$ . The parameter set we choose for the simulations corresponds to copper, $`Cu`$, at temperatures about 600 K. It should be noted that (i) the experimentally measured values of diffusivities could vary, according to different sources, by up to 3 orders of magnitude, and (ii) $`D_s`$ can be smaller than $`D_{gb}`$, due to, for example, surface contamination, thus giving $`r_d>1`$. Accordingly, we fix the value of $`D_s`$ and vary $`D_{gb}`$ in a wide range, thus varying the GB flux $`J_{gb}`$ (2.12).
* We study the advancing (elongating) grooves characterized by the positive values of $`J_{gb}`$ (matter flows out of the groove cavity and into the GB). In this case the electric field intensity vector is directed upwards (see Fig. 1), the positive potential being prescribed on the lower boundary of the computational box and the negative potential on the upper one. Reversing the direction of the field produces receding grooves or โridgesโ, Fig. 8, characterized by a negative GB flux (matter flows out of the GB into the groove cavity). The advancing grooves are of more practical interest, as explained in the Introduction.
* Most of our results are obtained with electrical model 2 (see section 2.1), based on the solution of Laplaceโs equation for the potential. As expected, this model produced more accurate results than the approximate model 1 based on the piecewise-constant electrical field. However, model 1 <sup>1</sup><sup>1</sup>1Or even more simplified model of the constant field throughout the entire domain $`C+O`$ (Fig. 1) used in most, if not all studies of electromigration (see for example) proved to be rather useful in our simulations, since it produced qualitatively good results in greatly reduced (comparable to the model 2) computational times; for very fine grids ($`80\times 80`$ resolution) the speedup achieved is by as much as a factor of 6.
Fig. 4 (a)-(d) shows the space/time evolution of the initially flat surface of the conductor for different values of $`r_d`$. The parameters are as in Figs. 2, 3 and $`U^+=U^{}=5.0\times 10^3V,k_{in}=10^8(\mathrm{\Omega }m)^1,k_{out}=10^7(\mathrm{\Omega }m)^1,\delta _{gb}=\delta _s=3.5\times 10^{10}m,z_s^{}=z_{gb}^{}=5`$. The surface profiles are dumped every 5000 time steps, the dimensions of the computational box are $`0.5\mu m\times 0.5\mu m`$, and the grid has a $`60\times 60`$ resolution.
In case $`r_d`$ is much less than one (Fig. 4 (a)), i.e. the GB flux is relatively small, and we observe that the evolution of the surface is similar to the one shown in Fig. 3. It slows down in time, providing to be not very dangerous in the sense of failure of the conductor. The evolution proceeds faster as $`r_d`$ increases (Fig. 4 (b)). After the GB grooves merge and form a single profile, this profile starts to advance slowly (note that the surface curvatures at the groove roots are still positive, at least during the time of the observation). Yet larger values of $`r_d`$ (Fig. 4 (c)) result in changes of the morphology of the surface profile in the near-groove-tip regions. The latter means that the sign of the surface curvatures at the groove tips changes from being positive to a negative one in a realively short time after the evolution starts, indicating the surface tendency to form so-called slits. In the case of Fig. 4 (c) this transition takes place at the time step $`n^{},15000<n^{}>20000`$; see also Figures 6, 7. No qualitative change in the surface shapes is observed as $`r_d`$ is further increased - up to the limit where the numerical method is applicable (note the significant losses of accuracy in Fig. 4 (d), which corresponds to $`r_d=22.424`$). Fig. 4 (d) differs from Fig. 4 (c) only in the increased velocity of the surfaceโs advance and in a more rapid transition from positive to negative curvature at the groove roots (the decreased $`n^{}`$). The evolution regime shown in Fig. 4 (c), (d) is known as the $`A`$-regime .
In Fig. 5 (a) - (d) we plot the distance $`d`$ traveled by the groove tip as a function of time. Fig. 5 (a) corresponds to the case of Fig. 4 (d) ($`r_d=22.424,L=0.5\mu m,U=U^+=U^{}=5.0\times 10^3V`$). One sees that the steady-state velocity of the surface advance is attained rather rapidly, within approximately 6 min. Also note that the groove tip traveled (in only 1 hour) a distance which is a little less than half the grain size. In Fig. 5 (b), (c) we investigate the influence of the grain size, $`L`$. When $`r_d`$ is small and the applied voltage is small too (Fig. 5 (b), $`r_d=2.424,U=1.0\times 10^4V`$), the evolution is driven mostly by the surface flux $`J_s^K`$. Then, smaller sizes of the grains result in larger velocities of the groove tip. This is because the curvature of the surface increases as the grain size decreases, resulting in the increase of $`J_s^K`$. The example of such a transitive grooving regime (from classical Mullins regime, Fig. 4 (a) to $`A`$-regime, Fig. 4 (c), (d)) is presented in Fig. 4 (b).
If electrically induced fluxes $`J_s^E`$ and $`J_{gb}`$ are dominant ($`A`$-regime, Fig. 5 (c), $`r_d=2.424,U=1.0\times 10^3V`$) then, in contrast with the previous case, larger grain sizes result in larger groove tip velocities, as predicted in . The dependence of the groove tip velocity on the applied voltage is illustrated by Fig. 5 (d) ($`r_d=2.424,L=0.5\mu m`$). GB grooving proceeds faster in strong electric fields due to the amplification of the electromigration and associated diffusion fluxes $`J_s^E`$ and $`J_{gb}`$.
For completeness, in Fig. 6 (a)-(d) and in Fig. 7 (a)-(d) we present plots of the surface curvature, the $`J_s^K`$ and $`J_s^E`$ diffusion fluxes, and the normal velocity function $`F`$ for the case of Fig. 4 (c). The data in Fig. 6 correspond to time step 1000 (transient stage), while the data in Fig. 7 correspond to time step 35000 (steady-state stage). Grooves develop faster during the transient stage (compare Fig. 6 (d) and Fig. 7 (d), see also Fig. 5), when the curvature of the surface in the near-groove tips regions is positive and both $`J_s^K>0`$ and $`J_s^E>0`$ fluxes tend to elongate the grooves, providing the flow of matter out of the groove tips. As the steady-state approaches, the curvature of the surface in the near-groove tips regions becomes negative (Fig. 7 (a)) and the $`J_s^E`$ flux is still into the GB but the flux $`J_s^K`$ changes the direction and slows the evolution down (Fig. 7 (b), (c)).
Fig. 9 shows the $`A`$-regime of GB grooving obtained with the use of the electrical model 1 (see section 2.1; to be compared with the Fig. 4 (d)). The computations are less accurate if this model is employed, resulting in highly asymmetric surface profiles. However, the dynamics of surface evolution could be predicted from these simulations and we made a heavy use of the electrical model 1 for trial numerical experiments. It is worth to note that the run time to obtain Fig. 9 is 2.2 hours on SGI workstation with 194 MHz IP25 processor, compared to 8.1 hours for Fig. 4 (d).
## 5 Conclusions
The Level Set method was used to model the GB grooving by surface/GB diffusion in an idealized polycrystalline interconnect. The diffusion is driven by surface curvature gradients and external applied electric field. The results demonstrate the high potential of the LS method for the simulation of complex failure phenomena in microelectronic interconnects. The plans for future research are:
* to obtain more physical results with the current version of the code and compare them to experimental ones
* to improve our numerical procedure to make possible the simulation of the propagation of slits ($`B`$-regime, ). The latter are characterized by (almost) straight vertical walls <sup>2</sup><sup>2</sup>2The flux $`J_s^K`$ is zero along these walls. and negative curvatures at the slit tips. The slits are formed at high values of $`r_d`$, and are supposed to propagate in a local steady state, leaving the rest of the surface behind. Physically, the surface of the conductor cannot accomodate a very large GB flux; the groove tips then become diffusively detached from the remaining surface. At the moment our numerical procedure does not allow to fully trace the evolution of the slits and the surface left behind. In our opinion, a locally refined grid is needed to provide high accuracy in the near-slit-tip regions; however, the adaptation of the LS method to such grids may be not straightforward. In our simulations (which make use of an uniform grid), the instability steps in shortly.
* to incorporate mechanical stress in the analysis
* to speed up the computations by making use of an implicit scheme for the solution of the equation (3.1). This will allow larger time steps.
## 6 Appendix
Derivation of the 2-dimensional electrostatic equation
We consider a conducting strip made of a thin metal film, attached to a strip of non zero conductivity substrate. The metal film may be continuous or it may be made of conducting patches with voids in between. We allow the metal film and substrate to have variable thickness. In the present formulation we neglect the interface resistance. The electrodes are attached to the strip and to the substrate. We may want to compute the local field strength which determines the resulting electromigration. This is a more realistic model then the model based on the assumption of a zero conductivity substrate. It allows us also to consider the behaviour of a metal film with varying effective thickness at no extra cost.
### 6.1 The 3-dimensional problem
The 3-dimensional Problem Ohmโs law implies: $`\stackrel{}{j}=\sigma \stackrel{}{E}=\sigma _3\varphi `$, where $`\stackrel{}{j}`$ is the electric current density vector, $`\stackrel{}{E}`$ is the electric field vector, $`\varphi `$ is the electric potential and $`\sigma `$ is the material conductivity. For steady fields Maxwellโs equations with vanishing space charge give:
$$_3\stackrel{}{j}=0,\text{where}_3=(\frac{}{x},\frac{}{y},\frac{}{z}).$$
(6.1)
Hence
$$_3\left(\sigma _3\varphi \right)=0.$$
(6.2)
At all external (lateral) boundaries there is no flux in the direction of the normal, $`\stackrel{}{n}`$, so that $`\stackrel{}{n}\stackrel{}{j}=0`$, and using (6.1) one gets:
$$\stackrel{}{n}_3\varphi =0.$$
(6.3)
The conditions (6.3), together with values of the potential specified at the strip and substrate edges and the continuity and jump conditions at the interface, constitute boundary conditions for equation (6.2) in the two layers. Thus the three dimensional potential can be found, in principle, as the solution of a well-posed three-dimensional boundary value problem. However such a solution can be very expensive to get in the present geometry, in particular as singularities in the solution will appear at sharp geometrical corners at crystal boundaries or voids, requiring high resolution or complicated integration formulae. To avoid this (probably unrealistic) behaviour of the solution and to avoid solving three dimensional problems many times, as required by the time development of the process, we proceed with an approximate approach suggested by (singular) perturbation theory.
### 6.2 The 2-dimensional equation
We assume that $`\varphi `$ and $`\sigma `$ change over a characteristic length scale $`L`$ in the horizontal directions $`x`$ and $`y`$ but over a scale $`H`$ in the vertical. Furthermore we assume that $`ฯต=H/L1`$. Using scaled variables in (6.2),
$$(X,Y,Z)=(x/L,y/L,z/H)$$
(6.4)
we get:
$$ฯต^2_2(\sigma _2\varphi )+\frac{}{Z}\left(\sigma \frac{\varphi }{Z}\right)=0,\text{where}_2=(\frac{}{X},\frac{}{Y}).$$
(6.5)
Singular perturbation analysis considers an expansion
$$\varphi =\varphi _0+ฯต^2\varphi _1+ฯต^4\varphi _2+\mathrm{}$$
(6.6)
where $`\varphi _k`$ are functions of order $`O(1)`$ in $`ฯต`$. Substitution of (6.6) in (6.5) gives relations for the functions $`\varphi _k`$ by grouping terms according to their order in $`ฯต`$ and equating each group to zero. The zeroth order term gives : $`^2\varphi _0/Z^2=0`$, thus $`\varphi _0`$ is a linear function in $`z`$ for every $`x`$ and $`y`$ while taking into account (6.3) kills off the $`z`$ dependence, so that:
$$\varphi _0=\varphi _0(X,Y).$$
(6.7)
Thus at this stage $`\varphi _0`$ is an arbitrary function of the horizontal coordinates $`X`$ and $`Y`$. The first order equation and the boundary conditions in $`Z`$ result ultimately in the approximate two-dimensional equation for $`\varphi _0`$ :
$$_2(h_1\sigma _1+h_2\sigma _2)_2\varphi _0=0$$
(6.8)
where $`h_1,\sigma _1`$ and $`h_2,\sigma _2`$ are, respectively, the heights and conductivities of the two layers under consideration. The equation (6.8) is solved with boundary conditions in the $`(X,Y)`$ plane. We remark that the approximate independence of the potential $`\varphi `$ on the $`Z`$ coordinate justifies also the two dimensional approach for the electromigration equation. This behaviour is a consequence of the small aspect ratio assumption and the normal derivative boundary conditions (6.3), where one must also involve a small slope assumption.
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# VLA Imaging of the Disk Surrounding the Nearby Young Star TW Hya
## 1 Introduction
The TW Hya association, which consists of at least 13 pre-main-sequence systems at a distance of about 50 pc, is thought to be the nearest region of recent star formation (Kastner et al. 1997, Webb et al. 1999a,b). The age of the association members has been estimated by various techniques at 5 to 20 million years, an important time period for the formation of planetary systems (Kastner et al. 1997, Jura et al. 1998, Soderblom et al. 1998, Webb et al. 1999b, Weintraub et al. 2000). Because of its close proximity and interesting age, the association has proved to be fertile ground for observations of phenomena related to the structure and evolution of circumstellar disks. Four of the association members were detected by IRAS indicating the presence of circumstellar dust, either precursors to the assemblage of larger bodies or perhaps the resulting debris. Recent observations with better angular resolution toward members of this association have revealed a disk around TW Hya visible in scattered near-infrared light (Weinberger et al. 1999) and detectable in various molecular tracers (Zuckerman et al. 1995, Kastner et al. 1997), a dust disk around HR 4796A with a cleared interior gap (Koerner et al. 1998, Jayawardhana et al. 1998, Schneider et al. 1999), a dust disk within the hierarchical quadruple system HD 98800 (Low et al. 1999), and direct images of a substellar companion to CoD -33 7795 (Lowrance et al. 1999, Webb et al. 1999b).
Among the TW Hya association members, TW Hya, a K7 star of $`0.5M_{}`$, has the strongest circumstellar dust emission and a disk mass perhaps comparable to the solar nebula prior to the formation of the planets. The TW Hya system is almost three times closer than the young stars in nearby dark clouds like Taurus, Ophiucus and Chameleon, with a distance of $`56\pm 7`$ pc determined by Hipparcos. The technique of millimeter interferometry can now probe emission from the dusty disk material at subarcsecond scales, comparable to the best resolution achieved at optical and infrared wavelengths (Mundy et al. 1996, Wilner et al. 1996, Dutrey et al. 1996, Wilner & Lay 2000). Dust emission at long millimeter wavelengths is nearly entirely optically thin and arises from the full disk column, including the innermost regions of planet formation. In this Letter, we describe observations of thermal emission from TW Hya obtained with the Very Large Array (VLA) of the National Radio Astronomy Observatory <sup>1</sup><sup>1</sup>1The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc.. These observations directly resolve the dusty disk surrounding the star and provide constraints on its structure.
## 2 Observations
We observed TW Hya with the VLA on two occasions in excellent weather conditions, first in a short observation in February 1998 in the D configuration as part of larger survey of dust in T Tauri systems, and then for a full track in October 1999 in the BnA configuration to obtain higher angular resolution. The VLA was divided into two subarrays in each observing session, one consisting of those antennas equipped with 7 mm receivers, and the other consisting of the remaining antennas, which were set to observe at 3.6 cm. Table 1 summarizes the observational parameters. All observations were made with the maximum bandwidth, two 50 MHz channels in each of two circular polarizations, for the best continuum sensitivity. For the observations in the BnA configuration, phase calibration at 7 mm was accomplished by rapid phase referencing within a 120 second cycle to the nearby calibrator 1037-295. The absolute flux scale was determined through observations of the standard source 3C286 and should be accurate to 10%. All data calibration and imaging were performed using the standard routines in the AIPS software package.
## 3 Results
### 3.1 7 mm
Figure 1 shows images made from the 7 mm data at two different resolutions obtained with different visibility weighting schemes. The lower resolution image in the left panel of Figure 1 was made with natural weighting and a 300 k$`\lambda `$ Gaussian taper to produce a $`0\stackrel{}{\mathrm{.}}6`$ beam that emphasizes the spatially extended low brightness emission. The region of detectable 7 mm emission at this resolution is $`100`$ AU in diameter. The total flux integrated within a box of size $`2^{\prime \prime }`$ is $`8\pm 1`$ mJy. It is possible that a fainter halo extends to even larger distances from the peak. The emission presents a roughly circular boundary, consistent with other suggestions that the emission arises in a disk oriented close to pole-on (Kastner et al. 1997), though the ellipticity is uncertain and does not provide a strong constraint on the inclination. The image in the right panel of Figure 1 was made with robust weighting to obtain a $`0\stackrel{}{\mathrm{.}}1`$ beam (5.6 AU), close to the highest resolution available from the longest baselines in these data. Little detectable emission is visible from TW Hya at this scale, though a weak signal remains at the center of the larger structure, with perhaps a comparable peak offset to the northwest. Note that while the flux sensitivity of the higher resolution image is a little better than that of the lower resolution image, the brightness temperature sensitivity is considerably worse because of the smaller synthesized beam.
### 3.2 3.6 cm
TW Hya was detected at 3.6 cm as an unresolved source with flux density $`0.20\pm 0.028`$ mJy. (The results from the two epochs are consistent within the noise.) The emission peak is located at $`11^\mathrm{h}01^\mathrm{m}51\stackrel{\mathrm{s}}{\mathrm{.}}91`$, $`34^{}42^{}16\stackrel{}{\mathrm{.}}96`$ (J2000) with an estimated uncertainty of $`\pm 0\stackrel{}{\mathrm{.}}1`$. Rucinski (1992) previously used the VLA to search for radio emission from TW Hya at 3.6, 6 and 20 cm, without success. He reports an rms noise level of 0.028 mJy at 3.6 cm, very similar to the value obtained here with fewer antennas and more integration time, and so TW Hya should have been detectable in the earlier observation (with a synthesized beam 15 times larger in area). If the observation of Rucinski (1992) is correct, then the radio emission from TW Hya must be time variable.
## 4 Discussion
### 4.1 Disk Spectrum
Figure 2 shows the broadband spectrum of TW Hya from mid-infrared to radio wavelengths. Like many T Tauri stars, the long wavelength emission, far in excess of the stellar photosphere, is well fitted by a family of thin disk models parameterized by radial power laws in temperature and surface density (Adams, Lada & Shu 1988, Beckwith et al. 1990). In these models, the slope of the spectrum in the infrared, where the disk is optically thick, constrains the temperature distribution. At millimeter wavelengths, where the disk is largely optically thin, the emission is proportional to the disk mass weighted by the temperature distribution. Irradiation from the star and low optical depths for the outer disk, together with flaring, tend to drive $`T(r)r^{0.5}`$ (Kenyon & Hartmann 1987, DโAlessio et al. 1998). For illustration, Figure 2 shows spectra from a series of face-on disk models with outer radius 100 AU and the usual (constant) dust opacity law $`\kappa =0.1(\lambda /250\mu m)^\beta `$ cm<sup>2</sup> g<sup>-1</sup>, $`\beta =1`$, and $`\mathrm{\Sigma }(r)(r/1\mathrm{AU})^p`$ g cm<sup>-2</sup>, $`p=0,0.5,1.0`$, and $`1.5`$ (with the mass of gas+dust adjusted from 0.044 to 0.034 $`\mathrm{M}_{}`$ to provide the best least squares fit) for $`T(r)=110(r/1\mathrm{AU})^q`$ K, $`q=0.55`$. Note that the spectrum is not very sensitive to $`p`$, the power law index of the surface density distribution.
Uncertainties in mass opacity coefficient dominate the uncertainties in the estimate of disk mass as the standard value derived from interstellar clouds may not apply (see the review by Beckwith et al. 2000). In addition to issues of grain composition, size and shape, the standard coefficient assumes a gas-to-dust ratio of 100 by mass. For TW Hya, Kastner et al. (1997) observe <sup>13</sup>CO emission and calculate a gas mass of $`3.5\times 10^5`$ $`M_{}`$, three orders of magnitude lower than implied by the millimeter continuum emission. One explanation for this commonly found discrepancy between tracers of total (gas+dust) disk mass is the spectre of severe molecular depletion (Dutrey et al. 1996). Another possibility is a dramatic decrease in the gas-to-dust ratio, which Zuckerman et al. (1995) argue may be especially appropriate for older systems like TW Hya where substantial processing of disk material may have occurred. The disk gas is lost on a timescale of $`10^7`$ yr, perhaps incorporated into giant planets, while dust apparently persists for $`10^8`$ yr or more. For TW Hya, the strong accretion signatures at ultraviolet and optical wavelengths indicate the presence of some gas close to the star. But all estimates of the total disk mass, which is likely dominated by molecular hydrogen, remain somewhat problematic.
At 3.6 cm, emission from the disk models falls far short of the observed flux. By analogy with other young stellar objects, the 3.6 cm emission may arise from hot plasma that originates in a stellar wind. The spectrum of this ionized component may be flat (if optically thin) or rising (if partially optically thick) with a spectral index perhaps as large as unity (Anglada et al. 1998). Alternatively, the 3.6 cm emission may be attributed to pre-main-sequence magnetic activity, an especially attractive source if the emission is time variable (Feigelson & Montmerle 1999). The same activity could also account for the variable X-ray emission from the star (Kastner et al. 1999). A dotted line in Fig 2 shows an extrapolation of the spectrum for a wind with spectral index unity, which maximally contributes to the emission at shorter wavelengths; even in this case, the hot plasma contribution at 7 mm is 1 mJy, less than 15% of the dust emission from the disk models.
### 4.2 Disk Structure
The resolved 7 mm images are sensitive to the degree of central concentration of the disk emission. The brightness in each beam is given by the product of temperature and opacity, where the latter quantity is proportional to the surface density when the optical depth is low. With the available sensitivity, the disk emission can be detected only where temperatures and optical depths are sufficiently high. The images offer hints of inhomogeneities in the disk emission, in particular along a southeast-northwest axis, but these features have marginal significance (less than twice the rms noise). We consider only axisymmetric models for the disk in deriving structural parameters. Figure 3 shows images of a series of four disk models that follow the standard power law description and match the TW Hya spectrum but have different surface density distributions. To mimic observations, the models have been imaged from the $`(u,v)`$ tracks obtained for TW Hya for two resolutions and deconvolved with the standard algorithms. The models in Figure 3 show that disks with steep surface density distributions produce detectable emission at small radii at high resolution while flatter distributions do not. The disks in these models are oriented face on and have outer radii of 100 AU, but disks with modest inclinations and outer radii from a few 10โs to a few 100โs of AU exhibit qualitatively similar behavior.
Comparison of the model images in Figure 3 with the images in Figure 1 suggest that values of $`p`$ as large as $`1.5`$ or as low as $`0.5`$ are not compatible with the 7 mm data. A model with $`p`$ as low as $`0.5`$ can more closely match the observations if the disk is made sufficiently small, with outer radius $`<50`$ AU, which raises the surface density and the resulting brightness. But the radial extent of scattered light in the Hubble Space Telescope NICMOS and WFPC2 images of TW Hya show that disk material extends well beyond this radius (Weinberger et al. 1999, Krist et al. 2000). The steeper power laws are especially in conflict with the observations if any of the emission visible at the highest resolution can be attributed to an ionized component. The preferred value of $`p`$ is near unity, where the power law model matches the 7 mm observations imaged at low and high resolution. This result requires that the dust properties within the disk are not a strong function of radial distance from the star.
The TW Hya spectrum and resolved 7 mm images are generally consistent with an irradiated gaseous accretion disk (DโAlessio et al. 1997, 1998, Chiang & Goldreich 1997) that follows the ShakuraโSunyaev $`\alpha `$ prescription. For steady accretion, the surface density is given by $`\mathrm{\Sigma }=\dot{M}/3\pi \nu `$ away from the boundaries, where $`\nu =\alpha c_sH`$ is the kinematic viscosity parameterized by a local velocity (the sound speed, $`c_s`$) and scale length (the scale height, $`H`$) and a dimensionless parameter ($`\alpha `$). At radii where the disk is optically thin to its own radiation, a condition that likely holds beyond a few AU for the TW Hya system, the disk becomes nearly vertically isothermal. For a midplane temperature distribution characterized by $`T_mr^{1/2}`$ due to irradiation, the surface density distribution can be approximated by $`\mathrm{\Sigma }(r^{3/2}T_m)^1r^1`$, in good agreement with the data. While inhomogeneities are likely present in the disk, and changes in disk composition and optical depth will modify the energy balance and structure close to the star, the overall structure of the TW Hya disk appears amenable to this simple power law description. Substituting the numerical factors in the expression for surface density gives $`\mathrm{\Sigma }(r)240(\dot{M}/10^8\mathrm{M}_{}\mathrm{yr}^1)(\alpha /0.01)^1(r/\mathrm{AU})^1`$ g cm<sup>-2</sup>, where the scalings for $`\alpha `$ and $`\dot{M}`$ are typical values for pre-main-sequence accretion disks (Hartmann et al. 1998). For the TW Hya disk, the surface density at 1 AU radius in the $`p=1`$ model shown in Figure 3 is 450 g cm<sup>-2</sup>, about a factor of two larger than the canonical value. If the accretion model is valid, then the agreement between the determinations of the surface densityโ within a factor of a fewโ provides support for the accuracy of the adopted mass opacity law at millimeter wavelengths. More stringent tests will be provided by independent measures of the accretion rate and resolved observations of the dust emission at additional wavelengths.
One consequence of the modest brightness observed within 5 AU of TW Hya is that the surface density is unlikely be very much higher in this region of the disk than expected from an $`r^1`$ extrapolation inward. A region of low viscosity and concomitant high surface density in the inner disk may be required by some mechanisms for accretion, planet formation, and also planet migration. For example, the layered accretion of Gammie (1996) piles up accreting mass in a โdead zoneโ where the magnetohydrodynamic instability does not operate very effectively. For TW Hya, there is no evidence for any such substantial mass reservoir at small radii, at least not at the present epoch.
## 5 Conclusions
The TW Hya system provides perhaps the closest analog to the early solar nebula, and VLA observations with subarcsecond resolution characterize conditions in the circumstellar environment at the size scales of giant planet formation. The images of TW Hya at $`\lambda =7`$ mm show extended emission that we attribute to dust in a disk viewed close to pole-on. The dust emission is best fit with a power law disk model having surface density falling off as $`r^1`$, as shallower or steeper power laws produce too little or too much signal at high resolution, especially at radii $`<10`$ AU. The observed structure conforms to that derived in self-consistent calculations of gaseous accretion disks with heating dominated by irradiation from the central star. The observations provide hints of substructure within the disk, but data with much better sensitivity will be required to confirm the reality of these features. This nearby disk system will be a prime target for the next generation of millimeter and submillimeter arrays.
This research was partially supported by NASA Origins of Solar Systems Program grant NAG5-8195. LFR acknowledges the support of CONACyT, Mexico We thank Lee Hartmann and Nuria Calvet for valuable discussions about irradiated accretion disks.
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# A critical-density closed Universe in Brans-Dicke Theory
## 1 Introduction
In the Friedman-Robertson-Walker (FRW) model, there are enough evidences that the Universe is flat, with a large component of negative pressure. This component was consider at first to be just the cosmological constant (or vacuum energy) (Tuner et al. 1984; Peebles 1984; Ostriker & Steinhard 1995; Liddle et al. 1996). Another possibility was to consider topological deffects (Vilenkin 1984), but, the one which has receiving a great deal of attention today, is related to a scalar field, $`Q`$, the so-called โquintessenceโ model or the QCDM model (Caldwell, Dave & Steinhardt 1998). This scalar field is characterized by a very negative pressure, i.e., $`P_Q=w\rho _Q`$, with $`w1/3`$.
Measurements of distant SNe Ia (at $`z1`$) indicate that the expansion of the Universe is in accelaration rather than deceleration (Perlmutter et al. 1998).
According to Garnavich et al. (1998) this is consistent with the existence of an unknown component for the energy density, which could be considered to correspond to the quintessence component. On the other hand, a test of the standard model, including spacetime geometry, galaxy peculiar velocities, structure formation, and early Universe physics, supports, in many of these cases, a flat Universe model with the presence of a cosmological constant (Peebles 1998). Specially, a model in which a short period of inflation occurs at a very early epoch in the evolution of the Universe.Most of these models predict that the total density parameter, $`\mathrm{\Omega }_o`$, be unity.
Given the idea that the Universe could be described by a flat geometry, an interesting question to ask is whether this flatness could be due to a local effect. This sort of question has been considered in the literature (\[Kamionkowski & Toumbas, 1996\]). There, a $`k=1`$ model was taken into account, together with a total density parameter corresponding to $`\mathrm{\Omega }_o<1`$. In this model, openness is obtained by considering a matter component whose equation of state is $`p={\displaystyle \frac{1}{3}}\rho `$. This sort of state of Equation has been reported in models in which topological defects, such as texture or tangled strings, are important components of the total energy of the Universe.
In this paper we analyse a FRW closed model, $`k=1`$, filled with dust and a effective density energy component characterized by a negative pressure, using the Brans-Dicke (BD) theory (\[Brans & Dicke, 1961\]) with a potential associated to the BD field. In particular, we investigate the conditions in order to have a model with positive curvature which mimics a flat universe at low redshift. This imply to determine the contributions of the scalar field $`Q(t)`$ and the BD field to the total energy density, which cancel the contribution due to curvature $`\mathrm{\Omega }_{ok}=1/a_0H_0`$. We should note that the obtained model is the generalization of the Einstein-de Sitter model ($`k=0`$) in the context of the BD theory. Our model is far from being realistic, since it gives an age for the universe,$`t_0`$, very close to $`2H_0/3`$, and a deceleration parameter, $`q_0`$, very close to $`1/2`$. These values, as we mentioned above, desagree with the measurements of distant supernovae. Nevertheless, we clearly pointed out that our proposal is to set the conditions under a closed universe mimics a flat universe in the BD theory. At the end of the following section we shall briefly discuss how in the BD theory a flat universe with negative pressures, presents acceleration in agreement with actual observations.
We studied this sort of model in a previous paper (\[Cruz, del Campo & Herrera, 1998\]), but that analysis was limited to the modelโs intrinsic characteristics, such that the explicit determination of the scalar fields, $`Q(t)`$, its potential $`V(Q)`$ and the BD potential $`V(\mathrm{\Phi })`$, where $`\mathrm{\Phi }`$ represents the BD scalar field. In this paper we investigate the modelโs cosmological characteristics, such that, the proper distance to the horizon, the luminosity distance, the angular size, the differential number of galaxies, and the ratio $`(1/z)\delta z/\delta \theta `$, where, all of these quantities are given as a function of the redshift $`z`$. We compare these parameters with that corresponding to the Einstein-de Sitter model.
We should note that our model, in the limit $`\omega _0\mathrm{}`$ and $`\mathrm{\Phi }=const`$, gives rise to the $`k=0`$ Einstein-de Sitter model.
## 2 Characteristics of the model
Assuming homogeneity and isotropy, the FRW metric for a closed universe is
$$ds^2=dt^2a(t)^2d\mathrm{\Omega }_{k=1}^2,$$
(1)
with $`d\mathrm{\Omega }_{k=1}^2`$ representing the spatial line element asociated to the hypersurfaces of homogeneity, corresponding to a three sphere. $`a(t)`$ represents the scale factor, which together with the assumption that the $`Q`$ scalar field is homogeneous,i.e., $`Q=Q(t)`$, we obtain the fundamental field Equations of the BD model, given by (where the dots representing derivatives with respect to time $`t`$. We use units in which $`c=\mathrm{}=1`$ )
$$H^2+H\left(\frac{\stackrel{}{\mathrm{\Phi }}}{\mathrm{\Phi }}\right)=\frac{\omega _o}{6}\left(\frac{\stackrel{}{\mathrm{\Phi }}}{\mathrm{\Phi }}\right)^2+$$
$$\frac{8\pi }{3\mathrm{\Phi }}\left(\rho _M+\rho _Q\right)\frac{1}{a^2}+\frac{U\left(\mathrm{\Phi }\right)}{6\mathrm{\Phi }},$$
(2)
$$\stackrel{}{\mathrm{\Phi }}+3H\stackrel{}{\mathrm{\Phi }}+\frac{\mathrm{\Phi }^3}{2\omega _o+3}\frac{d}{d\mathrm{\Phi }}\left(\frac{U(\mathrm{\Phi })}{\mathrm{\Phi }^2}\right)=$$
$$\frac{8\pi }{2\omega _o+3}\left[\rho _M+\left(13w\right)\rho _Q\right],$$
(3)
and
$$\stackrel{}{Q}+3H\stackrel{}{Q}=\frac{V(Q)}{Q},$$
(4)
The condition describing a model mimicing a flat Universe is (Cruz, del Campo & Herrera 1998) given by
$$\frac{d}{d\mathrm{\Phi }}\left(\frac{U\left(\mathrm{\Phi }\right)}{\mathrm{\Phi }^2}\right)\frac{1}{2\left(\omega _o+1\right)}\frac{U\left(\mathrm{\Phi }\right)}{\mathrm{\Phi }^3}$$
$$+\frac{3}{\mathrm{\Phi }_0^2(1+\omega _o)}\left(\frac{\mathrm{\Phi }_0}{\mathrm{\Phi }}\right)^{2\left(2+\omega _o\right)}=0,$$
(5)
where $`\mathrm{\Phi }_0`$ is the the actual value of the BD scalar field. Under these conditions Equation (2) and Equation (3)reduce to
$$H^2+H(\stackrel{}{\mathrm{\Phi }}/\mathrm{\Phi })=\frac{\omega _o}{6}(\stackrel{}{\mathrm{\Phi }}/\mathrm{\Phi })^2+\frac{8\pi }{3\mathrm{\Phi }}\rho _M,$$
(6)
and
$$\stackrel{}{\mathrm{\Phi }}+3H\stackrel{}{\mathrm{\Phi }}=\frac{8\pi }{2\omega _o+3}\rho _M,$$
(7)
respectivelly. Equation (4) remains unaltered. Note that this set of equations mimics a flat ($`\mathrm{\Omega }_m=\mathrm{\Omega }_0=1`$) Universe.
Assuming that the matter content, $`\rho _M`$, is dominated by dust, i.e. $`\rho _Ma^3`$, the solutions of Equations (6) and (7) are the well known power law solutions of a flat BD model, given by
$$a=a_o\left(\frac{t}{t_o}\right)^{\frac{2}{3}\left(\frac{\omega _o+1}{\omega _o+4/3}\right)},$$
(8)
and
$$\mathrm{\Phi }(t)=\mathrm{\Phi }_o\left(\frac{t}{t_o}\right)^{\left(\frac{2}{3\omega _o+4}\right)},$$
(9)
where $`t_o`$ is the current age of the Universe. Using the usual formula for the redshift $`z`$, $`1+z=a_o/a`$, we find that
$$t=\frac{2}{3}\left(\frac{\omega _o+1}{\omega _o+4/3}\right)H_o^1\left(1+z\right)^{\frac{2}{3}\left(\frac{\omega _o+4/3}{\omega _o+1}\right)},$$
(10)
where $`H_o\sqrt{8\pi \rho _M^o/3\mathrm{\Phi }_o}`$, is the current Hubble constant taken as $`H_o=100hkms^1Mpc^1`$, with $`h`$ in the range $`0.6h0.8`$ (Bureau, Mould & Staveley-Smith 1996). From Equation (10) we can obtain the age of the Universe, $`t_o={\displaystyle \frac{2}{3}}\left({\displaystyle \frac{\omega _o+1}{\omega _o+4/3}}\right)H_o^1=\left({\displaystyle \frac{\omega _o+1}{\omega _o+4/3}}\right)t_o^E`$, where $`t_o^E`$ is the Einstein-de Sitter value for the age of the Universe, given by $`\frac{2}{3}H_o^1`$. At first glance, the factor appearing in $`t_o`$ seems to decrease the age of the Universe, due to $`\left({\displaystyle \frac{\omega _o+1}{\omega _o+4/3}}\right)1`$. But, since $`\omega _o500`$ (Reasenberg et al. 1979), this factor will be small, and almost equal to one, and, therefore, we could use $`t_ot_o^E`$.
The growth of the Q-field is given by
$$\left(\frac{\stackrel{}{Q}}{Q}\right)=2H_o\left(\frac{\omega _o+3}{\omega _o+1}\right)\left(1+z\right),$$
(11)
and its present value is then determined by the values of the Hubble constant, $`H_o`$, and the BD-parameter $`\omega _o`$. At large values of $`\omega _o`$ $`(\omega _o\mathrm{})`$ it becomes $`(\stackrel{}{Q}/Q)_o=2H_o`$. On the other hand, the potential $`V(Q)`$ associated with this field is given by
$$V(Q)=M^4\left(\frac{M}{Q}\right)^{4\left(\frac{\omega _o+\frac{1}{2}}{\omega _o+1}\right)}$$
(12)
where $`M`$ is a free parameter.
Since $`\rho _Q={\displaystyle \frac{3}{2}}\left({\displaystyle \frac{\omega _o+1}{\omega _o+5/4}}\right)V(Q)`$ and the actual density parameter $`\mathrm{\Omega }_Q^{\omega _o}`$ is defined by $`\mathrm{\Omega }_Q^{\omega _o}={\displaystyle \frac{8\pi \rho _Q^o}{3\mathrm{\Phi }_oH_o^2}}`$, we find that $`\mathrm{\Omega }_Q^{\omega _o}=\left({\displaystyle \frac{\omega _o+1}{\omega _o+5/4}}\right){\displaystyle \frac{4\pi M^4}{\mathrm{\Phi }_oH_o^2}}=\left({\displaystyle \frac{\omega _o+1}{\omega _o+5/4}}\right)\mathrm{\Omega }_Q`$, where $`\mathrm{\Omega }_Q`$ corresponds to the scalar field density parameter in Einstein-de Sitterโs model of the universe.
If we consider the same range for $`\mathrm{\Omega }_Q`$ in our model, we find that this imposes a constrain on $`M`$. Since today it is required that $`V(QM_p\mathrm{})\rho _m10^{47}GeV^4`$, where $`M_p\mathrm{}`$ is the Planck mass, we obtain
$$M\left(\rho _mM_P\mathrm{}^{4\left(\frac{\omega _o+\frac{1}{2}}{\omega _o+1}\right)}\right)^{\frac{\omega _o+1}{8\left(\omega _o+\frac{3}{4}\right)}}.$$
(13)
This gives $`M10^4GeV`$ for $`\omega _o500`$. This value is comparable to particle physics scales. We mention that the evolution of $`Q`$ and the values of $`\mathrm{\Omega }_Q`$ today are very insensitive to the initial value of $`Q`$, due to its attractor solution (Zlater, Wang & Steinhardt 1998).
One point that needs to be taken into account in this kind of model is the fraction of the Universe in causal contact. To do so, we employ a comoving observer at coordinate $`(r_o=0,\theta ,\phi )`$ at time $`t`$. A light signal satisfies the geodesic equation of motion $`ds^2=0`$. Therefore, a light signal emitted from $`(r_1,\theta ,\phi )`$ at time $`t=0`$, following the line $`\theta =\varphi =const.`$, will reach the observer at time $`t`$
$$\underset{o}{\overset{t}{}}\frac{dt^{}}{a(t^{})}=\underset{o}{\overset{r_H}{}}\frac{dr}{\left(1r^2\right)^{\frac{1}{2}}},$$
(14)
and because the proper distance to the horizon is
$$d_H(t)=\underset{o}{\overset{r_H}{}}(g_{rr})^{\frac{1}{2}}๐r,$$
(15)
it is found that
$$d_H(z)=a_0\frac{\alpha (\omega _o,H_o)}{(1+z)}\left[1\left(1+z\right)^{\beta \left(\omega _o\right)}\right],$$
(16)
where $`\alpha (\omega _o,H_o)={\displaystyle \frac{2}{a_oH_o}}\sqrt{\left(\omega _o+3/2\right)\left(\omega _o+4/3\right)}/\left(\omega _o+2\right)`$ and $`\beta \left(\omega _o\right)={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\omega _o+2}{\omega _o+1}}\right)`$. In Fig. 1 we have plotted $`d_H(z)`$ as an function of $`z`$. In Einstein-de Sitter model we have nearly obtained the same curve.
Since, $`\omega _o1`$, we can compare $`d_H(z)`$ in the BD theory with the corresponding expression in Einsteinโs theory of general relativity (Einstein-de Sitter model). Expanding Equation (16) up to the first order in $`1/\omega _o`$, yields
$$d_H(z)d_H^E(z)+g(z,H_o)(1/\omega _o)+O(1/\omega _o)^2,$$
(17)
where $`g(z,H_o)`$ becomes defined by
$$g(z,H_o)=\frac{H_o^1}{(1+z)}\left[\frac{7}{3}\left(\frac{1}{\sqrt{1+z}}1\right)+\frac{\mathrm{}n\left(1+z\right)}{\sqrt{1+z}}\right],$$
(18)
and $`d_H^E(z)`$ represents the horizon distance in the Einstein-de Sitter model. This expression presents a maximum for $`z2`$, and there, the difference, $`\mathrm{\Delta }d_H(z)d_H(z)d_H^E(z)`$, computed up to order $`1/\omega _o`$, becomes a maximum. Its value is not significant, since it is less than one percent.
In the following we determine the deceleration parameter $`q_o`$ for our model. This parameter is defined by $`q_o=\left(\ddot{a}/aH^2\right)_o`$. By using the solution given by Eq. (8), we obtain that
$$q_o=\frac{1}{2}\left(\frac{\omega _o+2}{\omega _o+1}\right).$$
(19)
Note that we can write $`q_o=\gamma \left(\omega _o\right)1`$, where $`\gamma \left(\omega _o\right)`$ is the inverse of the exponent in the expression of the scale factor $`a(t)`$, i.e. $`\gamma \left(\omega _o\right)={\displaystyle \frac{3}{2}}\left({\displaystyle \frac{\omega _o+4/3}{\omega _o+1}}\right)`$. If we consider the lower bound for $`\omega _o`$, i.e. $`\omega _o500,`$ we find that the deceleration parameter has a value close to one half, as it should be in the Einstein-de Sitter model
Following the approach done by Uehara and Kim (1982), we may write directly an expression for the deceleration parameter $`q_o`$ given by
$$q_o=ฯต_o+\frac{\omega _o}{3}ฯต_o^2+\frac{1}{2}\left(\frac{\omega _o+3}{\omega _o+3/2}\right)\mathrm{\Omega }_m,$$
(20)
Notice that this factor is related to the present rate of change of the Newtonโs gravitational constant expressed by $`(\stackrel{}{G}/G)_o={\displaystyle \frac{\dot{\mathrm{\Phi }}}{\mathrm{\Phi }}}_0`$, since $`G(\mathrm{\Phi })={\displaystyle \frac{1}{\mathrm{\Phi }}}`$.
The contribution to the deceleration parameter in Equation (20) is small, since its experimental upper limit is given by $`(\stackrel{}{\mathrm{\Phi }}/\mathrm{\Phi })_o10^{10}yr^1`$ (Helling et al. 1983; Dickey, Newhall & Williams 1989; Shapiro 1989). More recent, measurements on white dwarfs, have decreased the upper limit to $`{\displaystyle \frac{\dot{G}}{G}}_0<10^{11}yr^1`$ (Garcรญa-Berro et al. 1995). Still smaller upper limits for this quantity have been reported (Mรผler et al. 1991), where lunar laser-ranging studies of the moonโs Earth orbit yields $`{\displaystyle \frac{\dot{G}}{G}}_0<10^{13}yr^1`$. Since $`ฯต_01`$, we can write, $`q_o{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\omega _o+3}{\omega _o+3/2}}\right)\mathrm{\Omega }_m`$, which becomes exactly one half the result of Einstein-de Sitterโs theory.
At this point we would like to mention that it is possible to describe in the same spirit, the situation in which the model represents an accelerating flat universe, i. e. a model in Brans-Dicke theory in which now $`\mathrm{\Omega }_0=\mathrm{\Omega }_m+\mathrm{\Omega }_\lambda =1`$. But, as we will see, this case becomes quite complicate to handle. We shall postpone the details of these studies for the near future. We shall restrict ourselves here to give a brief description of this situation. In this case the condition (5) becomes
$$\frac{d}{d\mathrm{\Phi }}\left(\frac{U\left(\mathrm{\Phi }\right)}{\mathrm{\Phi }^2}\right)+(13w)\frac{U\left(\mathrm{\Phi }\right)}{\mathrm{\Phi }^4}$$
$$=3(1w)\frac{\lambda }{\mathrm{\Phi }^2}+3(13w)\frac{1}{\mathrm{\Phi }^2a^2(\mathrm{\Phi })}.$$
(21)
If we want to get an explicit form of the Brans-Dicke scalar potential, $`U(\mathrm{\Phi })`$, we need to know the scale factor $`a`$ as a function of the Brans-Dicke scalar field $`\mathrm{\Phi }`$. In order to do this, we consider the set of basic field equations
$$H^2+H(\stackrel{}{\mathrm{\Phi }}/\mathrm{\Phi })=\frac{\omega _o}{6}(\stackrel{}{\mathrm{\Phi }}/\mathrm{\Phi })^2+\frac{8\pi }{3\mathrm{\Phi }}\rho _M+\frac{\lambda }{3},$$
(22)
and
$$\stackrel{}{\mathrm{\Phi }}+3H\stackrel{}{\mathrm{\Phi }}=\frac{8\pi }{2\omega _o+3}\rho _M+\frac{2\lambda }{2\omega _o+3}\mathrm{\Phi },$$
(23)
(together with equation (4)) that can be solved exactly (Uehara and Kim 1982). The solution for $`\lambda >0`$ is given by
$$a(t)=a_o\left[A\mathrm{cosh}\left(\eta \mathrm{}t\right)\frac{4\pi }{\lambda }\right]^{\alpha (\omega _o)}$$
$$\times \left[\frac{B\mathrm{tanh}\left(\frac{1}{2}\eta \mathrm{}t\right)\sqrt{(4\pi /\lambda )^2A^2}}{B\mathrm{tanh}\left(\frac{1}{2}\eta \mathrm{}t\right)+\sqrt{(4\pi /\lambda )^2A^2}}\right]^{\beta (\omega _o)},$$
(24)
where,$`\eta ^2`$, $`A^2`$, $`B`$, $`\alpha (\omega _o)`$ and $`\beta (\omega _0)`$ are given by $`{\displaystyle \frac{2(4+3\omega _o)}{3+2\omega _o}}`$, $`\left[{\displaystyle \frac{4\pi }{\lambda }}\right]]^2{\displaystyle \frac{3}{2\lambda }}{\displaystyle \frac{1}{4+3\omega _o}}[{\displaystyle \frac{\mathrm{\Phi }_o}{\rho _o}}\left]^2\right[(1+\omega _o)ฯต_o+H_o]^2(>0)`$,
$`{\displaystyle \frac{4\pi }{\lambda }}+A`$, $`{\displaystyle \frac{1+\omega _o}{4+3\omega _o}}`$ and $`{\displaystyle \frac{1}{4+3\omega _o}}\sqrt{{\displaystyle \frac{3+2\omega _o}{3}}}`$, respectively. The interval of time, $`\mathrm{}t=tt_c`$, is related to $`t_c={\displaystyle \frac{2}{\eta \sqrt{\lambda }}}\left[{\displaystyle \frac{4\pi /\lambda A}{\frac{4\pi }{\lambda }+A}}\right]`$.
Expressions (24) together with equation (22) allow us to determine the scalar field $`\mathrm{\Phi }`$ as a function of time. Thus, in principle, we could get the scale factor as a function of the Brans-Dicke scalar field.
We should note that the deceleration rate $`q_o`$ becomes in this case
$$q_o=ฯต_o+\frac{\omega _o}{3}ฯต_o^2+\frac{1}{2}\left(\frac{\omega _o+3}{\omega _o+3/2}\right)\mathrm{\Omega }_m\frac{2\omega _o}{2\omega _o+3}\mathrm{\Omega }_\lambda ,$$
(25)
with $`\mathrm{\Omega }_\lambda `$ defined by $`{\displaystyle \frac{\lambda }{3H_o^2}}`$. Notice that if the cosmological constant contribution to the total density parameter is significant, in agreement with the Supernova observations, the deceleration parameter becomes negative and thus the universe will show an acceleration. We stop here this analysis and, as we have mentioned above, we shall describe in more details this situation in near future.
In going back to our simple Einstein-de Sitter-like model in BD theory, we discuss in the following some kinematical properties.
## 3 Kinematics of the model
### 3.1 Luminosity Distance-Redshift
The โluminosity distanceโ is defined as the ratio of the emitted energy per unit time, $``$, and the energy received per unit time $``$
$$d_L^2=\frac{}{4\pi }$$
(26)
In the absence of an expansion, the luminosity distance is simply the physical distance to the source. In closed FRM cosmology, the luminosity distance to a source at coordinate $`(r_1,\theta ,\varphi )`$ is, assuming for convenience that the observer is located at $`r=0`$,
$$d_L^2(z)=a_0^2r_1^2(z)(1+z)^2,$$
(27)
The factor $`a_0^2r_1^2`$ in this expression is nothing but the โinverse square lawโ, since a two-sphere surrounding the source encompassing the observer has an area equal to $`4\pi a_0^2r_1^2(z)`$. The factor $`(1+z)^2`$ appears due to the redshift of the radiation between the time of emission and the time of detection. The parameter $`r_1`$ is determined by the expression
$$\underset{o}{\overset{r_1(z)}{}}\frac{dr}{\sqrt{1r^2}}=\frac{1}{H_o}\underset{o}{\overset{z}{}}\frac{dz}{E(z)},$$
(28)
where $`E(z)`$ is given by
$$E(z)=\frac{\left(1+z\right)^{\gamma \left(\omega _o\right)}}{a_oH_o\beta \left(\omega _o\right)\alpha (\omega _o,H_o)}.$$
(29)
Thus, we find for the Luminosity distance-redshift relation
$$d_L(z)=a_o(z+1)\mathrm{sin}\left\{\alpha (\omega _o,H_o)\left[1(1+z)^{\beta (\omega _o)}\right]\right\}.$$
(30)
Fig. 2 shows how the luminosity distance, $`d_L\left(z\right)`$, changes with redshift parameter $`z`$, for Einstein-de Sitter model (dotted line) and BD closed model (solid line) . From this Figure we observe that the BD-theory begins to differ from Einstein-de Sitterโs theory at $`z100`$. If we take the value for the redshift $`z`$ corresponding to the value associated to the last scattering surface, i.e. $`zz_{LS}1.100`$, the luminosity distance in Einstein-de Sitterโs theory differs approximately from its analogous expression in the BD theory by a few percentage points.
Note that for small $`z`$, Equation (30) yields
$$H_od_L(z)=z+\frac{1}{2}\left(1\beta \left(\omega _o\right)\right)z^2+o(z^3),$$
(31)
where we have rescaled the luminosity distance by a factor of $`a_oH_o\beta \left(\omega _o\right)\alpha (\omega _o,H_o).`$ The first term of this expression is nothing but the Hubble law, and from the second term, we can read the $`q_o`$ parameter, $`q_o=\beta \left(\omega _o\right)={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\omega _o+2}{\omega _o+1}}\right)`$, which coincides with the result obtained above.
Also, from Equation (30) we can use the apparent bolometric magnitude $`m(z)`$ of a standard candle (with $`M`$ equal to the absolute bolometric magnitude) as a function of the redshifts
$$m(z)=M+5log[d_L(z)]+25.$$
(32)
If we define $`D_L(z)H_od_L(z)`$, then we obtain, from the latter Equation
$$m(z)=+5log[D_L(z)],$$
(33)
where $`M5log(H_o)+25`$, is the magnitude โzero pointโ, which can be determined observationally. Hamuy et al. (1996) have determined the value $`=3.17\pm 0.03`$ using 18 supernovae discovered by the Calรกn/Tololo researchers. This value is supposed to be independent of the redshift $`z`$. Therefore, if we take two different values for the redshift, we can obtain a variation of the apparent bolometric magnitude which is given by $`\mathrm{\Delta }m(z_2,z_1)5\mathrm{log}\left[d_L(z_2)/d_L(z_1)\right]`$, or explicitly
$$\mathrm{\Delta }m(z_2,z_1)=5\mathrm{log}$$
$$\left\{\frac{\left(1+z_2\right)\mathrm{sin}\left[\alpha (\omega _o,H_o)\left(1\left(1+z_2\right)^{\beta \left(\omega _o\right)}\right)\right]}{\left(1+z_1\right)\mathrm{sin}\left[\alpha (\omega _o,H_o)\left(1\left(1+z_1\right)^{\beta \left(\omega _o\right)}\right)\right]}\right\}.$$
(34)
This expression could be used to restrict the value of the BD parameter $`\omega _o`$.
By using the values $`z_1=0.5`$ and $`z_2=1.0`$, we find the $`\mathrm{\Delta }m`$ increases monotonically with the BD-parameter $`\omega _o`$. When this parameter reaches a value close to $`\omega _o20`$, the growth of $`\mathrm{\Delta }m`$ with respect to $`\omega _o`$ increases, but slowly now, approaching asymptotically a limiting value. At $`\omega _o=500`$ this difference yields the value $`\mathrm{\Delta }m3.32`$.
On the other hand, if we consider the values specified by Goobar & Perlmutter (1995) in which the apparent magnitude (R-band) was $`m_R=22.17\pm 0.05`$ at redshift $`z=0.5`$ and $`m_R=25.20\pm 0.05`$ at redshift $`z=1.0`$, we obtain $`\mathrm{\Delta }m_R=3.03\pm 0.05`$. This value represents a similar order of magnitude as that obtained for $`\omega _o=500`$.
### 3.2 The Angular diameter-redshift
The angular diameter distance $`d_A`$ between a source at redshift $`z_2`$ and $`z_1<z_2`$ is defined by
$$d_A(z_1,z_2)=\frac{a_o\mathrm{sin}\left[\mathrm{\Delta }\chi (z_1,z_2)\right]}{\left(1+z_2\right)}$$
(35)
where $`\mathrm{\Delta }\chi (z_1,z_2)`$ is the polar-coordinate distance between a source at $`z_1`$ and another at $`z_2`$, in the same line of sight
$$\mathrm{\Delta }\chi (z_1,z_2)=\alpha (\omega _o,H_o)\left[\left(1+z_1\right)^{\beta \left(\omega _o\right)}\left(1+z_2\right)^{\beta \left(\omega _o\right)}\right]$$
(36)
The corresponding angular size of an object of proper length $`\mathrm{}`$ at a redshift $`z`$ results in $`\theta \left(z\right)\mathrm{}/d_A(0,z)`$, which becomes (in units of $`\mathrm{}H_o`$)
$$\mathrm{\Theta }\left(z\right)=\frac{1}{a_oH_o}\frac{\left(1+z\right)}{\mathrm{sin}\left\{\alpha (\omega _o,H_o)\left[1\left(1+z\right)\right]^{\beta \left(\omega _o\right)}\right\}}$$
(37)
Fig. 3 shows the plot of $`\mathrm{\Theta }`$ as a function of $`z`$ in Einstein-de Sitterโs theory and the BD theory with $`\omega _o=500`$. These two theories coincide in the $`0zz_{LS}`$. At $`z=z_{LS}`$, it is found that in the first approximation of $`1/\omega _o`$, their difference is $`\mathrm{\Delta }\mathrm{\Theta }\mathrm{\Theta }_{BD}\mathrm{\Theta }_E147\left(1/\omega _o\right)`$, for any large value of $`\omega _o`$. In this plot, we have taken $`a_0H_0=\sqrt{2/5}`$ \[Kamionkowski & Toumbas, 1996\].
### 3.3 The number count-redshift
The number of galaxies in a comoving volume element in an angular solid area $`d\mathrm{\Omega }`$ with redshift between $`z`$ and $`z+dz`$ is sensitive to the number of galaxies, $`n`$, in a comoving volume element $`dV_c`$ and the spatial curvature:
$$dN_{gal}=ndV_c=n\frac{r^2}{\left(1r^2\right)^{\frac{1}{2}}}drd\mathrm{\Omega }$$
(38)
ยฟFrom this relation we can write the differential number of galaxies per steradian per unit redshift,
$$\frac{dN_{gal}}{dzd\mathrm{\Omega }}\left(z\right)=\frac{na_o^2\mathrm{sin}^2\left[\mathrm{\Delta }\chi \left(z\right)\right]}{H_oE(z)}$$
(39)
Using the expression for $`\mathrm{\Delta }\chi (z)`$ from Equation 36 and $`E(z)`$ from Equation (29), we obtain
$$\frac{dN_{gal}}{dzd\mathrm{\Omega }}\left(z\right)=na_o^3\alpha (\omega _o,H_o)\beta (\omega _o)\times $$
$$\frac{\mathrm{sin}^2\left\{\alpha (\omega _o,H_o)\left[1\left(1+z\right)^{\beta \left(\omega _o\right)}\right]\right\}}{\left(1+z\right)^{\gamma \left(\omega _o\right)}}.$$
(40)
In Fig. 4 we have plotted the number-redshift relation (given by Equation (39)), for $`\omega _o=500`$. Here, we have confirmed that, for a large range of the redshift z, the function of the number of galaxies per unit of steradian per unit redshift in Einstein-de Sitterโs model becomes almost indistinguishable from its similar expression in the BD-theory.
For small redshift one obtaines
$$\frac{dN_{gal}}{dzd\mathrm{\Omega }}\left(z\right)n\left(a_o\alpha (\omega _o,H_o)\beta (\omega _o)\right)^3z^2+O\left(z^3\right).$$
(41)
This expression indicates that, for small redshift, the slope in the plot $`{\displaystyle \frac{dN_{gal}}{dzd\mathrm{\Omega }}}`$ v/s $`z^2`$ could provide some information about the $`\omega _o`$ BD parameter, assuming that the present value of the Hubble constant $`H_o`$ is known. This, in principle, could be checked by the corresponding observations.
### 3.4 The ratio of the redshift thickness-angular size
The redshift thickness $`\delta z`$ and the angular size $`\delta \theta `$ of a spherical structure that grows with the expansion of the Universe have a dependence on $`z`$ given by
$$\frac{1}{z}\frac{\delta z}{\delta \theta }(z)=a_oH_o\frac{E(z)\mathrm{sin}\left[\mathrm{\Delta }\chi (z)\right]}{z},$$
(42)
or explicitly,
$$\frac{1}{z}\frac{\delta z}{\delta \theta }(z)=\frac{\left(1+z\right)^{\gamma (\omega _o)}\mathrm{sin}\left\{\alpha (\omega _o,H_o)\left[1\left(1+z\right)^{\beta (\omega _o)}\right]\right\}}{a_oH_o\alpha (\omega _o,H_o)\beta (\omega _o)z}$$
(43)
Fig. 5 (lower panel) shows how $`{\displaystyle \frac{1}{z}}{\displaystyle \frac{\delta z}{\delta \theta }}`$ changes with redshift $`z`$ in the BD-theory, for $`\omega _o=500`$. We see that, in the range of $`0z10`$, there is no difference when it is compared with its analogous expression in Einstein-de Sitter model. However, at large redshift i.e. $`zz_{LS}1.100`$, their difference becomes significant, as we can see from Fig. 5 (upper panel). Thus, the two expressions differ at large redshift values.
## 4 Conclusions
We have studied a model in which the โquintessenceโ component was included. In this model we have chosen a Universe with closed topology $`\left(k=1\right)`$ in the BD theory, and we have added a scalar potential for the BD-scalar field. The quintessence contribution to the matter component was fine tuned to exactly cancel the curvature term (with the help of the BD scalar potential) in the corresponding field Equations.
Under these conditions, different cosmological expressions were calculated as a function of the redshift $`z`$. For instance, the luminosity distant, the angular diameter, the number count and ratio of the redshift thickness-angular size. All these quantities were compared with their analogous expression related to Einstein-de Sitter model which seems flat at low redshift. In some cases (Luminosity distance and the ratio of the redshift thickness-angular size), the corresponding expressions become distinguishable at a high enough redshift. At that redshift, these differences are difficult to be directly detected . Perhaps, if we consider other observational facts, such as, the anisotropy of the cosmic microwave background radiation, we might say something about these differences, and use these results for establishing a limit for the BD parameter, $`\omega _o`$, since the two theories differ in the value of this parameter.
## acknowledgements
SdC was supported from the COMISION NACIONAL DE CIENCIAS Y TECNOLOGIA through grant FONDECYT $`N^0`$1000305 and also from UCV-DGIP 123.744/00. NC was supported by USACH-DICYT under grant $`N^0`$ 0497-31CM.
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# 1 Introduction
## 1 Introduction
Noncommutativity of coordinates is not a surprising occurrence in physics, quantum phase space being the first example that comes to mind. In fact some early considerations on its โquantizedโ differential geometry can be found in . This particular operator algebra has inspired the idea of spacetime coordinates as noncommuting operators. The idea has been explored since quite some time in various directions, one main motivation being that the relation:
$$[x^\mu ,x^\nu ]=i\theta ^{\mu \nu }$$
(1.1)
embodies an uncertainty principle that smears the spacetime picture at distances shorter than $`\sqrt{\theta }`$, and therefore a natural cutoff when using a quantum field theory to describe natural phenomena. Since โmeasuringโ spacetime geometry under distances smaller than the Planck length $`L_P`$ is not accessible even to Gedanken experiments (at this scale the curvature radius of spacetime becomes of the order of a probe particle wavelength), relation (1.1) seems to make good physical sense when $`\sqrt{\theta }L_P`$. Thus a quantum theory of gravity containing or predicting relation (1.1) would have a good chance to be intrinsically regulated.
String theories have been pointing towards a noncommuting scenario already in the 80โs . More recently Yang-Mills theories on noncommutative spaces have emerged in the context of $`M`$-theory compactified on a torus in the presence of constant background three-form field , or as low-energy limit of open strings in a backround $`B`$-field -, describing the fluctuations of the $`D`$-brane world volume . As observed for example in , noncommutativity in open string theories is to be expected at some level, since open string vertex operators are inserted along a one-dimensional line, i.e. the boundary of the world sheet: the points of insertion are canonically ordered, so that the product of two such operators depends on their order of insertion. For a comprehensive account on noncommutativity in string theory and M-theory we refer to D. Bigattiโs lectures , and to earlier reviews (for ex. ).
The first part of this review concerns a short description of noncommutative Yang-Mills theories, with emphasis on the algebraic structure, that is on the (noncommutative) Moyal product, and with some remarks on the relations between deformed products and quantization rules. Recent results on perturbative aspects of noncommutative scalar field theories are recalled.
The second part is devoted to the differential geometry of finite groups. The general theory is illustrated in the case of $`Z_2`$. As a physical application, we construct a gauge theory on $`M_4\times Z_2`$, obtaining via a Kaluza-Klein mechanism a Higgs field (with the correct spontaneous symmetry-breaking potential and Yukawa couplings) in $`d=4`$ Minkowski spacetime $`M_4`$.
Noncommutative geometry (NCG) has a vast literature that we do not even attempt to cite. Reviews can be found in . We just mention some of its uses in physics not discussed in these two lectures: Connesโ program of reconstructing the standard model from the NCG of suitable operator algebras ; quantum groups, i.e. continuous deformations of Lie groups, and their NCG applied to gauge and gravity theories (see for ex. ), deformed quantum mechanics and solid state physics.
To find the geometry corresponding to a given algebraic structure is a fascinating and usually difficult task, whereas the inverse route is often much easier. A constructive starting point for NCG is to reformulate as much as possible the geometry of a manifold in terms of an algebra of functions defined on it <sup>1</sup><sup>1</sup>1For example tangent vectors on a manifold $`V`$ can be seen as derivations on the functions on $`V`$, etc. , and then to generalize the corresponding results of differential geometry to the case of a noncommutative algebra of functions. The main notion which is lost in this generalization is that of a point (โnoncommutative geometry is pointless geometryโ).
## 2 From sets of points to algebras of functions: $`C^{}`$ algebras
The primordial arena for geometry and topology are sets $`V`$ of points with some particular structure. Such a set we call โspaceโ. In many cases this set is completely characterized by an algebra of functions on it, so that all the information about $`V`$ can be retrieved from the functions alone.
Let us start with an elementary example: a finite dimensional vector space $`V`$. The functionals
$$f:Vor$$
(2.1)
constitute the dual vector space $`V^{}`$ isomorphic to $`V`$, a basis in $`V^{}`$ being given by the functionals $`x^i`$, dual to the basis vectors $`v_j`$ of $`V`$: $`x^i(v_j)=\delta _j^i`$. The study of $`V^{}`$ is completely equivalent to the study of $`V`$.
More generally consider a set $`V`$ of points, and the algebra of complex valued functions on $`V`$, $`A=Fun(V)`$. This algebra is clearly associative and commutative, with the usual pointwise product and sum: $`(fg)(v)=f(v)g(v),(f+g)(v)=f(v)+g(v),(\lambda f)(v)=\lambda f(v),\lambda `$. The unit $`I`$ of the algebra is given by the function $`I(v)=1,vV`$. As a simple example suppose again that $`V`$ has a finite number of elements. Then $`A`$ is of finite dimension as a vector space, and any $`fA`$ is expressible as
$$f=f_ix^i,x^i(v_j)=\delta _j^i$$
(2.2)
where now $`v_j`$ are the elements of $`V`$. Note the multiplication rule:
$$x^ix^j=\delta ^{ij}x^i$$
(2.3)
A norm can be defined in $`A`$ : $`fmax_{vV}|f(v)|`$. Let $`f^{}`$ be the complex conjugate of $`f`$, then
$$ff^{}=f^2$$
(2.4)
A normed algebra with an involution $`ff^{}`$ satisfying (2.4) is called a $`C^{}`$ algebra. Thus $`A=Fun(V)`$ is a (commutative) $`C^{}`$ algebra.
Conversely any $`n`$-dimensional commutative $`C^{}`$ algebra can be considered as algebra of functions on a set of $`n`$ points. Note that commutativity is essential to interpret it as an algebra of functions on a set of points.
The finite dimensional example extends to infinite sets if they have a topology. In fact if $`V`$ is a compact space, then the algebra $`C^{}(V)`$ of continuous functions on $`V`$ is a $`C^{}`$ algebra. Conversely any $`C^{}`$ algebra $`A`$ with a unit element is isomorphic to the algebra of continuous complex functions on some compact space $`V`$. This space is just the space of homomorphisms $`\chi `$ from $`A`$ to $`C`$ such that $`\chi (I)=1`$. The points of $`V`$ are then in 1-1 correspondence with irreducible representations of $`A`$. This is essentially the commutative Gelโfand-Naimark theorem.
Replacing now the commutative $`A`$ with a noncommutative $`A`$, the โspaceโ may be hard to find: in most cases these algebras have non nontrivial homomorphisms into $``$, so that the reconstruction of a space fails. But the existence of such a space may not be necessary, if one has transferred all the relevant information for a physical theory into the algebra $`A`$.
There are various ways to generalize to the noncommuting case. Continuous deformations of commutative $`A`$ into noncommutative $`A`$ include quantum groups (and quantum coset spaces) and deformations of Poisson structures, of which the noncommutative torus is a simple example. In these cases there is a set of continuous parameters that control the noncommutativity, and one recovers the commuting case (the โclassical limitโ) for some specified values of these parameters.
On the other hand there are noncommutative algebras that are not connected to a commutative limit, as in the case of matrices with entries in $`Fun(V)`$. An example that we will work out in some detail in Section 4 is the differential geometry of finite groups: in this case $`Fun(V)`$ is commutative, but the differentials do not simply anticommute between themselves and do not commute with functions: hence a noncommutative differential geometry.
## 3 Deformation quantization
Consider the algebra of smooth functions on phase space. Deformation quantization essentially consists in deforming the usual commutative product between functions into an associative noncommutative product, the โstarโ product:
$$AB=AB+i\frac{\mathrm{}}{2}\{A,B\}_{PB}+0(\mathrm{}^2)$$
(3.1)
where $`\mathrm{}`$ is a parameter ($`\mathrm{}0`$ corresponds to the commutative limit), and $`\{A,B\}_{PB}`$ is the Poisson bracket of the two phase space functions $`A(q,p),B(q,p)`$. Imposing associativity of the star product determines the higher $`0(\mathrm{}^2)`$ terms up to equivalences that we discuss in next paragraph.
More generally, on a given manifold $`X`$ with a Poisson structure there is essentially one star product, modulo gauge equivalences that amount to linear redefinitions of the functions:
$$AD(\mathrm{})AA+\mathrm{}D_1(A)+\mathrm{}^2D_2(A)+\mathrm{},$$
(3.2)
$`D_i:Fun(X)Fun(X)`$ being differential operators. This result was proved, in the sense of formal series expansions, in ref.s . That the linear automorphisms (3.2) are gauge transformations with respect to the star product can be understood as follows: consider the deformation of the ordinary product $`AB`$ due to $`D(\mathrm{})`$:
$$AB=D(\mathrm{})^1(D(\mathrm{})(A)D(\mathrm{})(B))$$
(3.3)
This product is still commutative, and not essentially different from the ordinary one. Two $``$ products related by $`D(\mathrm{})`$ may therefore be considered equivalent. Thus deformation quantization yields a noncommutative algebra of functions for each Poisson structure on the manifold $`X`$. Poisson structures $`\{,\}`$ can be parametrized by an antisymmetric tensor $`\theta ^{ij}(x)`$ such that $`\{A,B\}\theta ^{ij}(x)(_iA)(_jB)`$, satisfying differential identities corresponding to the Jacobi identities of the Poisson bracket. The simplest Poisson structure is of course the Poisson bracket of ordinary (flat) phase-space, whose noncommutative algebra we consider in the following.
Historically the deformations (3.1) arose in studying the noncommutative structure of quantum mechanics, and this explains the word โquantizationโ and the appearance of the symbol $`\mathrm{}`$ as deformation parameter. Consider for example the Weyl quantization rule $`W`$ (a linear map from the classical phase-space functions to the quantum operators) of the basic phase space monomial:
$$q^mp^nW(q^mp^n)=\frac{1}{2^n}\underset{k=0}{\overset{n}{}}\left(\begin{array}{c}n\\ k\end{array}\right)\widehat{p}^{nk}\widehat{q}^m\widehat{p}^k$$
(3.4)
where $`\widehat{q},\widehat{p}`$ are the quantum phase space operators. This rule amounts to sum on the permutations of all $`\widehat{p}`$ and $`\widehat{q}`$ considered as different objects, thus producing an hermitian operator. For example
$$W(qp^2)=\frac{1}{4}(\widehat{p}^2\widehat{q}+2\widehat{p}\widehat{q}\widehat{p}+\widehat{q}\widehat{p}^2)$$
(3.5)
Note that this rule can be efficiently restated as
$$W(q^mp^n)=\left[\mathrm{exp}[\frac{1}{2}i\mathrm{}(\frac{^2}{qp})]q^mp^n\right]_{q\widehat{q},p\widehat{p}}$$
(3.6)
where the substitution $`q\widehat{q},p\widehat{p}`$ occurs on each monomial $`q^rp^s`$ with $`q`$โs ordered to the left. This formula may be checked to hold on the basic monomial, and extends therefore to any phase-space function $`A(q,p)`$ expressible as a power series:
$$W(A(q,p))=:\mathrm{exp}[\frac{1}{2}i\mathrm{}(\frac{^2}{qp})]A(q,p):$$
(3.7)
$`::`$ indicating normal ordering ($`q`$ preceding $`p`$) and substitution $`q\widehat{q},p\widehat{p}`$. The map $`W`$ is invertible, i.e. there is a 1-1 correspondence between quantum operators and functions on phase-space. This is essentially the core of Moyal formalism , enabling the study of quantum systems within the classical arena of phase-space via the inverse map $`W^1`$.
Consider the product of two quantum operators $`W(A)`$, $`W(B)`$: the classical image $`W^1`$ of their product is what is called the Moyal product $`AB`$, and is given by <sup>2</sup><sup>2</sup>2In fact the product was introduced by H. Groenewold (and even earlier, less explicitly, by J. von Neumann ). We thank C. Zachos for bringing this to our attention.
$$W^1(W(A)W(B))AB=A(q,p)\mathrm{exp}[i\frac{\mathrm{}}{2}\mathrm{}]B(q,p)$$
(3.8)
where $`\mathrm{}`$ is the bidifferential operator defining the Poisson bracket:
$$A\mathrm{}B\{A,B\}_{PB}$$
(3.9)
i.e. $`\mathrm{}=(\frac{\stackrel{}{}}{q}\frac{\stackrel{}{}}{p}\frac{\stackrel{}{}}{p}\frac{\stackrel{}{}}{q})`$. Clearly the Moyal product inherits the properties of the operator product, i.e. it is associative and noncommutative (unless the operators $`W(A),W(B)`$ happen to commute), and gives an explicit instance of the star product (3.1).
The Moyal bracket $`\{A,B\}_M`$ is given by the commutator:
$$\{A,B\}_MABBA=2iA\mathrm{sin}[\frac{\mathrm{}}{2}\mathrm{}]B$$
(3.10)
and obviously has all the properties of a Lie bracket: it is bilinear, antisymmetric and satisfies Jacobi identities. The Moyal bracket is the image in classical phase-space of the commutator between quantum operators:
$$\{A,B\}_M=W^1([W(A),W(B)])$$
(3.11)
cf. eq. (3.8).
Of course the Weyl map is not the only possible quantization rule. A classification of quantization rules and the construction of the corresponding noncommutative $``$ products and brackets can be found in . In fact different quantization rules correspond to $``$ products connected by the gauge transformations (3.2).
Similarly we can introduce noncommutativity in ordinary $`^d`$ spacetime via a new product on the $`C^{}`$ algebra of $`C^{\mathrm{}}`$ complex functions:
$$AB(x)A(x)\mathrm{exp}[\frac{i}{2}\underset{\mu }{\overset{}{}}\theta ^{\mu \nu }\underset{\nu }{\overset{}{}}]B(x)$$
(3.12)
where $`\theta ^{\mu \nu }`$ is constant, real and antisymmetric. Then the commutator of the coordinates $`x^\mu `$ computed with the star product yields precisely relation (1.1). By a change of coordinates $`\theta `$ can be reduced to the symplectic form:
$$\theta =\left(\begin{array}{ccc}\hfill 0& \hfill 1& \\ \hfill 1& \hfill 0& \\ & & \hfill \mathrm{}\end{array}\right)$$
(3.13)
Thus if $`\theta `$ has rank $`r`$ the relations (1.1) describe a spacetime with $`\frac{r}{2}`$ pairs of noncommuting coordinates (with the same algebraic structure as an $`r`$-dimensional phase-space) and $`dr`$ coordinates that commute with all the others. In the $`r`$-dimensional subspace the star product coincides with the Moyal product discussed previously.
A noncommutative torus is obtained by considering periodic coordinates $`0x^\mu <2\pi `$. In the periodic case it is convenient to redefine the star product (3.12) as $`AB(x)A(x)\mathrm{exp}[\pi i\underset{\mu }{\overset{}{}}\theta ^{\mu \nu }\underset{\nu }{\overset{}{}}]B(x)`$ (which amounts to multiply $`\theta `$ by $`2\pi `$), and to change variables:
$$U_\mu e^{ix^\mu }$$
(3.14)
The product between these new variables becomes:
$$U_\mu U_\nu =e^{\pi i\theta ^{\mu \nu }}e^{i(x^\mu +x^\nu )}=e^{2\pi i\theta ^{\mu \nu }}U_\nu U_\mu $$
(3.15)
Notice that two noncommutative tori related by $`\theta ^{\mu \nu }\theta ^{\mu \nu }+\mathrm{\Lambda }^{\mu \nu }`$, where $`\mathrm{\Lambda }^{\mu \nu }`$ is antisymmetric with integer entries, are equivalent.
Quantum field theories on noncommutative spacetime (for a very partial list of ref.s see -) are then obtained by considering their ordinary action and replacing the usual product between fields with the $``$ product. Indeed the algebra of functions on noncommutative $`R^d`$ can be viewed as the algebra of ordinary functions on the usual $`R^d`$ with a deformed $``$ product. Thus we transfer the noncommutativity of spacetime to the noncommutativity of the product between functions, and then apply the usual perturbation theory. Because of the nonpolynomial character of the star product the resulting field theory is nonlocal. This kind of theories is under active study. Weโll mention here only a few results.
Noncommutative scalar theories at the perturbative level have been investigated for example in . The quadratic part of the action is the same as in the noncommutative theory, since $`d^dx\varphi \varphi =d^dx\varphi \varphi `$ and likewise for the kinetic term (assuming suitable boundary conditions on $`\varphi `$ that allow to drop total derivatives). Therefore propagators are the same as in the commutative case. The interactions however are modified: in momentum space an interaction vertex $`\varphi ^n`$ gives rise to an additional phase factor:
$$V(k_1,k_2,\mathrm{},k_n)=e^{\frac{i}{2}_{i<j}k_i\times k_j}$$
(3.16)
where $`k_i`$ is the momentum flowing into the vertex through the $`i`$th $`\varphi `$ and $`k_i\times k_j(k_i)_\mu \theta ^{\mu \nu }(k_j)_\nu `$. This is the only modification to the Feynman rules and its consequences have been investigated in , finding that $`\theta `$ dependence factorizes in planar graphs (i.e. the phase factor associated with the planar diagram does not contain internal momenta), yielding then essentially the same behaviour as in the noncommutative theory. Interesting differences arise in the non-planar diagrams: the one-loop diagrams turn out to be finite at generic values of the external momenta, a consequence of the rapid oscillations of phase factors of the type $`e^{ip\times k}`$ where $`p`$ is an external momentum and $`k`$ is the loop momentum. These factors disappear when $`p_\mu \theta ^{\mu \nu }0`$, and the nonplanar graphs become singular in this limit. This can be interpreted as a mixing between UV and IR divergences: turning on $`\theta `$ replaces the UV divergence with a $`p0`$ IR divergence. Moreover, the commutative limit $`\theta 0`$ is not smooth .
The quantum analysis is more complicated for noncommutative Yang-Mills theories. Classically the noncommutative $`U(N)`$ YM action is :
$$S=\frac{1}{4g^2}Tr(F_{\mu \nu }F_{\mu \nu })$$
(3.17)
where
$`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu i(A_\mu A_\nu A_\nu A_\mu )`$ (3.18)
$`A=A^at^a,Tr(t^at^b)=\delta ^{ab}`$ (3.19)
As already noticed in , the noncommutative gauge transformations:
$`\delta _\epsilon A_\mu =_\mu \epsilon i(A_\mu \epsilon \epsilon A_\mu )`$ (3.20)
$`\delta _\epsilon F_{\mu \nu }=i(F_{\mu \nu }\epsilon \epsilon F_{\mu \nu })`$ (3.21)
leave the action invariant. The perturbative quantum theory is the object of current research.
The extension of the AdS/CFT correspondence to backgrouds with constant $`B`$ field can shed some light on the nonperturbative regime of noncommutative field theories, see for ex. ref. .
Deformation quantization has been applied to numerous other physical systems, besides scalar and gauge theories and their supersymmetric versions . We mention for example gravity and the bosonic string action .
## 4 Dynamics on finite groups from their noncommutative geometry
In this Section we present a systematic method for constructing field theories on finite groups. This method is based on the (noncommutative) differential geometry of finite groups, studied in ref.s . The general theory is applied to the simplest possible finite group, i.e. $`Z_2`$. The example of the simplest nonabelian finite group $`S_3`$ can be found in , and in a gravity-like theory on $`S_3`$ is discussed. Here we will use the NCG on $`Z_2`$ to formulate a $`U(N)`$ gauge theory coupled to Dirac fermions on $`M_4\times Z_2`$, yielding in $`M_4`$ (Minkowski spacetime) a Yang-Mills theory coupled to Dirac matter plus a Higgs sector with symmetry-breaking potential and Yukawa couplings to the fermions.
Differential calculi can be constructed on spaces that are more general than differentiable manifolds. Indeed the general algebraic construction of differential calculus in terms of Hopf structures allows to extend the usual differential geometric quantities (connection, curvature, metric, vielbein etc.) to a variety of interesting spaces that include quantum groups, noncommutative spacetimes (i.e. quantum cosets), and discrete spaces.
In this lecture we concentrate on the differential geometry of finite group โmanifoldsโ. As discussed in , these spaces can be visualized as collections of points, corresponding to the finite group elements, and connected by oriented links according to the particular differential calculus we build on them. Although functions $`fFun(G)`$ on finite groups $`G`$ commute, the calculi that are constructed on $`Fun(G)`$ by algebraic means are in general noncommutative, in the sense that differentials do not commute with functions, and the exterior product does not coincide with the usual antisymmetrization of the tensor product.
Among the physical motivations for finding differential calculi on finite groups we mention the possibility of using finite group spaces as internal spaces for Kaluza-Klein compactifications of Yang-Mills, supergravity or superstring theories ( for example Connesโ reconstruction of the standard model in terms of noncommutative geometry can be recovered as Kaluza-Klein compactification of Yang-Mills theory on an appropriate discrete internal space). Differential calculi on discrete spaces can be of use in the study of integrable models, see for ex. ref. . Finally gauge and gravity theories on finite group spaces may be used as lattice approximations. For example the action for pure Yang-Mills $`F{}_{}{}^{}F`$ considered on the finite group space $`Z^N\times Z^N\times Z^N\times Z^N`$, yields the usual Wilson action of lattice gauge theories, and $`N\mathrm{}`$ gives the continuum limit . New lattice theories can be found by choosing different finite groups.
While the construction of the differential calculus on finite groups in ref.s uses the Hopf algebraic approach of Woronowicz , here this calculus will be presented without recourse to Hopf algebra techniques. Most of the content of next Section can be found in , and is included here for self-consistency.
### 4.1 Differential calculus on finite groups
Let $`G`$ be a finite group of order $`n`$ with generic element $`g`$ and unit $`e`$. Consider $`Fun(G)`$, the set of complex functions on $`G`$. An element $`f`$ of $`Fun(G)`$ is specified by its values $`f_gf(g)`$ on the group elements $`g`$, and can be written as
$$f=\underset{gG}{}f_gx^g,f_g๐ช$$
(4.1)
where the functions $`x^g`$ are defined by
$$x^g(g^{})=\delta _g^{}^g$$
(4.2)
Thus $`Fun(G)`$ is a n-dimensional vector space, and the $`n`$ functions $`x^g`$ provide a basis. $`Fun(G)`$ is also a commutative algebra, with the usual pointwise sum and product, and unit $`I`$ defined by $`I(g)=1,gG`$. In particular:
$$x^gx^g^{}=\delta _{g,g^{}}x^g,\underset{gG}{}x^g=I$$
(4.3)
Consider now the left multiplication by $`g_1`$:
$$L_{g_1}g_2=g_1g_2,g_1,g_2G$$
(4.4)
This induces the left action (pullback) $`_{g_1}`$ on $`Fun(G)`$:
$$_{g_1}f(g_2)f(g_1g_2)|_{g_2},_{g_1}:Fun(G)Fun(G)$$
(4.5)
where $`f(g_1g_2)|_{g_2}`$ means $`f(g_1g_2)`$ seen as a function of $`g_2`$. Similarly we can define the right action on $`Fun(G)`$ as:
$$(_{g_1}f)(g_2)=f(g_2g_1)|_{g_2}$$
(4.6)
For the basis functions we find easily:
$$_{g_1}x^g=x^{g_1^1g},_{g_1}x^g=x^{gg_1^1}$$
(4.7)
Moreover:
$`_{g_1}_{g_2}=_{g_1g_2},_{g_1}_{g_2}=_{g_2g_1},`$ (4.8)
$`_{g_1}_{g_2}=_{g_2}_{g_1}`$ (4.9)
Differential calculus
A first-order differential calculus on $`A`$ is defined by
i) a linear map $`d`$: $`A\mathrm{\Gamma }`$, satisfying the Leibniz rule
$$d(ab)=(da)b+a(db),a,bA;$$
(4.10)
The โspace of 1-formsโ $`\mathrm{\Gamma }`$ is an appropriate bimodule on $`A`$, which essentially means that its elements can be multiplied on the left and on the right by elements of $`A`$ \[more precisely $`A`$ is a left module if $`a,bA,\rho ,\rho ^{}\mathrm{\Gamma }`$ we have: $`a(\rho +\rho ^{})=a\rho +a\rho ^{},(a+b)\rho =a\rho +b\rho ,a(b\rho )=(ab)\rho ,I\rho =\rho `$. Similarly one defines a right module. A left and right module is a bimodule if $`a(\rho b)=(a\rho )b`$\]. From the Leibniz rule $`da=d(Ia)=(dI)a+Ida`$ we deduce $`dI=0`$.
ii) the possibility of expressing any $`\rho \mathrm{\Gamma }`$ as
$$\rho =\underset{k}{}a_kdb_k$$
(4.11)
for some $`a_k,b_k`$ belonging to $`A`$.
To build a first order differential calculus on $`Fun(G)`$ we need to extend the algebra $`A=Fun(G)`$ to a differential algebra of elements $`x^g,dx^g`$ (it is sufficient to consider the basis elements and their differentials). Note however that the $`dx^g`$ are not linearly independent. In fact from $`0=dI=d(_{gG}x^g)=_{gG}dx^g`$ we see that only $`n1`$ differentials are independent. Every element $`\rho =adb`$ of $`\mathrm{\Gamma }`$ can be expressed as a linear combination (with complex coefficients) of terms of the type $`x^gdx^g^{}`$. Moreover $`\rho b\mathrm{\Gamma }`$ (i.e. $`\mathrm{\Gamma }`$ is also a right module) since the Leibniz rule and the multiplication rule (4.3) yield the commutations:
$$dx^gx^g^{}=x^gdx^g^{}+\delta _g^{}^gdx^g$$
(4.12)
allowing to reorder functions to the left of differentials.
Partial derivatives
Consider the differential of a function $`fFun(g)`$:
$$df=\underset{gG}{}f_gdx^g=\underset{ge}{}f_gdx^g+f_edx^e=\underset{ge}{}(f_gf_e)dx^g\underset{ge}{}_gfdx^g$$
(4.13)
We have used $`dx^e=_{ge}dx^g`$ (from $`_{gG}dx^g=0`$). The partial derivatives of $`f`$ have been defined in analogy with the usual differential calculus, and are given by
$$_gf=f_gf_e=f(g)f(e)$$
(4.14)
Not unexpectedly, they take here the form of finite differences (discrete partial derivatives at the origin $`e`$).
Left and right covariance
A differential calculus is left or right covariant if the left or right action of $`G`$ ($`_g`$ or $`_g`$) commutes with the exterior derivative $`d`$. Requiring left and right covariance in fact defines the action of $`_g`$ and $`_g`$ on differentials: $`_gdbd(_gb),bFun(G)`$ and similarly for $`_gdb`$. More generally, on elements of $`\mathrm{\Gamma }`$ (one-forms) we define $`_g`$ as:
$$_g(adb)(_ga)_gdb=(_ga)d(_gb)$$
(4.15)
and similar for $`_g`$. Computing for example the left and right action on the differentials $`dx^g`$ yields:
$`_g(dx^{g_1})d(_gx^{g_1})=dx^{g^1g_1},`$ (4.16)
$`_g(dx^{g_1})d(_gx^{g_1})=dx^{g_1g^1}`$ (4.17)
A differential calculus is called bicovariant if it is both left and right covariant.
Left invariant one forms
As in usual Lie group manifolds, we can introduce a basis in $`\mathrm{\Gamma }`$ of left-invariant one-forms $`\theta ^g`$:
$$\theta ^g\underset{hG}{}x^{hg}dx^h(=\underset{hG}{}x^hdx^{hg^1}),$$
(4.18)
It is immediate to check that indeed $`_k\theta ^g=\theta ^g`$. The relations (4.18) can be inverted:
$$dx^h=\underset{gG}{}(x^{hg}x^h)\theta ^g$$
(4.19)
From $`0=dI=d_{gG}x^g=_{gG}dx^g=0`$ one finds:
$$\underset{gG}{}\theta ^g=\underset{g,hG}{}x^hdx^{hg^1}=\underset{hG}{}x^h\underset{gG}{}dx^{hg^1}=0$$
(4.20)
Therefore we can take as basis of the cotangent space $`\mathrm{\Gamma }`$ the $`n1`$ linearly independent left-invariant one-forms $`\theta ^g`$ with $`ge`$ (but smaller sets of $`\theta ^g`$ can be consistently chosen as basis, see later).
The $`x,\theta `$ commutations (bimodule relations) are easily derived:
$$x^hdx^g=x^h\theta ^{g^1h}=\theta ^{g^1h}x^g(hg)x^h\theta ^g=\theta ^gx^{hg^1}(ge)$$
(4.21)
implying the general commutation rule between functions and left-invariant one-forms:
$$f\theta ^g=\theta ^g_gf$$
(4.22)
Thus functions do commute between themselves (i.e. $`Fun(G)`$ is a commutative algebra) but do not commute with the basis of one-forms $`\theta ^g`$. In this sense the differential geometry of $`Fun(G)`$ is noncommutative.
The right action of $`G`$ on the elements $`\theta ^g`$ is given by:
$$_h\theta ^g=\theta ^{ad(h)g},hG$$
(4.23)
where $`ad`$ is the adjoint action of $`G`$ on itself, i.e. $`ad(h)ghgh^1`$. Then bicovariant calculi are in 1-1 correspondence with unions of conjugacy classes (different from $`\{e\}`$) : if $`\theta ^g`$ is set to zero, one must set to zero all the $`\theta ^{ad(h)g},hG`$ corresponding to the whole conjugation class of $`g`$.
We denote by $`G^{}`$ the subset corresponding to the union of conjugacy classes that characterizes the bicovariant calculus on $`G`$ ($`G^{}=\{gG|\theta ^g0\}`$). Unless otherwise indicated, repeated indices are summed on $`G^{}`$ in the following.
A bi-invariant (i.e. left and right invariant) one-form $`\mathrm{\Theta }`$ is obtained by summing on all $`\theta ^g`$ with $`ge`$:
$$\mathrm{\Theta }=\underset{ge}{}\theta ^g$$
(4.24)
Exterior product
For a bicovariant differential calculus on a Hopf algebra $`A`$ an exterior product, compatible with the left and right actions of $`G`$, can be defined by
$$\theta ^{g_1}\theta ^{g_2}=\theta ^{g_1}\theta ^{g_2}\theta ^{g_1^1g_2g_1}\theta ^{g_1}\theta ^{g_1}\theta ^{g_2}\mathrm{\Lambda }_{g_3g_4}^{g_1g_2}\theta ^{g_3}\theta ^{g_4}$$
(4.25)
where the tensor product between elements $`\rho ,\rho ^{}\mathrm{\Gamma }`$ is defined to have the properties $`\rho a\rho ^{}=\rho a\rho ^{}`$, $`a(\rho \rho ^{})=(a\rho )\rho ^{}`$ and $`(\rho \rho ^{})a=\rho (\rho ^{}a)`$.
Note that:
$$\theta ^g\theta ^g=0\text{(no sum on }g\text{)}$$
(4.26)
Left and right actions on $`\mathrm{\Gamma }\mathrm{\Gamma }`$ are simply defined by:
$`_h(\rho \rho ^{})=_h\rho _h\rho ^{},`$ (4.27)
$`_h(\rho \rho ^{})=_h\rho _h\rho ^{}`$ (4.28)
(with the obvious generalization to $`\mathrm{\Gamma }\mathrm{}\mathrm{\Gamma }`$) so that for example:
$`_h(\theta ^i\theta ^j)=\theta ^i\theta ^j,`$ (4.29)
$`_h(\theta ^i\theta ^j)=\theta ^{ad(h)i}\theta ^{ad(h)j}`$ (4.30)
Compatibility of the exterior product with $``$ and $``$ means that
$$(\theta ^i\theta ^j)=\theta ^i\theta ^j,(\theta ^i\theta ^j)=\theta ^i\theta ^j$$
(4.31)
only the second relation being nontrivial, and verifiable using the definition (4.25).
We can generalize the definition (4.25) to exterior products of $`n`$ one-forms:
$$\theta ^{i_1}\mathrm{}\theta ^{i_n}W_{j_1..j_n}^{i_1..i_n}\theta ^{j_1}\mathrm{}\theta ^{j_n}$$
(4.32)
or in short-hand notation:
$$\theta ^1\mathrm{}\theta ^n=W_{1..n}\theta ^1\mathrm{}\theta ^n$$
(4.33)
where the labels 1..n in $`W`$ refer to index couples. The numerical coefficients $`W_{1\mathrm{}n}`$ are given through the recursion relation
$$W_{1\mathrm{}n}=_{1\mathrm{}n}W_{1\mathrm{}n1},$$
(4.34)
where
$$_{1\mathrm{}n}=1\mathrm{\Lambda }_{n1,n}+\mathrm{\Lambda }_{n2,n1}\mathrm{\Lambda }_{n1,n}\mathrm{}(1)^n\mathrm{\Lambda }_{12}\mathrm{\Lambda }_{23}\mathrm{}\mathrm{\Lambda }_{n1,n}$$
(4.35)
$`\mathrm{\Lambda }`$ being defined in (4.25) and $`W_1=1`$. The space of $`n`$-forms $`\mathrm{\Gamma }^n`$ is therefore defined as in the usual case but with the new permutation operator $`\mathrm{\Lambda }`$, and can be shown to be a bicovariant bimodule (see for ex. ), with left and right action defined as for $`\mathrm{\Gamma }\mathrm{}\mathrm{\Gamma }`$ with the tensor product replaced by the wedge product.
Exterior derivative
Having the exterior product we can define the exterior derivative
$$d:\mathrm{\Gamma }\mathrm{\Gamma }\mathrm{\Gamma }$$
(4.36)
$$d(a_kdb_k)=da_kdb_k,$$
(4.37)
which can easily be extended to $`\mathrm{\Gamma }^n`$ ($`d:\mathrm{\Gamma }^n\mathrm{\Gamma }^{(n+1)}`$), and has the following properties:
$$d(\rho \rho ^{})=d\rho \rho ^{}+(1)^k\rho d\rho ^{}$$
(4.38)
$$d(d\rho )=0$$
(4.39)
$$_g(d\rho )=d_g\rho $$
(4.40)
$$_g(d\rho )=d_g\rho $$
(4.41)
where $`\rho \mathrm{\Gamma }^k`$, $`\rho ^{}\mathrm{\Gamma }^n`$, $`\mathrm{\Gamma }^0Fun(G)`$ . The last two properties express the fact that $`d`$ commutes with the left and right action of $`G`$.
Tangent vectors
Using (4.19) to expand $`df`$ on the basis of the left-invariant one-forms $`\theta ^g`$ defines the (left-invariant) tangent vectors $`t_g`$:
$`df={\displaystyle \underset{gG}{}}f_gdx^g={\displaystyle \underset{hG^{}}{}}(_{h^1}ff)\theta ^h`$
$`{\displaystyle \underset{hG^{}}{}}(t_hf)\theta ^h`$ (4.42)
so that the โflatโ partial derivatives $`t_hf`$ are given by
$$t_hf=_{h^1}ff$$
(4.43)
The Leibniz rule for the flat partial derivatives $`t_g`$ reads:
$$t_g(ff^{})=(t_gf)_{g^1}f^{}+ft_gf^{}$$
(4.44)
In analogy with ordinary differential calculus, the operators $`t_g`$ appearing in (4.42) are called (left-invariant) tangent vectors, and in our case are given by
$$t_g=_{g^1}id$$
(4.45)
They satisfy the composition rule:
$$t_gt_g^{}=\underset{h}{}C_{g,g^{}}^ht_h$$
(4.46)
where the structure constants are:
$$C_{g,g^{}}^h=\delta _{g^{}g}^h\delta _g^h\delta _g^{}^h$$
(4.47)
and are $`ad(G)`$ invariant:
$$C_{ad(h)g_2,ad(h)g_3}^{ad(h)g_1}=C_{g_2,g_3}^{g_1}$$
(4.48)
Note 4.1 : The exterior derivative on any $`fFun(G)`$ can be expressed as a commutator of $`f`$ with the bi-invariant one-form $`\mathrm{\Theta }`$:
$$df=[\mathrm{\Theta },f]$$
(4.49)
as one proves by using (4.22) and (4.42).
Note 4.2 : From the fusion rules (4.46) we deduce the โdeformed Lie algebraโ (cf. ref.s ):
$$t_{g_1}t_{g_2}\mathrm{\Lambda }_{g_1,g_2}^{g_3,g_4}t_{g_3}t_{g_4}=๐ช_{g_1,g_2}^ht_h$$
(4.50)
where the $`๐ช`$ structure constants are given by:
$$๐ช_{g_1,g_2}^gC_{g_1,g_2}^g\mathrm{\Lambda }_{g_1,g_2}^{g_3,g_4}C_{g_3,g_4}^g=C_{g_1,g_2}^gC_{g_2,g_2g_1g_2^1}^g=\delta _{g_1}^{ad(g_2^1)g}\delta _{g_1}^g$$
(4.51)
and besides property (4.48) they also satisfy:
$$๐ช_{g_1,g_2}^g=๐ช_{g,g_2^1}^{g_1}$$
(4.52)
Moreover the following identities hold:
i) deformed Jacobi identities:
$$๐ช_{h_1,g_1}^k๐ช_{k,g_2}^{h_2}\mathrm{\Lambda }_{g_1,g_2}^{g_3,g_4}๐ช_{h_1,g_3}^k๐ช_{k,g_4}^{h_2}=๐ช_{g_1,g_2}^k๐ช_{h_1,k}^{h_2}$$
(4.53)
ii) fusion identities:
$$๐ช_{h_1,g}^k๐ช_{k,g^{}}^{h_2}=C_{g,g^{}}^h๐ช_{h_1,h}^{h_2}$$
(4.54)
Thus the $`๐ช`$ structure constants are a representation (the adjoint representation) of the tangent vectors $`t`$.
Cartan-Maurer equations, connection and curvature
From the definition (4.18) and eq. (4.22) we deduce the Cartan-Maurer equations:
$$d\theta ^g+\underset{g_1,g_2}{}C_{g_1,g_2}^g\theta ^{g_1}\theta ^{g_2}=0$$
(4.55)
where the structure constants $`C_{g_1,g_2}^g`$ are those given in (4.47).
Parallel transport of the vielbein $`\theta ^g`$ can be defined as in ordinary Lie group manifolds:
$$\theta ^g=\omega _g^{}^g\theta ^g^{}$$
(4.56)
where $`\omega _{g_2}^{g_1}`$ is the connection one-form:
$$\omega _{g_2}^{g_1}=\mathrm{\Gamma }_{g_3,g_2}^{g_1}\theta ^{g_3}$$
(4.57)
Thus parallel transport is a map from $`\mathrm{\Gamma }`$ to $`\mathrm{\Gamma }\mathrm{\Gamma }`$; by definition it must satisfy:
$$(a\rho )=(da)\rho +a\rho ,aA,\rho \mathrm{\Gamma }$$
(4.58)
and it is a simple matter to verify that this relation is satisfied with the usual parallel transport of Riemannian manifolds. As for the exterior differential, $``$ can be extended to a map $`:\mathrm{\Gamma }^n\mathrm{\Gamma }\mathrm{\Gamma }^{(n+1)}\mathrm{\Gamma }`$ by defining:
$$(\phi \rho )=d\phi \rho +(1)^n\phi \rho $$
(4.59)
Requiring parallel transport to commute with the left and right action of $`G`$ means:
$`_h(\theta ^g)=(_h\theta ^g)=\theta ^g`$ (4.60)
$`_h(\theta ^g)=(_h\theta ^g)=\theta ^{ad(h)g}`$ (4.61)
Recalling that $`_h(a\rho )=(_ha)(_h\rho )`$ and $`_h(\rho \rho ^{})=(_h\rho )(_h\rho ^{}),aA,\rho ,\rho ^{}\mathrm{\Gamma }`$ (and similar for $`_h`$), and substituting (4.56) yields respectively:
$$\mathrm{\Gamma }_{g_3,g_2}^{g_1}๐ช$$
(4.62)
and
$$\mathrm{\Gamma }_{ad(h)g_3,ad(h)g_2}^{ad(h)g_1}=\mathrm{\Gamma }_{g_3,g_2}^{g_1}$$
(4.63)
Therefore the same situation arises as in the case of Lie groups, for which parallel transport on the group manifold commutes with left and right action iff the connection components are $`ad(G)`$ \- conserved constant tensors. As for Lie groups, condition (4.63) is satisfied if one takes $`\mathrm{\Gamma }`$ proportional to the structure constants. In our case, we can take any combination of the $`C`$ or $`๐ช`$ structure constants, since both are $`ad(G)`$ conserved constant tensors. As we see below, the $`C`$ constants can be used to define a torsionless connection, while the $`๐ช`$ constants define a parallelizing connection.
As usual, the curvature arises from $`^2`$:
$$^2\theta ^g=R_g^{}^g\theta ^g^{}$$
(4.64)
$$R_{g_2}^{g_1}d\omega _{g_2}^{g_1}+\omega _{g_3}^{g_1}\omega _{g_2}^{g_3}$$
(4.65)
The torsion $`R^g`$ is defined by:
$$R^{g_1}d\theta ^{g_1}+\omega _{g_2}^{g_1}\theta ^{g_2}$$
(4.66)
Using the expression of $`\omega `$ in terms of $`\mathrm{\Gamma }`$ and the Cartan-Maurer equations yields
$`R_{g_2}^{g_1}=(\mathrm{\Gamma }_{h,g_2}^{g_1}C_{g_3,g_4}^h+\mathrm{\Gamma }_{g_3,h}^{g_1}\mathrm{\Gamma }_{g_4,g_2}^h)\theta ^{g_3}\theta ^{g_4}=`$
$`=(\mathrm{\Gamma }_{h,g_2}^{g_1}๐ช_{g_3,g_4}^h+\mathrm{\Gamma }_{g_3,h}^{g_1}\mathrm{\Gamma }_{g_4,g_2}^h\mathrm{\Gamma }_{g_4,h}^{g_1}\mathrm{\Gamma }_{g_4g_3g_4^1,g_2}^h)\theta ^{g_3}\theta ^{g_4}`$
$`R^{g_1}=(C_{g_2,g_3}^{g_1}+\mathrm{\Gamma }_{g_2,g_3}^{g_1})\theta ^{g_2}\theta ^{g_3}=`$
$`(๐ช_{g_2,g_3}^{g_1}+\mathrm{\Gamma }_{g_2,g_3}^{g_1}\mathrm{\Gamma }_{g_3,g_3g_2g_3^1}^{g_1})\theta ^{g_2}\theta ^{g_3}`$ (4.67)
Thus a connection satisfying:
$$\mathrm{\Gamma }_{g_2,g_3}^{g_1}\mathrm{\Gamma }_{g_3,g_3g_2g_3^1}^{g_1}=๐ช_{g_2,g_3}^{g_1}$$
(4.68)
corresponds to a vanishing torsion $`R^g=0`$ and could be referred to as a โRiemannianโ connection.
On the other hand, the choice:
$$\mathrm{\Gamma }_{g_2,g_3}^{g_1}=๐ช_{g_3,g_2^1}^{g_1}$$
(4.69)
corresponds to a vanishing curvature $`R_g^{}^g=0`$, as can be checked by using the fusion equations (4.54) and property (4.52). Then (4.69) can be called the parallelizing connection: finite groups are parallelizable.
Tensor transformations
Under the familiar transformation of the connection 1-form:
$$(\omega _j^i)^{}=a_k^i\omega _l^k(a^1)_j^l+a_k^id(a^1)_j^k$$
(4.70)
the curvature 2-form transforms homogeneously:
$$(R_j^i)^{}=a_k^iR_l^k(a^1)_j^l$$
(4.71)
The transformation rule (4.70) can be seen as induced by the change of basis $`\theta ^i=a_j^i\theta ^j`$, with $`a_j^i`$ invertible $`x`$-dependent matrix (use eq. (4.58) with $`a\rho =a_j^i\theta ^j`$).
Metric
The metric tensor $`\gamma `$ can be defined as an element of $`\mathrm{\Gamma }\mathrm{\Gamma }`$:
$$\gamma =\gamma _{i,j}\theta ^i\theta ^j$$
(4.72)
Requiring it to be invariant under left and right action of $`G`$ means:
$$_h(\gamma )=\gamma =_h(\gamma )$$
(4.73)
or equivalently, recalling $`_h(\theta ^i\theta ^j)=\theta ^i\theta ^j`$, $`_h(\theta ^i\theta ^j)=\theta ^{ad(h)i}\theta ^{ad(h)j}`$ :
$$\gamma _{i,j}๐ช,\gamma _{ad(h)i,ad(h)j}=\gamma _{i,j}$$
(4.74)
These properties are analogous to the ones satisfied by the Killing metric of Lie groups, which is indeed constant and invariant under the adjoint action of the Lie group.
On finite $`G`$ there are various choices of biinvariant metrics. One can simply take $`\gamma _{i,j}=\delta _{i,j}`$, or $`\gamma _{i,j}=๐ช_{l,i}^k๐ช_{k,j}^l`$.
For any biinvariant metric $`\gamma _{ij}`$ there are tensor transformations $`a_j^i`$ under which $`\gamma _{ij}`$ is invariant, i.e.:
$$a_h^{}^h\gamma _{h,k}a_k^{}^k=\gamma _{h^{},k^{}}a_h^{}^h\gamma _{h,k}=\gamma _{h^{},k^{}}(a^1)_k^k^{}$$
(4.75)
These transformations are simply given by the matrices that rotate the indices according to the adjoint action of $`G`$:
$$a_h^{}^h(g)=\delta _h^{}^{ad(\alpha (g))h}$$
(4.76)
where $`\alpha (g):GG`$ is an arbitrary mapping. Then these matrices are functions of $`G`$ via this mapping, and their action leaves $`\gamma `$ invariant because of the its biinvariance (4.74). Indeed substituting these matrices in (4.75) yields:
$$a_h^{}^h(g)\gamma _{h,k}a_k^{}^k(g)=\gamma _{ad([\alpha (g)]^1)h^{},ad([\alpha (g)]^1)k^{}}=\gamma _{h^{},k^{}}$$
(4.77)
proving the invariance of $`\gamma `$.
Consider now a contravariant vector $`\phi ^i`$ transforming as $`(\phi ^i)^{}=a_j^i(\phi ^j)`$. Then using (4.75) one can easily see that
$$(\phi ^k\gamma _{k,i})^{}=\phi ^k^{}\gamma _{k^{},i^{}}(a^1)_i^i^{}$$
(4.78)
i.e. the vector $`\phi _i\phi ^k\gamma _{k,i}`$ indeed transforms as a covariant vector.
Lie derivative and diffeomorphisms
The notion of diffeomorphisms, or general coordinate transformations, is fundamental in gravity theories. Is there such a notion in the setting of differential calculi on finite groups ? The answer is affirmative, and is based on general results obtained for Hopf algebras , of which finite groups are a simple example. As for differentiable manifolds, it relies on the existence of the Lie derivative.
Let us review the situation for Lie group manifolds. The Lie derivative $`l_{t_i}`$ along a left-invariant tangent vector $`t_i`$ is related to the infinitesimal right translations generated by $`t_i`$:
$$l_{t_i}\rho =\underset{\epsilon 0}{lim}\frac{1}{\epsilon }[_{\mathrm{exp}[\epsilon t_i]}\rho \rho ]$$
(4.79)
$`\rho `$ being an arbitrary tensor field. Introducing the coordinate dependence
$$l_{t_i}\rho (y)=\underset{\epsilon 0}{lim}\frac{1}{\epsilon }[\rho (y+\epsilon t_i)\rho (y)]$$
(4.80)
identifies the Lie derivative $`l_{t_i}`$ as a directional derivative along $`t_i`$. Note the difference in meaning of the symbol $`t_i`$ in the r.h.s. of these two equations: a group generator in the first, and the corresponding tangent vector in the second.
For finite groups the Lie derivative takes the form:
$$l_{t_g}\rho =[_{g^1}\rho \rho ]$$
(4.81)
so that the Lie derivative along $`t_g`$ coincides with the tangent vector $`t_g`$:
$$l_{t_g}=_{g^1}id=t_g$$
(4.82)
cf. the definition of $`t_g`$ in (4.45). For example
$$l_{t_g}(\theta ^{g_1}\theta ^{g_2})=\theta ^{ad(g^1)g_1}\theta ^{ad(g^1)g_2}\theta ^{g_1}\theta ^{g_2}$$
(4.83)
As in the case of differentiable manifolds, the Cartan formula for the Lie derivative acting on p-forms holds:
$$l_{t_g}=i_{t_g}d+di_{t_g}$$
(4.84)
see ref.s .
Exploiting this formula, diffeomorphisms (Lie derivatives) along generic tangent vectors $`V`$ can also be consistently defined via the operator:
$$l_V=i_Vd+di_V$$
(4.85)
This requires a suitable definition of the contraction operator $`i_V`$ along generic tangent vectors $`V`$, discussed in ref. .
We have then a way of defining โdiffeomorphismsโ along arbitrary (and x-dependent) tangent vectors for any tensor $`\rho `$:
$$\delta \rho =l_V\rho $$
(4.86)
and of testing the invariance of candidate lagrangians under the generalized Lie derivative.
Haar measure and integration
Since we want to define actions (integrals on $`p`$-forms), we must now define integration of $`p`$-forms on finite groups.
Let us start with integration of functions $`f`$. We define the integral map $`h`$ as a linear functional $`h:Fun(G)๐ช`$ satisfying the left and right invariance conditions:
$$h(_gf)=0=h(_gf)$$
(4.87)
Then this map is uniquely determined (up to a normalization constant), and is simply given by the โsum over $`G`$โ rule:
$$h(f)=\underset{gG}{}f(g)$$
(4.88)
Next we turn to define the integral of a p-form. Within the differential calculus we have a basis of left-invariant 1-forms, which may allow the definition of a biinvariant volume element. In general for a differential calculus with $`n`$ independent tangent vectors, there is an integer $`pn`$ such that the linear space of $`p`$-forms is 1-dimensional, and $`(p+1)`$\- forms vanish identically <sup>3</sup><sup>3</sup>3with the exception of $`Z_2`$, see Section 4.2. This means that every product of $`p`$ basis one-forms $`\theta ^{g_1}\theta ^{g_2}\mathrm{}\theta ^{g_p}`$ is proportional to one of these products, that can be chosen to define the volume form $`vol`$:
$$\theta ^{g_1}\theta ^{g_2}\mathrm{}\theta ^{g_p}=ฯต^{g_1,g_2,\mathrm{}g_p}vol$$
(4.89)
where $`ฯต^{g_1,g_2,\mathrm{}g_p}`$ is the proportionality constant. Note that the volume $`p`$-form is obviously left invariant. We can prove that it is also right invariant with the following argument. Suppose that $`vol`$ be given by $`\theta ^{h_1}\theta ^{h_2}\mathrm{}\theta ^{h_p}`$ where $`h_1,h_2,\mathrm{}h_p`$ are given group element labels. Then the right action on $`vol`$ yields:
$$_g[\theta ^{h_1}\mathrm{}\theta ^{h_p}]=\theta ^{ad(g)h_1}\mathrm{}\theta ^{ad(g)h_p}=ฯต^{ad(g)h_1,\mathrm{}ad(g)h_p}vol$$
(4.90)
Recall now that the โepsilon tensorโ $`ฯต`$ is necessarily made out of products of the $`\mathrm{\Lambda }`$ tensor of eq. (4.25), defining the wedge product. This tensor is invariant under the adjoint action $`ad(g)`$, and so is the $`ฯต`$ tensor. Therefore $`ฯต^{ad(g)h_1,\mathrm{}ad(g)h_p}=ฯต^{h_1,\mathrm{}h_p}=1`$ and $`_gvol=vol`$. This can be verified in the $`S_3`$ example of , and in the $`Z_2`$ case of next Section.
Having identified the volume $`p`$-form it is natural to set
$$fvolh(f)=\underset{gG}{}f(g)$$
(4.91)
and define the integral on a $`p`$-form $`\rho `$ as:
$`{\displaystyle \rho }={\displaystyle \rho _{g_1,\mathrm{}g_p}\theta ^{g_1}}\mathrm{}\theta ^{g_p}=`$
$`{\displaystyle \rho _{g_1,\mathrm{}g_p}ฯต^{g_1,\mathrm{}g_p}vol}`$
$`{\displaystyle \underset{gG}{}}\rho _{g_1,\mathrm{}g_p}(g)ฯต^{g_1,\mathrm{}g_p}`$ (4.92)
Due to the biinvariance of the volume form, the integral map $`:\mathrm{\Gamma }^p๐ช`$ satisfies the biinvariance conditions:
$$_gf=f=_gf$$
(4.93)
Moreover, under the assumption that $`d(\theta ^{g_2}\mathrm{}\theta ^{g_p})=0`$, i.e. that any exterior product of $`p1`$ left-invariant one-forms $`\theta `$ is closed, the important property holds:
$$๐f=0$$
(4.94)
with $`f`$ any $`(p1)`$-form: $`f=f_{g_2,\mathrm{}g_p}\theta ^{g_2}\mathrm{}\theta ^{g_p}`$. This property, which allows integration by parts, has a simple proof. Rewrite $`๐f`$ as:
$`{\displaystyle ๐f}={\displaystyle (df_{g_2,\mathrm{}g_p})\theta ^{g_2}}\mathrm{}\theta ^{g_p}+`$
$`+{\displaystyle f_{g_2,\mathrm{}g_p}d(\theta ^{g_2}\mathrm{}\theta ^{g_p})}`$ (4.95)
The second term in the r.h.s. vanishes by assumption. Using now (4.42) and (4.91):
$`{\displaystyle ๐f}={\displaystyle (t_{g_1}f_{g_2,\mathrm{}g_p})\theta ^{g_1}}\theta ^{g_2}\mathrm{}\theta ^{g_p}=`$
$`{\displaystyle [_{g_1^1}f_{g_2,\mathrm{}g_p}f_{g_2,\mathrm{}g_p}]ฯต^{g_1,\mathrm{}g_p}vol}=`$
$`ฯต^{g_1,\mathrm{}g_p}{\displaystyle \underset{gG}{}}[_{g_1^1}f_{g_2,\mathrm{}g_p}(g)f_{g_2,\mathrm{}g_p}(g)]=`$
$`=0`$ (4.96)
Q.E.D.
### 4.2 Bicovariant calculus on $`Z_2`$
In this Section we illustrate the general theory on $`Z_2`$, the simplest possible example.
Elements: $`e,u`$
with $`u^2=e`$.
Conjugation classes: $`\{e\},\{u\}`$.
There is therefore only one bicovariant calculus, corresponding to the only nontrivial congugation class $`\{u\}`$, of dimension 1.
Basis of functions on $`Z_2`$: $`\{x^e,x^u\}`$. Any function $`f`$ can be expanded as: $`f=f_ex^e+f_ux^u`$.
The left action of $`Z_2`$ on the functions is simply: $`_e=id,_u(x^e)=x^u,_u(x^u)=x^e`$. The right action coincides with the left action since $`Z_2`$ is abelian.
Partial derivatives:
$$df=\underset{gG}{}f_gdx^g=\underset{ge}{}f_gdx^g+f_edx^e=(f_uf_e)dx^u_uf=f_uf_e$$
(4.97)
Left-invariant one-forms:
$`\theta ^u=x^udx^e+x^edx^u=(x^ex^u)dx^u,`$ (4.98)
$`(\theta ^e=x^edx^e+x^udx^u=(x^ux^e)dx^u=\theta ^u)`$ (4.99)
Inversion formula:
$$dx^u=(x^ex^u)\theta ^u,(dx^e=dx^u)$$
(4.100)
Commutations $`x,\theta `$:
$$f\theta ^u=\theta ^u_ufx^e\theta ^u=\theta ^ux^u,x^u\theta ^u=\theta ^ux^e$$
(4.101)
Commutations $`x,dx`$:
$$x^edx^u=dx^ux^u,x^udx^u=dx^ux^efdx^u=dx^u_uf$$
(4.102)
Left and right action on $`\theta ^u`$:
$$_u\theta ^u=_u\theta ^u=(x^ux^e)dx^e=\theta ^u$$
(4.103)
Exterior product:
$$\theta ^u\theta ^u=0$$
(4.104)
when using the general formula (4.25). However the case of $`Z_2`$ is special in this respect: $`\theta ^u\theta ^u`$ can be set different from zero consistently with the differential calculus. Indeed for $`Z_2`$ (and only for this case) the whole calculus is consistent with the exterior product:
$$\theta ^u\theta ^u=\theta ^u\theta ^u$$
(4.105)
For example taking the exterior derivative of both members of the commutation relations $`x^udx^u=dx^ux^e`$ yields an identity (using $`d^2=0`$ and $`dx^e=dx^u`$), and does not imply $`dx^udx^u=0`$. Using then the expression of $`dx^u`$ in terms of $`\theta ^u`$ and the $`x,\theta `$ commutations, one finds
$$dx^udx^u=\theta ^u\theta ^u$$
(4.106)
so that $`\theta ^u\theta ^u`$ can be different from zero. In the case of $`Z_N,N>2`$ the situation is different since taking the exterior derivative of the $`x,dx`$ commutations implies the vanishing of the exterior product of a left-invariant one-form with itself: then one has to adopt the canonical definition as given in (4.25).
For $`Z_2`$, we denote the two possibilities calculus I ($`dx^udx^u=0`$) and calculus II ($`dx^udx^u0`$).
Tangent vector
$$t_uf=(_uid)f,t_ut_u=(_uid)(_uid)=_e2_u+id=2(id_u)=2t_u$$
(4.107)
Cartan-Maurer equations
Calculus I:
$$d\theta ^u=0$$
(4.108)
Calculus II:
$$d\theta ^u=dx^udx^e+dx^edx^u=2dx^udx^u=2\theta ^u\theta ^u$$
(4.109)
Connection
$$\omega _u^u=\mathrm{\Gamma }_{u,u}^u\theta ^u$$
(4.110)
where $`\mathrm{\Gamma }_{u,u}^u=constant=c`$ satisfies left and right invariance.
Curvature and torsion
Calculus I:
$`R_u^u=d\omega _u^u+\omega _u^u\omega _u^u=cd\theta ^u+c^2\theta ^u\theta ^u=0`$ (4.111)
$`T^u=d\theta ^u+c\theta ^u\theta ^u=0`$ (4.112)
Calculus II:
$`R_u^u=cd\theta ^u+c^2\theta ^u\theta ^u=(2c+c^2)\theta ^u\theta ^u`$ (4.113)
$`T^u=d\theta ^u+c\theta ^u\theta ^u=(2+c)\theta ^u\theta ^u`$ (4.114)
In this case $`c=2`$ gives a flat and torsionless connection.
Integration
For calculus I, the volume form is the one-form $`\theta ^u`$, and the integral of a one-form $`\rho =\rho _u\theta ^u`$ is simply:
$$\rho =\rho _u\theta ^u=\rho _uvol=\underset{gG}{}\rho _u(g)=\rho _u(e)+\rho _u(u)$$
(4.115)
Integration by parts holds since:
$$๐f=(t_uf)\theta ^u=[(_uid)f]vol=\underset{gG}{}(_uff)(g)=0$$
(4.116)
for $`f`$ = 0-form.
In the special case of $`Z_2`$, choosing calculus II, there is no upper limit to the degree of a $`p`$-form, since all the products $`\theta ^u\theta ^u\mathrm{}\theta ^u`$ are nonvanishing. Then any one of these products, being bi-invariant, can be chosen as volume form ! Supposing to take the $`p`$-form volume as volume form, the integral of a p-form $`\rho _{u,u,\mathrm{}u}\theta ^u\theta ^u\mathrm{}\theta ^u`$ is then simply $`\rho _{u,u,\mathrm{}u}(e)+\rho _{u,u,\mathrm{}u}(u)`$. Choosing $`\theta ^u\theta ^u`$ as volume form, we find $`๐\sigma 0`$ (where $`\sigma `$ is a $`1`$ form); indeed:
$$๐\sigma =d(\sigma _u\theta ^u)=(t_u\sigma _u)\theta ^u\theta ^u+2\sigma _u\theta ^u\theta ^u=2[\sigma _u(e)+\sigma _u(u)]$$
(4.117)
Note 4.3: Choosing higher volume forms one retrieves the integration by parts rule, essentially because an exterior product of two or more $`\theta ^u`$โs is closed.
### 4.3 Kaluza-Klein gauge theory on $`M_4\times Z_2`$
In this example we label the $`M_4`$ coordinates as $`x^\mu `$ and the $`Z_2`$ coordinate as $`y`$. Field theories (and in particular gauge theories) on discrete spaces have been considered by many authors. The treatment of this Section is closer in spirit to the works of
Calculus on $`M_4\times Z_2`$
The $`y`$ coordinate can take the values $`e,u`$, and any function $`f`$ on $`M_4\times Z_2`$ is expanded as:
$$f(x,y)=f_e(x)y^e+f_u(x)y^u$$
(4.118)
where $`y^e,y^u`$ are defined as usual to be โdualโ to the $`Z_2`$ points: $`y^e(e)=y^u(u)=1,y^e(u)=y^u(e)=0`$. We will frequently use the notation:
$$\stackrel{~}{f}_uf=f_u(x)y^e+f_e(x)y^u$$
(4.119)
Thus $`\stackrel{~}{f}`$ is obtained from $`f`$ simply by exchanging its components along $`y^e,y^u`$.
The only independent $`Z_2`$ differential $`dy^u`$ will be simply denoted by $`dy`$.
Note that
$$\stackrel{~}{f}dy=dyf$$
(4.120)
cf. eq. (4.102).
To define completely the differential geometry on $`M_4\times Z_2`$ we need the rules:
$$dx^\mu dy=dydx^\mu ,fdx^\mu =dx^\mu f$$
(4.121)
A basis of differentials is given by $`dx^M=(dx^\mu ,dy)`$, so that any one-form $`A(x,y)`$ is expanded as:
$$A(x,y)=A(x,y)_Mdx^M=A_\mu (x,y)dx^\mu +A_{}(x,y)dy$$
(4.122)
Finally, integration of a function $`f(x,y)`$ on $`M_4\times Z_2`$ is defined by:
$$_{M_4\times Z_2}fvol_{M_4}\underset{Z_2}{}f(x,y)d^4x=_{M_4}[f_e(x)+f_u(x)]d^4x$$
(4.123)
Gauge potential
Consider now the one-form $`A`$ to be the potential 1-form of a gauge theory: then it must be also matrix valued. For example in ordinary Yang-Mills theory, $`A(x)=A_\mu ^IT_Idx^\mu `$ where $`T_I`$ are the generators of the gauge group G in some irreducible representation.
As in the usual case, we define G-gauge transformations on the potential $`A(x,y)`$ as:
$$A^{}=(dG)G^1+GAG^1$$
(4.124)
where $`G=G(x,y)`$ is a group element is some irrep, depending on the point $`(x,y)M_4\times Z_2`$. In components:
$$A_\mu ^{}=(_\mu G)G^1+GA_\mu G^1,A_{}^{}=(_{}G)\stackrel{~}{G}^1+GA_{}\stackrel{~}{G}^1$$
(4.125)
the derivative along $`y`$ being denoted by $`_{}`$. Note that
$$_{}f(x,y)=f_u(x)f_e(x)=\frac{1}{2}(y^ey^u)(\stackrel{~}{f}f)J(\stackrel{~}{f}f)$$
(4.126)
where we have introduced the function $`J\frac{1}{2}(y^ey^u)`$.
The transformation laws tell us something about the matrix structure of the gauge potential $`A`$. The potential components $`A_\mu `$ must belong to the Lie algebra of G, since $`(_\mu G)G^1Lie(G)`$. On the other hand $`A_{}`$ does not belong to Lie(G) but rather to the group algebra of G. Indeed $`_{}G`$ is a finite difference of group elements, and thus $`(_{}G)G^1`$ belongs to the group algebra; then the second eq. in (4.125) implies that $`A_{}`$ is matrix valued in the group algebra of G.
For definiteness, we consider unitary groups, so that $`G^{}=G^1`$. Then $`A_\mu `$ is antihermitian (since the generators $`T_I`$ are antihermitian), while $`A_{}`$, being in the $`U(N)`$ group algebra is a sum of $`U(N)`$ matrices.
We can consistently incorporate hermitian conjugation in the $`M_4\times Z_2`$ \- differential calculus by setting:
$`(dx^\mu )^{}=dx^\mu ,(dy)^{}=dy`$ (4.127)
$`(fdy)^{}=dyf^{}`$ (4.128)
Matter fields
Matter fields $`\psi `$ are taken to transform in an irrep of G:
$$\psi ^{}=G\psi ,(\psi ^{})^{}=\psi ^{}G^{}=\psi ^{}G^1$$
(4.129)
and their covariant derivative, defined by
$$D\psi =d\psi +A\psi ,D\psi ^{}=d\psi ^{}\psi ^{}A$$
(4.130)
transforms as it should: $`(D\psi )^{}=G(D\psi ),(D\psi ^{})^{}=(D\psi ^{})G^1`$. Requiring compatibility of hermitian conjugation with the covariant derivative , i.e. $`(D\psi )^{}=D\psi ^{}`$, implies:
$$A^{}=A$$
(4.131)
that is, $`A`$ must be an antihermitian connection. This is compatible with its transformation rule (4.124). In components the antihermitian condition reads:
$$A_\mu ^{}=A_\mu ,A_{}^{}=\stackrel{~}{A}_{}$$
(4.132)
Field strength
The field strength $`F`$ is formally defined as usual:
$$F=dA+AA$$
(4.133)
so that it transforms as:
$$F^{}=GFG^1$$
(4.134)
The components of the 2-form $`F`$ are labelled as follows:
$$FF_{MN}dx^Mdx^NF_{\mu \nu }dx^\mu dx^\nu +2F_\mu dx^\mu dy+F_{}dydy$$
(4.135)
Therefore the $`F`$ components are given by:
$`F_{\mu \nu }={\displaystyle \frac{1}{2}}(_\mu A_\nu _\nu A_\mu +A_\mu A_\nu A_\nu A_\mu )`$ (4.136)
$`F_\mu ={\displaystyle \frac{1}{2}}(_\mu A_{}_{}A_\mu +A_\mu A_{}A_{}\stackrel{~}{A}_\mu )`$ (4.137)
$`F_{}=_{}A_{}+A_{}\stackrel{~}{A}_{}`$ (4.138)
and transform as:
$`F_{\mu \nu }^{}(x,y)=G(x,y)F_{\mu \nu }(x,y)G^1(x,y)`$ (4.139)
$`F_\mu ^{}(x,y)=G(x,y)F_\mu (x,y)\stackrel{~}{G}^1(x,y)`$ (4.140)
$`F_{}^{}(x,y)=G(x,y)F_{}(x,y)G^1(x,y)`$ (4.141)
The gauge action
Formally the gauge action has the same expression as in the usual case:
$$A_{YM}=_{M_4\times Z_2}Tr_G[F_{AB}F_{AB}^{}]vol$$
(4.142)
When expanded into components:
$$A_{YM}=_{M_4\times Z_2}\underset{Z_2}{}Tr_G[F_{\mu \nu }F_{\mu \nu }^{}+2F_\mu F_\mu ^{}+F_{}F_{}^{}]d^4x$$
(4.143)
This action is invariant under the G gauge transformations (4.141). We now rewrite it in a suggestive way, by introducing the โlinkโ field $`U(x,y)`$:
$$U(x,y)1\text{I}+J^1A_{}$$
(4.144)
Then
$$F_\mu =\frac{1}{2}J(_\mu U+A_\mu UU\stackrel{~}{A}_\mu )\frac{1}{2}JD_\mu U,F_{}=\frac{1}{4}(1\text{I}U\stackrel{~}{U})$$
(4.145)
Using the transformation rules (4.125) one finds that the link field $`U`$ and its covariant derivative vary homogeneously:
$$U^{}=GU\stackrel{~}{G}^1,(D_\mu U)^{}=G(D_\mu U)\stackrel{~}{G}^1$$
(4.146)
Moreover the antihermiticity of $`A`$ (4.132) implies:
$$U^{}=\stackrel{~}{U}$$
(4.147)
(use $`\stackrel{~}{J}=J`$). Expanding $`U(x,y)`$ as $`U_e(x)y^e+U_u(x)y^u`$, relation (4.147) becomes $`U_e^{}=U_u`$.
Using the expressions (4.145) for the field strength components finally yields the action in the form:
$$A_{YM}=d^4xTr_G\underset{Z_2}{}[F_{\mu \nu }F_{\mu \nu }+\frac{1}{16}D_\mu U(D_\mu U)^{}+\frac{1}{16}(1\text{I}UU^{})^2]$$
(4.148)
The sum on $`Z_2`$ is easy to perform, and taking into account $`U_e^{}=U_u`$ we find:
$$A_{YM}=2d^4xTr_G[F_{\mu \nu }F_{\mu \nu }+\frac{1}{16}D_\mu U(D_\mu U)^{}+\frac{1}{16}(1\text{I}UU^{})^2]$$
(4.149)
where now $`U(x)U_e(x)`$ can be seen as a complex Higgs field, with a symmetry-breaking potential. The cyclic property of $`Tr_G`$ has been used to achieve this final form of $`A_{YM}`$. Moreover we have identified for simplicity $`A_\mu A_\mu (u)=A_\mu (e)`$ so that the sum on $`Z_2`$ of the usual Yang-Mills term just gives a factor of 2.
Coupling to fermion matter
We can add a Dirac term $`_{Dirac}`$ to the integrand of $`A_{YM}`$:
$$_{Dirac}=Re[i\psi ^{}\gamma _0\gamma _MD_M\psi ]$$
(4.150)
where now the matter field $`\psi (x,y)`$ is a $`d=5`$ Dirac spinor and has therefore 4 complex spinor components. Splitting the sum on the index M:
$$_{Dirac}=Re[i\psi ^{}\gamma _0\gamma _\mu D_\mu \psi ]+Re[i\psi ^{}\gamma _0\gamma _5_{}\psi +i\psi ^{}\gamma _0\gamma _5A_{}\stackrel{~}{\psi }]$$
(4.151)
The first term is just the usual kinetic term in $`d=4`$; the last two terms give:
$$Re[iJ\overline{\psi }\gamma _5\psi +iJ\overline{\psi }\gamma _5U\stackrel{~}{\psi }]$$
(4.152)
The first term in square parentheses disappears (since its real part vanishes) and the second is:
$$Re[i\overline{\psi }_e\gamma _5U_e\psi _uy^ei\overline{\psi }_u\gamma _5U_u\psi _ey^u]$$
(4.153)
Summing on $`Z_2`$ and redefining $`\psi \psi _e,\chi i\gamma _5\psi _u,UU_e`$ one finds finally:
$$_{Dirac}=i(\overline{\psi }\gamma _\mu D_\mu \psi +\overline{\chi }\gamma _\mu D_\mu \chi )+\overline{\psi }U\chi +\overline{\chi }U^{}\psi $$
(4.154)
that is a kinetic term for the Dirac fields, and Yukawa couplings Higgs-Fermi-Fermi. We emphasize the appearance of the correct Higgs couplings to the Fermi fields as an output of the Kaluza-Klein mechanism on $`M_4\times Z_2`$ rather than an ad hoc addition to the Lagrangian. Also, the Higgs sector appears in (4.149) with the right form of the potential. This provides a nice interpretation of the Higgs appearance in the $`d=4`$ theory in terms of a Kaluza-Klein gauge theory coupled to Dirac fermions on $`M_4\times Z_2`$. The Higgs field is the component of the potential 1-form along the discrete dimension.
Note that the Kaluza-Klein mechanism on discrete internal spaces yields a finite number of fields in $`d=4`$: there is no infinite tower of massive modes ! The โharmonicโ analysis (4.1) on finite groups is elementary.
Acknowledgements
I have benefited from numerous discussions with G. Arcioni, P. Aschieri, F. Lizzi and M. Tarlini.
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# INSTABILITIES AND CLUMPING IN TYPE Ia SUPERNOVA REMNANTS
## 1 INTRODUCTION
Type Ia supernovae (SNe) are thought to arise from the thermonuclear explosion of a white dwarf that accretes mass from its binary companion (Whelan & Iben 1973; Chevalier 1981a; Nomoto, Thielemann, & Yokoi 1984; Harkness 1991). Both Chandrasekhar and sub-Chandrasekhar mass models have been considered for the mass of the white dwarf at the time of the explosion. In the sub-Chandrasekhar mass models, He detonation occurs on the surface of the white dwarf (Woosley & Weaver 1994; Livne & Arnett 1995). Although these models have some attractive features, they seem less able than Chandrasekhar mass models to explain the light curves and spectra of Type Ia supernovae (e.g., Hรถflich & Khokhlov 1996). In Chandrasekhar mass models, ignition is initiated in the high density central regions of the white dwarf. In one type of model, the burning front propagates subsonically through the white dwarf. A popular example of this deflagration model is the W7 model of Nomoto et al. (1984), which can approximately reproduce the observed light curves. Branch et al. (1985) found that the model can reproduce the observed spectra if there is mixing of matter with a velocity $`>8,000\mathrm{km}\mathrm{s}^1`$, although Harkness (1991), in more detailed models, found that mixing may not be necessary. In another class of Chandrasekhar mass models, the initial subsonic burning front, or deflagration, makes a transition to a supersonic burning front, or detonation. This is the delayed detonation, or DD, model (Khokhlov 1991).
Although observations of extragalactic Type Ia supernovae have given insights into the nature of the explosions, a clear determination of the supernova model has not been possible. Observations of Galactic supernova remnants provide another window on the structure and composition of supernovae. Likely remnants of Type Ia explosions are Tychoโs supernova and SN 1006. Tychoโs supernova (SN 1572) appears to be a Type Ia supernova based on its reconstructed light curve (Baade 1945), although it may have been subluminous (van den Bergh 1993). Schaefer (1996) noted that the light curve allows a Type Ib origin, but the lack of evidence for stellar wind or photoionization effects in the surroundings of Tychoโs supernova argues against a massive star origin. At X-ray and radio wavelengths, Tycho shows a similar morphology, characterized by a limb-brightened circular shell. The outer edge is sharply defined and is interpreted as the forward shock front. Seward, Gorenstein, & Tucker (1983) found evidence for 400 X-ray clumps in the Einstein observatory image. At optical wavelengths, an H$`\alpha `$ emission filament, which is thought to be from a fast shock moving into a surrounding neutral material, is observed along the NE perimeter, supporting the shock interpretation. The remnant outline spans $`8^{}`$ in diameter, corresponding to a linear radius of $``$ 3 pc for an assumed distance of 2.5 kpc (Raymond 1984). The HI absorption show that the environment surrounding Tycho is fairly uniform, without a prominent density gradient (Reynoso et al. 1999). This property makes it especially suitable for detailed study.
Hamilton et al. (1986) calculated a one-dimensional hydrodynamical model to study the X-ray spectrum of Tychoโs remnant and obtained good agreement with spectra from the Einstein Observatory and other X-ray observatories. They took a three-layer composition structure and were successful in reproducing the observed Si line emission. However, they chose a constant density profile for the supernova. When this type of supernova density profile interacts with a constant density ambient medium, the outer shocked layers (where Si was placed in the model) end up dense and cool. A constant density profile is probably a poor representation of an exploded Type Ia supernova. Dwarkadas & Chevalier (1998, hereafter DC98) recently reviewed the types of explosion models that reproduce the basic properties of Type Ia supernovae and found that an exponential profile is the best approximation for the density profile overall, although there can be significant deviations from this profile. When such a profile interacts with a constant density medium, the shocked supernova gas has an approximately constant temperature. In order to reproduce the observed spectral properties of Tychoโs remnant, DC98 suggested that either the supernova gas has dense clumps or there was an early phase of interaction with a circumstellar medium. This was on the basis of one-dimensional hydrodynamic models for the interaction of a Type Ia supernova with a surrounding medium.
In order to make further progress, we extend the hydrodynamic modeling effort to two dimensions. One motivation for this is the investigation of instabilities that result from the deceleration of the supernova ejecta by the ambient medium. Chevalier, Blondin, & Emmering (1992) previously investigated the interaction of a power law supernova density profile with an ambient medium. In this case, the flow approaches a self-similar state, although turbulent motions are continuously fed by the instability. In the case of an exponential profile, the flow is no longer self-similar; the density profile effectively evolves from a steep power law profile to a flatter one. Chevalier et al. (1992) investigated the instability for a range of power law indices and found that the flow remained qualitatively similar. In view of this, we expect the exponential case to be qualitatively similar while the steep part of the density profile is interacting with the surroundings, even though the flow is no longer self-similar. Our aim is to determine the nature of the evolution. Dwarkadas (2000) has already made some numerical studies of this phenomenon.
Another part of our effort is to include density inhomogeneities in the freely expanding supernova gas. The ASCA observatory has made possible the imaging of Tychoโs remnant in moderately narrow X-ray bands (Hwang & Gotthelf 1997). There are two X-ray emitting knots on the eastern side of the remnant where there is a protrusion in the remnantโs outline (Vancura et al. 1995; Hughes 1997; Hwang & Gotthelf 1997). One of these knots is strong in Si line emission and the other is strong in Fe line emission. With the ROSAT HRI, Hughes (1997) has measured the proper motion of features in the X-ray emission. He found that these two knots show approximately undecelerated motion. On the other hand, radio observations over the whole remnant show substantial deceleration, with $`Vt/R0.47`$ (Reynoso et al. 1997 and references therein), where $`V`$ is the shock velocity and $`R`$ is the radius. The implication is that the knots represent not only an inhomogeneous composition, but also higher density clumps that have not been decelerated. At a distance of 2.5 kpc, the velocities of the knots are about $`8,300\mathrm{km}\mathrm{s}^1`$ (Hughes 1997).
The other well observed Type Ia supernova remnant, SN 1006, does not show evidence for the expansion of X-ray emitting clumps. However, absorption line studies of the Schweizer-Middleditch star behind the remnant have provided a powerful probe of the supernova remnant structure and composition (Wu et al. 1983, 1997; Fesen et al. 1988). These observations show evidence for central cool Fe II that is compatible with expectations for the unshocked Fe rich matter in the central region of a Type Ia supernova, although the amount of Fe is smaller than expected (Hamilton et al. 1997 and references therein). In addition, there is evidence for Si II absorption that is redshifted over the velocity range $`3,0007,000\mathrm{km}\mathrm{s}^1`$ (Wu et al. 1997; Hamilton et al 1997). Fesen et al. (1988) initially argued that the Si rich matter is in a clump on the far side of SN 1006. However, Hamilton et al. (1997) more recently argued for a spherically symmetric model for the supernova ejecta in which there is a transition from Fe to Si rich gas at a velocity of $`5,600\mathrm{km}\mathrm{s}^1`$ and the supernova remnant is considerably more extended on the back side because of a low interstellar density in that direction. They argue against a Si clump, but appear to assume that the clump gas would have been shocked and cooled. Another possibility is that low ionization absorption is due to unshocked matter in the clump. The clump may be moving at about $`7,000\mathrm{km}\mathrm{s}^1`$ and matter stripped from the clump may lead to the matter extending to lower velocities. The remnant can then have a shape that is closer to spherical symmetry; the presence of a clump implies that the ejecta are not in spherical shells, but this is also indicated in Tychoโs remnant.
Our aim here is to examine how inhomogeneities in the ejecta interact with the decelerated interaction region. The observed properties can place constraints on the nature of the inhomogeneities, which can be important for the explosion models.
The plan of our paper is as follows. In ยง 2, we discuss the density structure of a Type Ia supernova and consider mechanisms that might give rise to inhomogeneities in the ejecta. Our method is given in ยง 3. Results on the growth of instabilities from small perturbations are in ยง 4. The evolution of nonlinear inhomogeneities is given in ยง 5. Our results are discussed in the context of observations of Tychoโs supernova remnant in ยง 6. Implications for models for Type Ia supernovae are in ยง 7 and the conclusions are in ยง 8.
## 2 DENSITY STRUCTURE IN TYPE Ia SUPERNOVAE
Soon after a supernova explosion, the ejecta are expected to be freely expanding, with velocity $`v=r/t`$. In this phase, DC98 found that an exponential density profile generally describes the density distribution obtained from numerical 1-D (one-dimensional) explosion models, including delayed detonation, pulsating delayed detonation, He detonation, and deflagration models. The exponential density profile is given by
$$\rho _{SN}=A\mathrm{exp}(v/v_e)t^3,$$
(1)
where $`A`$ is a constant and $`v_e`$ is the velocity scale height, which is determined by the total ejecta mass $`M`$ and explosion energy $`E`$. This expression allows for spherical expansion; an element of supernova gas moves with a constant velocity and has its density drop as $`t^3`$. Integrating the density and energy density over space and time gives (DC98)
$$M=4\pi Av_e^3,E=48\pi Av_e^5,$$
(2)
or
$$v_e=(E/6M)^{1/2}=2.44\times 10^8E_{51}^{1/2}\left(\frac{M}{M_{ch}}\right)^{1/2}\mathrm{cm}\mathrm{s}^1,$$
(3)
$$A=\frac{6^{3/2}}{8\pi }\frac{M^{5/2}}{E^{3/2}}=7.67\times 10^6\left(\frac{M}{M_{ch}}\right)^{5/2}E_{51}^{3/2}\mathrm{g}\mathrm{s}^3\mathrm{cm}^3,$$
(4)
where $`M_{ch}1.4M_{}`$ is the Chandrasekhar mass and $`E_{51}`$ is the explosion energy in units of $`10^{51}`$ ergs. A larger ratio of $`E/M`$ gives a larger velocity scale height and produces a flatter density profile at a given velocity.
In some Type Ia models and core collapse models, the outer density distribution has been described by a power law. Then, the ejecta density is given by $`\rho v^nt^3`$, where $`v=r/t`$ and $`n`$ is a constant $`>5`$. If the ejecta interact with a shallow power law medium $`\rho r^s`$ with $`s<3`$, there exist self-similar solutions for the shocked flow (Chevalier 1982a). The problem contains only two independent dimensional parameters: the coefficients of the supernova density and the external density. In the exponential case there is an additional parameter, the velocity scale height $`v_e`$, so that the problem is not self-similar.
The exponential profile is an oversimplification of the density profile, as can be seen in fig. 1 of DC98. In a thermonuclear explosion, the elements of gas receive most of their energy from the burning of that gas, as opposed to a shock wave as in the case of core collapse supernovae. The expansion of these heated elements is what leads to the approximately exponential density distribution. If there is incomplete burning of the material, the internal energy per unit mass of the gas is lower that than in neighboring regions where complete burning has taken place. The gas is expected to evolve toward pressure equilibrium, compressing the incompletely burned region. The magnitude of the effect can be estimated by considering the relative energy densities produced by different burning processes, and allowing the gas with the greater burning to adiabatically expand to come into pressure equilibrium with the other gas. The difference in burning a C/O mixture to Si or to Ni does not lead to a large density difference: $`40`$% higher density in the Si gas. We will find that this is not sufficient to explain the clumping indicated in Tychoโs remnant. If some C or O is left unburned, a greater degree of compression is expected, but this typically applies to only the outermost layers of the supernova. This effect can apparently be seen in models (e.g., the W7 model density distribution shown in Fig. 1 of Branch et al. 1985).
On a longer timescale, the density structure can be changed by the Ni bubble effect. The occurrence of SN 1987A brought the realization that the energy release by radioactive <sup>56</sup>Ni and <sup>56</sup>Co after the initial explosion could affect the density and composition structure within the supernova (e.g., Woosley 1988). Li, McCray, & Sunyaev (1993) analyzed the Fe/Ni/Co lines from SN 1987A and showed that the filling factor of Fe with velocity $`2500\mathrm{km}\mathrm{s}^1`$ was $`0.3`$ although the mass fraction of this element was only $`0.01`$ in the same region. They attributed this effect to the expansion of gas heated by radioactive power deposition. Basko (1994) made more detailed calculations of the expansion effect. Li et al. (1993) noted that the structure should show up in young supernova remnants and briefly mentioned the Crab Nebula. Chevalier (2000) contrasted the operation of the effect in Type Ia and in core collapse supernovae, and we follow his discussion of the Type Ia events here.
We examine the expansion of a spherical volume of radioactive gas, following Basko (1994). To get a crude estimate of the effect, we assume that immediately after the explosion, the supernova gas is freely expanding and has constant density throughout. The velocity at the edge of a spherical clump of <sup>56</sup>Ni with mass $`M_{\mathrm{Ni}}`$ is
$$U_0=7.7\times 10^3\left(\frac{M_{\mathrm{Ni}}}{0.5M_{}}\right)^{1/3}\left(\frac{E}{10^{51}\mathrm{ergs}}\right)^{1/2}\left(\frac{M}{M_{ch}}\right)^{5/6}\mathrm{km}\mathrm{s}^1.$$
(5)
The deposition of radioactive energy in the <sup>56</sup>Ni gas leads to expansion so that the velocity of the outer edge of the region becomes $`U_{\mathrm{}}`$. Conservation of energy leads to the equation
$$X^5X^2=\frac{5Q}{U_0^2},$$
(6)
where $`X=U_{\mathrm{}}/U_0`$ and $`Q`$ is the energy per gram deposited by radioactivity. The assumption of full deposition of energy in the gas implies $`Q=3.69\times 10^{16}`$ ergs g<sup>-1</sup> for <sup>56</sup>Ni$`^{56}`$Co and $`Q=7.87\times 10^{16}`$ ergs g<sup>-1</sup> for <sup>56</sup>Co$`^{56}`$Fe. The expansion is limited by the diffusion of either the $`\gamma `$-rays or the photons out of the high pressure bubble so that the full values of $`Q`$ listed above may not be achieved. The bubble expansion is expected to drive a shell in the incompletely burned matter on the outside. Basko (1994) showed that the shell is subject to the Rayleigh-Taylor instability, but that the instability does not have much effect on the shell expansion. The shell may break into clumps. Equation (6) shows that if $`U_0`$ is large, the expansion of the Ni bubble is relatively small and the factor increase in the filling factor ($`X^3`$) is small. Clumps can still be created in the swept up gas. On the other hand, small regions of <sup>56</sup>Ni have the potential to substantially increase their filling factor in the supernova gas.
The reference values used in equation (5) are those thought to be typical of Type Ia supernovae. These supernovae typically reach maximum light on a timescale of 15 days. This is the timescale on which diffusion of the internal radiative energy becomes effective, which has the effect of smoothing pressure gradients and stopping Ni bubble expansion. The mean life of <sup>56</sup>Ni is 8.8 days and that of <sup>56</sup>Co is 114 days, so we estimate that the <sup>56</sup>Ni radioactive energy is deposited in the radioactive gas, but not that from <sup>56</sup>Co. After the expansion due to the radioactivity, the outer velocity of the Fe-rich region is $`8,400\mathrm{km}\mathrm{s}^1`$ and about $`0.13M_{}`$ of matter has been swept into a shell around the Fe-rich region. These numbers are approximate and will change for the density profile of a specific supernova.
As discussed in ยง 1, Tychoโs remnant shows evidence for two knots that have moved with undecelerated motion. At a distance of 2.5 kpc, the velocities of the knots are about $`8,300\mathrm{km}\mathrm{s}^1`$ (Hughes 1997). This velocity is consistent with that expected in a clumpy shell formed by the expansion of radioactive gas. However, Fe is expected in a clump only if it formed as non-radioactive Fe. The Fe clump may have formed as <sup>54</sup>Fe. X-ray spectra of Tychoโs remnant have been obtained with the result that, in general, Fe line emitting gas is at a higher temperature and lower density than Si line emitting gas (Hwang, Hughes, & Petre 1998 and references therein). This is also expected as a result of the Ni bubble effect.
These considerations show that there is some evidence in Type Ia supernova remnants for clumping surrounding a large Ni bubble. The agreement between the knot velocities in Tychoโs remnant and those expected for the Ni bubble shell is probably fortuitous. There is some evidence for both Tycho and SN 1006 that they were subluminous supernovae (e.g., van den Bergh 1993; Schaefer 1996), which would imply a smaller amount of <sup>56</sup>Ni. For the subluminous Type Ia SN 1991bg, Mazzali et al. (1997, and references therein) estimate the synthesis of $`0.07M_{}`$ of <sup>56</sup>Ni extending to a velocity of $`5,000\mathrm{km}\mathrm{s}^1`$. This mass is consistent with the upper limit on the Fe mass ($`<0.16M_{}`$) found by Hamilton et al. (1997) for SN 1006.
The density contrast that is created by the nickel bubble is not well determined. In the calculations of Basko (1994), the density contrast at the edge of the bubble is a factor $`5`$ above the surroundings. However, in the case of a Type Ia supernova with a relatively much greater amount of radioactive Ni, clumps of non-radioactive material may be immersed in the radioactive Ni and subjected to a greater compression (Hรถflich 2000). We have considered a simple model for this process in order to estimate the compression. We assume that the expanding supernova gas has a constant density and that the free expansion during the first day is negligibly affected by the addition of energy, as above. The energy equation can be written as
$$\frac{1}{\gamma 1}\frac{dpV}{dt}=\frac{dQ}{dt}p\frac{dV}{dt},$$
(7)
where $`\gamma `$ is the adiabatic index, $`V`$ is the volume, and $`dQ/dt`$ is the heat addition from radioactivity. Over most of the evolution, the gas is radiation dominated and $`\gamma =4/3`$, although the thermal pressure may be significant in the initial phases. Provided $`t1`$ day, the power input from <sup>56</sup>Ni is approximately constant at $`4.8\times 10^{10}`$ ergs g<sup>-1</sup> s<sup>-1</sup>. Consideration of eq. (7) shows that in non-radioactive gas the pressure falls as $`t^4`$; this also initially applies to the radioactive gas, but it switches to $`t^2`$ evolution (at $`t2\times 10^3`$s for typical parameters). If the radioactive gas can keep non-radioactive gas at its pressure, the compression of the non-radioactive gas increases as $`t^{3/2}`$ and can reach $`200`$ after a day. However, a typical clump cannot be maintained in pressure equilibrium at the high pressure and a compression wave is driven into the clump from the outside. When the radiative diffusion timescale becomes less than the age, there is the possibility that the gas is more compressible and larger density inhomogeneities are created. This effect has been found in calculations of shock breakout from a star in supernova models (e.g., Chevalier 1981b). An important aspect of the effect is that radiation pressure dominates the gas pressure, which is the case here. These issues will have to be examined in detailed computations.
This discussion shows that it is possible that Type Ia supernovae do contain inhomogeneities which could be observed in a supernova remnant. The Ni bubble effect occurs on a timescale of 10 days, or a scale of $`10^{15}`$ cm for a velocity of $`10^4\mathrm{km}\mathrm{s}^1`$. The structure resulting from the Ni bubble should become frozen into the ejecta at this early time.
## 3 METHOD
We used the 2-D (two-dimensional) code ZEUS2D based on a finite difference scheme to carry out the hydrodynamical simulations (Stone & Norman 1992). The code uses an artificial viscosity to smooth shock transitions. We first ran 1-D simulations with an inner exponential density profile (ejecta) interacting with an outer constant density ambient medium, starting 2 weeks after the supernova explosion. The initial density distribution evolved into an intershock structure consisting of reverse-shocked ejecta in the inner region and a forward-shocked ambient medium on the outside, separated by a contact discontinuity (DC98). The 1-D intershock profile was then used to initiate 2-D simulations. The inner ejecta gas was freely-expanding, using an inflow boundary condition. The outer gas was at rest. Gas pressure in the unshocked gas was unimportant. We neglected the effects of magnetism, heat conduction, and radiation. Magnetic fields may play a role in suppressing heat conduction and viscosity, but are not expected to be dynamically significant in the supernova remnant. The radiative cooling time for the optically thin gas at an age $`300`$ yrs is greater than the age of the remnant. We assumed that the energy losses through X-ray emission are small so that the dynamics of the flow are not affected by gas cooling. We used an adiabatic index $`\gamma =5/3`$.
In 1-D runs, the starting dynamical age determines the initial ejecta-ISM interface position. Simulations with different initial ages showed that the evolution quickly converges within a few radial doubling times. The solution is unique, independent of when the simulation starts (see DC98). More generally, since the ejecta density profile only depends on its explosion mass and kinetic energy, the interaction with an ambient medium can be described by a set of scaling parameters $`R^{}`$, $`V^{}`$, $`T^{}`$ using $`M`$, $`E`$, and $`\rho _{am}`$ (see Truelove & McKee 1999; DC98)
$$R^{}=\left(\frac{3M}{4\pi \rho _{am}}\right)^{1/3}2.19\left(\frac{M}{M_{ch}}\right)^{1/3}n_0^{1/3}\mathrm{pc},$$
(8)
$$V^{}=\left(\frac{2E}{M}\right)^{1/2}8.45\times 10^3\left(\frac{E_{51}}{M/M_{ch}}\right)^{1/2}\mathrm{km}\mathrm{s}^1,$$
(9)
$$T^{}=\frac{R^{}}{V^{}}248E_{51}^{1/2}\left(\frac{M}{M_{ch}}\right)^{5/6}n_0^{1/3}\mathrm{yr},$$
(10)
where $`n_0=\rho _{am}/(2.34\times 10^{24}`$ gm cm<sup>-3</sup>) and $`\rho _{am}`$ is the ambient density. The dimensional variables $`r`$, $`v`$, $`t`$ can be expressed in terms of the nondimensional quantities: $`r^{}=r/R^{}`$, $`t^{}=t/T^{}`$, and $`v^{}=v/V^{}`$. Nondimensional solutions can be conveniently returned to dimensional ones by re-scaling, and new dimensional solutions for different $`M`$, $`E`$, and $`\rho _{am}`$ can be calculated. For a particular density distribution, one evolutionary sequence in the nondimensional variables represents all possible dimensional solutions.
For the exponential density profile, the 1-D solution (Fig. 1) entered the Sedov-Taylor self-similar blast wave evolution with $`rt^{0.4}`$ at $`t^{}1`$. As the ejecta density decreased with time, the reverse shock, initially moving outward in the stellar frame, began to move inward at $`t^{}2.5`$. It reached the stellar center at $`t^{}8`$.
## 4 INSTABILITIES FROM SMALL PERTURBATIONS
We simulated the Rayleigh-Taylor instabilities in the intershock region in two-dimensional spherical polar coordinates assuming $`\varphi `$-symmetry. The 2-D numerical grid was initialized with the 1-D solution at $`t^{}=0.00054`$ (0.134 yrs for the standard parameters $`E_{51}=1`$, $`M=M_{ch}`$, and $`n_0=1`$) across the angular domain, with perturbation seeds in density, pressure, and velocity placed between the contact discontinuity and the reverse shock wave. The grid was radially expanding, following the intershock boundaries until the reverse shock radius turned over to the stellar center. The evolution was tracked for five orders of magnitude in time. Simulations initiated earlier presented numerical problems because the higher density contrast across the contact discontinuity reduced the numerical time step determined by the Courant condition. The 2-D grid was moved based on the 1-D results. We used grid wiggling in the 1-D simulations, tracking the forward and the reverse shocks. In this process, the grid expands and contracts so that the shock locations are kept at the same grid numbers. We fit the evolution of the radii and velocities with a linear function and a logarithmic polynomial, respectively. We then applied the 1-D smooth fits on the 2-D grid boundaries, avoiding significant numerical noise coming from the wiggling scheme. The radial zone boundaries had a velocity distribution linearly varying with distance to the inner boundary, so that the zone spacing was kept uniform during the expansion. We used 500 zones in the radial direction and varied the number of angular zones; the radial zone spacing was $`3\times 10^6`$ $`R^{}`$ at $`t^{}`$=0.01, and increased to $`1.6\times 10^3`$ $`R^{}`$ at $`t^{}`$=1.73.
We investigated the growth of instabilities with various grid resolutions and perturbations. To reduce the computational expense, only a fraction of a quadrant centered at $`\theta =45^{}`$ was used. The run with the finest resolution had 2000 angular zones in 1/2 of a quadrant. The radial boundary conditions were inflow and outflow at the inner and outer sides, respectively. The angular boundaries were reflecting. We generally used a spherical harmonic function for the perturbation,
$$Y_{lm}(\theta ,\varphi )=P_l^m(\mathrm{cos}\theta )e^{im\varphi },$$
(11)
where $`m=0`$ and the harmonic $`l`$ (perturbation mode) is even, considering the $`\varphi `$-symmetry and the reflection symmetry about the equator. The angular perturbation is essentially the associated Legendre polynomial $`P_l^m`$, which has an increasing amplitude toward the polar axis. We applied perturbations such that the amplitude was between 1% and 50% near $`\theta =45^{}`$. The perturbed width was less than 40% of the distance between the reverse shock and the contact discontinuity, about 2.5% of the whole intershock width.
### 4.1 Evolution
Fig. (INSTABILITIES AND CLUMPING IN TYPE Ia SUPERNOVA REMNANTS) illustrates the evolution of the instability perturbed by a 1%, $`l=100`$ mode using 2000 angular zones on 1/2 of a quadrant centered on $`\theta =45^{}`$. Small perturbations in the density grow and form spikes that protrude from the contact discontinuity into the forward-shocked region. As the remnant expands, the Rayleigh-Taylor instabilities continue to grow to form mushroom-shaped caps while the stems become narrower. Secondary fingers develop at the same time, so that there are more fingers than the initial perturbation mode suggests. After reaching their maximum extent, the fingers fall to their sides and interfere with neighboring fingers. The forward growth of the caps is blocked by the drag of the flow and the fingers do not extend farther from the dense shell at the contact discontinuity.
In the later stages, the Kelvin-Helmholtz instability takes over at the mushroom caps, creating vortex rings. The relative motion of flows between a finger and its surroundings bends the stem and disrupts the flow. Since the effect of drag increases as the finger grows outward, the continual shedding of mass at the fingerโs top eventually leaves insuffient mass to overpower the drag of the countersteaming flow. The mushroom cap falls off to the side, and the remaining filaments are swept back. Vorticity develops in the less dense regions left by the original mushroom caps. Globally, the vortex rings gradually come to dominate over the spikes, and only the stems remain recognizable.
The Rayleigh-Taylor instability consistently builds up fingers with long wavelengths (low $`l`$). The dynamical stages of evolution can thus be distinguished by the dominant mode. We found four stages of evolution. The first stage is characterized by linear growth of the Rayleigh-Taylor instability; the perturbations grow out from the initial seeds and evolve towards a mushroom shape. Second, nonlinear growth of secondary fingers among primary ones becomes important. The fingers multiply and become congested. Third, mutual interference among fingers increases the wavelength of the dominant mode; vortex rings develop as the Kelvin-Helmholtz instability becomes active, although new fingers still develop, as in a convective roll. The instability is fully developed at this point. In the final stage, the intershock density contrast is reduced as the expansion decelerates; vortex rings and all the other features in the flow are dispersed, and the instabilities fade.
Figs. (2) and (3) show the angle-averaged density profiles with various resolutions at an early phase and at $`t^{}=1.6`$ (approximately the present epoch for Tycho). Compared to the 1-D solution, the reverse shock front smears out somewhat because it is corrugated, and the density peak at the contact discontinuity is not present. The forward shock front is less affected. Various resolutions give basically the same profile, although finer grid resolutions reveal larger fluctuations. The late phase shows less density contrast and the decline of instabilities.
We compare the exponential case to the power law case of $`n=7`$ with $`s=0`$ (interstellar medium) and $`s=2`$ (circumstellar wind) in Chevalier et al. (1992). In the power law case, the conditions are self-similar, so that the evolution is quasi-steady. Convection cells continually stir up a region near the contact discontinuity. The density structure remains qualitatively similar with time. The difference in the two $`n=7`$ cases with different ambient media is that for the $`s=2`$ case the density structure is characterized by groups of slender blobs with a pyramid outline, while for the $`s=0`$ case it is characterized by individual mushroom shapes. The reason is that a constant density surrounding medium causes a smaller density contrast across the the unstable region, and a negative gradient of entropy. Thus, the $`s=2`$ case allows the initial blobs to grow and stretch out into narrow fingers, contrary to the $`s=0`$ case in which the flow is blocked from propagating further and ends in a mushroom cap.
The $`s=0`$ power law case studied by Chevalier et al. (1992) shows that a larger power law index $`n`$ increases the dominant mode $`l`$ and the evolution goes from quasi-steady convection to intermittent growth of slender fingers, approaching in a sense the general profile of the low-$`n`$, $`s=2`$ cases. This is due to the larger density contrast between the shocked ambient medium and the shocked ejecta at larger $`n`$, and to the fact that the gradients in the shocked ambient region become shallower. The exponential case exhibits evolutionary properties similar to the power law case with decreasing index with time, in that longer wavelengths are built up throughout the evolution after the initial growth stage.
We show in Fig. (4) the approximate power law index $`n_{ej}=d\mathrm{ln}\rho /d\mathrm{ln}r`$ in the unshocked ejecta just inside the reverse shock for the exponential case. In the 2-D simulations, the reverse shock surface is corrugated, but its position does not deviate significantly from its position in the 1-D case. The position of the contact discontinuity found by dimensional analysis is $`R_ct^{(n3)/(ns)}`$ for a power law ejecta density profile and an ambient medium $`\rho r^s`$. The power law index for $`n=5`$ is the same as that for a Sedov-Taylor blast wave. For $`n5`$ the forward shock is still expected to expand with the Sedov law; however, the contact discontinuity expands more slowly than $`rt^{0.4}`$ and steadily distances itself from the forward shock in the self-similar frame (Chevalier 1982a). Self-similar solutions are no longer possible and the reverse shock starts to move inward towards the center. For the exponential case, we expect that the deceleration at the contact discontinuity is reduced when the reverse shock is in ejecta with $`d\mathrm{ln}\rho /d\mathrm{ln}r5`$. The instabilities should then lose the deceleration that drives them. The index $`n`$ starts to fall below 5 at $`t^{}0.325`$, while the turnover of the reverse shock actually occurs at $`t^{}2.5`$; the motion of the reverse shock takes time to respond to the changing density profile.
Both the density contrast and acceleration decline rapidly with time in the exponential case; the decline of the instabilities is an inevitable consequence. At the end of our simulation ($`t^{}=2.16`$), there is little continuing development of the instability.
### 4.2 Variation of the Perturbation and the Resolution
The evolution of the instability is insensitive to the initial linear perturbation. We have applied the initial perturbation in various areas: between the contact discontinuity and the forward shock, ahead the forward shock, and inside the reverse shock. After the growth stage, the sizes of the unstable regions are similar in all the cases. Perturbations placed at the shocked ejecta near the contact discontinuity are the most effective in exciting the instability.
On a low resolution grid of $`500\times 200`$ zones on 1/2 of a quadrant, a 10% perturbation with $`l=100`$ initially has 12 Rayleigh-Taylor fingers, which subsequently double in number. The $`l=50`$ mode generates three secondary fingers between every two primary fingers, resulting in 24 fingers like the $`l=100`$ mode (Fig. INSTABILITIES AND CLUMPING IN TYPE Ia SUPERNOVA REMNANTS). There appears to be a preferred mode shortly after the linear growth, independent of the initial conditions. However, with finer resolution more details between the primary fingers appear and the presence of a dominant mode becomes vague. More highly resolved fingers appear to be more slender and sharp, and the overall flow patterns thus appear different. On finer grids, the flow pattern approaches the pyramid shape of an $`n=7`$, $`s=2`$ power law density case; mushroom structures become congested and twisted at the base near the contact discontinuity. Toward the end of the evolution, a preferred mode in increasing wavelength is built up independent of grid coarseness or the initial conditions. The presence of the late-stage dominant mode can be understood as the exponential model having a continually decreasing power law index, although limited resolution can also fabricate a longer wavelength dominant mode. Our results are consistent with the studies of Rayleigh-Taylor instability growth in a stellar explosion by Fryxell, Mรผller and Arnett (1991), who used a PPM (piecewise parabolic method) code and cylindrical coordinates.
The morphology of the fingers on our higher resolution runs is similar to that of Kane, Drake, & Remington (1999) using the PPM code PROMETHEUS. The results of Chevalier et al. (1992) using the piecewise parabolic code VH-1, Jun & Norman (1996a) using ZEUS2D, and Fryxell et al. (1991) using PROMETHEUS appear to have lower resolution than our results and do not show the narrow structure that we find. The high resolution calculation of Mรผller, Fryxell, & Arnett (1991) is closer to the morphology that we obtain. Higher resolution calculations introduce shorter wavelength numerical noise and allow small perturbations to develop. Small scale structures have a faster growth rate of the Rayleigh-Taylor instability in the linear regime, and this can be seen in our simulations. Fig. (INSTABILITIES AND CLUMPING IN TYPE Ia SUPERNOVA REMNANTS) shows that the initial growth of Rayleigh-Taylor fingers leads to fingers with a larger extent at high resolution, presumably because the smaller cross section gives less drag. However, this trend may begin to reverse at the highest resolution because the stems of the Rayleigh-Taylor fingers become more unstable (see $`t^{}=0.003`$ in Fig. INSTABILITIES AND CLUMPING IN TYPE Ia SUPERNOVA REMNANTS). As mentioned above, the extent of the unstable region shows little dependence on resolution in the fully developed regime. On finer grids, the fingers in the nonlinear regime show an increasing tendency to bend, which prompts stronger interaction among fingers and produces more complex structures. As the resolution increases, the mushroom caps break up into filaments; this trend continues at our highest resolution. The clear mushroom shapes found in lower resolution calculations may be an artifact of the limited resolution.
In independent work that was concurrent with our study, Dwarkadas (2000) simulated the instabilities resulting from the interaction of supernova ejecta (with an exponential profile) with a surrounding medium. He presented the evolution of the instability for a case in which it grew from numerical noise. This perturbation is much smaller than the cases that we considered and the instability remained in the linear regime until the late damping stage of the instability. When he initiated the simulations with a 2 % perturbation, his results are consistent with ours.
We have also simulated the evolution for an initial perturbation with 100% amplitude. After a transition phase during which the intershock region became more than twice as broad as in the standard case, the evolution returned to the standard case. As found by Kane, Drake, & Remington (1999) in simulations with a power law supernova density profile, the standard evolutionary track is robust.
## 5 EVOLUTION OF CLUMPS
As discussed in ยง 1, recent X-ray (Hughes 1997; Hwang & Gotthelf 1997) and radio (Velรกzquez et al. 1998) observations of Tychoโs remnant show two knots juxtaposed near the eastern edge. The outer outline of radio (Velรกzquez et al. 1998) and X-ray (Vancura et al. 1995) emission, which presumably defines the forward shock front, shows an outward protrusion surrounding the knots. The properties of undecelerated motion (Hughes 1997) and an outer shock protrusion cannot be explained by instabilities generated by linear pertubations. The Rayleigh-Taylor fingers from linear perturbations are substantially decelerated and do not affect the outer shock front. Although our computations are two-dimensional, we believe that going to three dimensions would not change this conclusion. Rayleigh-Taylor fingers are known to show somewhat greater growth in three dimensions during the initial growth phase (e.g., Kane et al. 2000 and references therein). However, in computations with power law supernova ejecta, Jun & Norman (1996b) find that the thickness of the mixing layer is very similar in two and three dimensions when the instability is fully developed (see their figs. 2 and 3). The instability is probably limited by the properties of the interaction region. Jun & Norman (1996b) do find that the stellar ejecta show smaller scale structure in three dimensions.
We believe that the origin of the knots is best explained by nonlinear clumps of supernova ejecta expanding into the intershock region. We thus simulated the hydrodynamic evolution of ejecta clumps interacting with the intershock gas. We represented clumps as denser spheres superposed on the smooth exponential density profile. In the simulations, we initialized a clump at the polar angle $`\theta =45^{}`$ on 2-D computations like those described in the previous section. In two dimensions, the sphere is in fact a 3-D torus around the polar axis. A 3-D sphere can be simulated by placing the clump on the polar axis $`\theta =0`$ on a 2-D grid; however, there are then singularity problems on the symmetry axis of the grid, which lead to long Rayleigh-Taylor fingers on the axis due to the instabilities discussed in the previous section. Comparisons of the hydrodynamics of ISM clouds in spherical and toroidal morphologies indicate that geometry contributes little to the interaction (Jun & Jones 1999). To examine the effect of geometry, we carried out clump interaction simulations with clumps at $`\theta =15^{}`$ and $`\theta =80^{}`$; the clump evolution did not substantially change.
Under the same boundary conditions as used for the instability computations, the clump was included in the freely expanding ejecta. The clump first ran into the reverse shock and then moved forward through the intershock region. The initial clump as well as the ejecta were cold compared to the high temperature shocked ejecta. The gas pressure was unimportant, so that the size and the density contrast of the clump remained unchanged until the clump expanded into the reverse shock front.
The basic physics of the interaction is similar to that for the interaction of nonradiative blast waves with an interstellar cloud (Klein et al. 1994). As the shock wave moves past the cloud, a reflected wave moves back into the shocked intercloud gas and creates a bow shock. The incident shock also creates a transmitted โcloud shockโ that moves into the cloud and crushes it. The cloud shock propagates with a velocity $`v_c(\rho _i/\rho _c)^{1/2}v_s`$, where $`\rho _i`$, $`\rho _c`$, and $`v_s`$ are the density of the intercloud medium, the density of the cloud, and the shock velocity in the intercloud medium, respectively. A larger density contrast between the cloud and the intercloud medium causes a smaller velocity ratio of $`v_c/v_s`$, which helps the development of a shear flow, and consequently the Kelvin-Helmholtz instability at the cloud-intercloud interface. The Richtmyer-Meshkov instability powered by the impulsive acceleration takes place due to the impact of the shock wave on the upstream side of the cloud-intercloud surface. When the cloud shock exits the cloud, a rarefaction wave moves back into the cloud and causes expansion. This acceleration leads to the Rayleigh-Taylor instability on the upstream side of the cloud. The combined instabilities lead to the destruction of the cloud on a timescale of several times of the cloud-crushing time, $`t_{cc}=R_c/v_c`$.
The clump-young supernova remnant interaction is more complicated than the cloud-blast wave interaction because of the structure of the intershock structure, including the pre-existing large scale instability. The most important factor measuring the crushing strength of a cloud-shock interaction is the impact area per unit mass, which is inversely proportional to the initial density contrast, $`\chi =\rho _c/\rho _i`$. We have explored the evolution of a single clump, varying its size, density contrast, and initial impact time with the reverse shock. We examined four initial impact times, with the clump density contrast ranging between 3 and 100, and the clump size $`a_0`$ (as a fraction of the intershock width) below 1/3 (Table 1). We estimate the size of the Tycho X-ray knots to be about 1/5 the intershock width, although this is uncertain because the clumps are poorly resolved and we see only shocked parts of the clumps. For a size below 30% of the intershock width at the corresponding interaction time, we found that (a) an early clump initiated at $`t^{}=0.011`$ must have a density contrast $`\chi 10`$ to cause a protrusion on the remnant outline; a clump with $`\chi =30`$ causes a protrusion but the effect would have subsided long before the present epoch for Tycho; (b) a clump initiated at $`t^{}=0.22`$ with $`\chi =50`$ has a bulging effect on the remnant outline, but the outline returns to spherical symmetry by $`t^{}=1`$; (c) a later clump initiated at $`t^{}=0.86`$ with $`\chi =100`$ is not dense enough to reach the forward shock at present epoch $`t^{}=1.6`$ for Tycho. A light clump with $`\chi =3`$ develops vorticity when encountering the reverse shock and is quickly destroyed.
Figs. (INSTABILITIES AND CLUMPING IN TYPE Ia SUPERNOVA REMNANTS) and (INSTABILITIES AND CLUMPING IN TYPE Ia SUPERNOVA REMNANTS) show the interaction with a single clump for $`\chi =100`$ and two different interaction times. After passing through the reverse shock, the cloud gradually became flattened and curved like a crescent. Material streamed out from the horns of the crescent; the ram pressure difference between the axis and the side of the cloud drove the mass loss. The pressure near the cloud axis was higher because of the additional ram pressure on the front face of the clump. At the rear of the clump the flow became turbulent and left a trail of vorticies. The clump expanded laterally as it approached the forward shock. The front of the clump snowplowed the material ahead. The forward shock wave was distorted and a bulge formed. The material on the shock front became distributed into two lumps that gradually receded from each other. As the expansion of the supernova remnant continued, the bulge as well as the crescent clump immediately behind it lost their identities. Even a clump with $`\chi =100`$ could not penetrate the forward shock when initiated at a later time (Fig. INSTABILITIES AND CLUMPING IN TYPE Ia SUPERNOVA REMNANTS), but snowplowed the material ahead and caused a protrusion.
The velocity of a shocked clump is determined by the drag of the surrounding material. A shocked clump becomes comoving with the postshock flow in approximately the time that the cloud sweeps up a column density of intercloud gas equal to its initial column density. The drag time is proportional to $`t_{drag}\chi ^{1/2}t_{cc}`$, so that a denser clump travels faster in the remnant. The lateral expansion of a shocked clump into a cresent shape increases the cross section area and the drag, significantly decelerating the clump.
For an exponential density profile of the ejecta, the clump evolution depends on the time of clump interaction because of the evolution of the supernova remnant. An earlier clump requires less density contrast to have a protrusion effect on the remnant outline (Fig. INSTABILITIES AND CLUMPING IN TYPE Ia SUPERNOVA REMNANTS). Early interactions give the clump a crescent shape without the action of the Rayleigh-Taylor instability on the clumpโs front because the higher flow velocity in the remnant delays the rarefaction wave traveling back upstream as the cloud-shock exits the clump. Later clumps tend to develop instabilities on the clump surface. The exponential model gives a larger velocity difference and a larger density difference between the ejecta and the reverse shock front at earlier epochs. The early clumps rapidly move into a lower density medium and are thus more robust, for a given value of $`\chi `$.
For $`n_o=0.8`$ cm<sup>-3</sup>, $`E=10^{51}`$ ergs, and $`M=1.4M_{}`$ (DC98), Tychoโs remnant has $`t^{}=1.6`$ at its present age of 427 years. If this model applies and the knots are in undecelerated motion, they must have crossed the reverse shock front at $`t^{}1.2`$. This places Tycho in the regime where $`\chi 100`$ is required.
In addition to the isolated knots, we considered groups of clumps. The Ni bubble process is likely to create a shell of clumpy ejecta and, as discussed in ยง 6, the observations of ejecta in Tychoโs remnant suggest that widespread clumpiness may be present. In Fig. (INSTABILITIES AND CLUMPING IN TYPE Ia SUPERNOVA REMNANTS) we show two simulations with multiple clumps. In the top one, there is one band of clumps and in the second, there are two bands. The evolution of each forward clump is similar to that in the isolated clump case, but they combine to push the forward shock wave to larger radius. The ejecta move out to a relatively large radius, as may be required in Tychoโs remnant.
## 6 THE REMNANT OF TYCHOโS SUPERNOVA
As discussed in ยง 1, Tychoโs supernova was likely to be of Type Ia. The expansion rate of the remnant has been measured in the radio, giving an expansion parameter $`m=0.47\pm 0.05`$ (Strom, Shaver, & Goss 1982), $`0.462\pm 0.024`$ (Tan & Gull 1985), and $`0.471\pm 0.028`$ (Reynoso et al. 1997). The optical filament yields an expansion parameter of $`0.39\pm 0.01`$ (Kamper & van den Bergh 1978), close to the adiabatic blast wave case. The slight discrepancy in these measurements can be explained by sampling variance; the optical measurement samples only the densest regions where the emissivity is the highest. HI observations do indicate a higher surrounding density to the NE (Reynoso et al. 1999). These radio expansion values indicate a global pre-Sedov stage and exclude a surrounding circumstellar medium, which would give an expansion parameter $`0.7`$ (Chevalier 1982b; DC98). The exponential model is successful in reproducing the observed shock position and deceleration of Tychoโs remnant. At Tychoโs present age of 427 years ($`t^{}=1.7`$) the model places the reverse shock and the forward shock at radii of 2.08 pc ($`r^{}=0.98`$) and 3.09 pc ($`r^{}=1.45`$), while they decelerate with expansion parameters of $`m=0.15`$ and $`m=0.47`$, respectively.
The radio observations of Tycho by VLA by Reynoso et al. (1997) show that the NE and particularly the adjacent SE parts protrude from the circular outline. The SE protrusion corresponds to the two X-ray knots. Near the sharp NE edge, Velรกzquez et al. (1998) note the presence of regularly spaced structure. The spatial regularity is interpreted by Velรกzquez et al. as Rayleigh-Taylor fingers still in the linear regime with a mode $`l`$ 30; they attribute the preferred wavenumber to the effects of viscosity. However, the importance of viscosity is uncertain because magnetic fields can inhibit its effect. We believe that a linear regime for the Rayleigh-Taylor instability is unlikely. Regardless of the supernova density model, the shock front has expanded by a sufficently large factor to allow the instabilities to evolve to their saturated phase if there are initial perturbations with amplitude $`1`$%.
The regular structure observed in the radio image by Velรกquez et al. (1998) is suggestive of a Rayleigh-Taylor instability, but there are problems with this interpretation in the context of our models. One is that at $`t^{}=1.6`$, the region occupied by the unstable ejecta does not extend beyond 85 % of the remnant radius. The structure observed by Velรกquez et al. is typically at a larger radius and in fact extends to the edge of the remnant. Another problem is that the instability at this late stage does not have a clear regular finger structure in our models, but is dominated by vortex rings with an irregular distribution. The large radius of the structure suggests that it may be connected to perturbations in the surrounding medium as opposed to the ejecta. However, the outer outline of the shock front is smooth on the scale of the structure, as imaged at both optical and radio wavelengths. In addition, the X-ray image of Tychoโs remnant (e.g., Vancura et al. 1995), which is dominated by line emission, shows strong emission at a large fractional radius. Because the image is likely to be dominated by emission from the ejecta (Vancura et al. 1995; Hwang & Gotthelf 1997), the evidence is that the outer structure is related to the ejecta.
We believe that a solution to this problem is indicated by the presence of the X-ray knots with undecelerated motion. Two knots are clearly observed in one section, but it is likely that they represent a more widespread phenomenon. We suggest that there was a spherical shell of clumpy ejecta at a velocity $`6,0008,000\mathrm{km}\mathrm{s}^1`$ in the freely expanding ejecta that gives rise to the clumpy structure. Kane et al. (1999) previously suggested that the radio structure observed in Tycho near the forward shock wave may be related to nonlinear density variations in the freely expanding ejecta. There are reasons to believe that the clumps are restricted to a region of the ejecta and not spread through all velocities. First, the deceleration parameter measured at X-ray wavelengths is significantly larger than that measured at radio or optical wavelengths (Hughes 1997). If the clumps are widespread, they should come and go in a band that expands as does the emission at lower wavelengths. However, a band of clumps could give the observed difference in deceleration parameters. Another reason is that in SN 1006, the Fe II absorption line profiles are consistent with a distribution of gas at velocities $`5,000\mathrm{km}\mathrm{s}^1`$ like that in the W7 model (Fesen et al. 1988), which has an approximately exponential density distribution. The Si line shows absorption at redshifted velocities $`3,0007,000\mathrm{km}\mathrm{s}^1`$, which can be attributed to moderately high velocity Si interacting with the surrounding medium. There is no blueshifted Si absorption. The data are consistent with the picture of smoothly distributed Fe on the inside, outside of which is Si gas that may be clumpy.
The presence of a band of clumps is expected in the Ni bubble scenario for clump formation (see ยง 2). One prediction of this model is that the outer ejecta emission should be very clumpy when observed at high X-ray resolution (see Fig. (INSTABILITIES AND CLUMPING IN TYPE Ia SUPERNOVA REMNANTS). In addition to this component, there should be emission from a smaller radius, near the reverse shock front. This emission should be representative of lower density, higher temperature gas. We believe that the emission in the Fe K line, which is at a smaller radius (Hwang & Gotthelf 1997), is primarily from gas close to the reverse shock wave. This gas should have some structure because of the earlier instabilities, but much less than the outer clumpy ejecta.
Our computations showed that in order for knots to survive to close to the outer shock wave, they must have an initial density contrast with respect to their surroundings of $`100`$. Hwang et al. (1998) in fact find that the ionization time ($`n_et`$, where $`t`$ is the time since the gas was heated) is 100 times larger for Si than for Fe, based on line emission from the entire remnant. The time since heating should be comparable for the two elements, so this result is consistent with the Si being primarily in shocked clumps (created outside the Ni bubble) and the Fe primarily in interclump gas (inside the Ni bubble). However, the situation is complicated by the fact that there is evidence for some Fe in knots, which we argue had a nonradiogenic origin (e.g., <sup>54</sup>Fe). In addition, we estimated in ยง 2 that $`0.13M_{}`$ of Si could be swept into a shell around the Ni bubble, so there could be some Si emission from the ejecta that is not clumped. Investigation of these issues will require spatially resolved X-ray spectroscopy, as will be possible with the Chandra and XMM observatories.
## 7 IMPLICATIONS FOR TYPE Ia SUPERNOVA MODELS
One result from observations of the Type Ia remnants is a lower limit on the velocity of Si in the explosion. For Tychoโs supernova, there is an undecelerated knot of Si moving at $`8,300(D/2.5\mathrm{kpc})\mathrm{km}\mathrm{s}^1`$, where $`D`$ is the distance. In SN 1006, there is Si freely expanding with a minimum velocity in the range $`5,6007,000\mathrm{km}\mathrm{s}^1`$. In both of these cases, the velocities appear to be less than the minimum velocity of the Si-dominant layer in the W7 deflagration model, about $`10,000\mathrm{km}\mathrm{s}^1`$ (Thielemann et al. 1986). However, if these are the remnants of subluminous supernovae, the W7 model may not be the appropriate one. Also, the initial burning front is unlikely to be precisely spherically symmetric, so that some material that has not burned to <sup>56</sup>Ni may lag behind.
Our most significant results involve the evidence for clumping in Type Ia supernovae. We have found that density contrasts $`100`$ are necessary to reproduce the large radius of X-ray emitting clumps and their velocities. A plausible mechanism for the formation of the clumps is the action of an expanding Ni bubble, although calculations of this process have not yet shown that such large compressions are feasible. We estimate that $`0.1M_{}`$ of matter can be swept into clumps around the Ni bubble. For the two undecelerated X-ray emitting clumps in Tychoโs remnant, Hughes (1997) estimates a mass of order $`0.002M_{}`$ for the Si $`+`$ S clump and $`0.0004M_{}`$ for the Fe clump, so there can be many comparable clumps. The presence of Fe in a clump requires that it have a non-radioactive origin. The 1-dimensional calculations of nucleosynthesis associated with the W7 Type Ia model (Thielemann et al. 1986; Iwamoto et al. 1999) show that <sup>54</sup>Fe is present along with Si and S in the region immediately outside of the <sup>56</sup>Ni region. However, the <sup>54</sup>Fe does not dominate the composition, unlike what is indicated by the observations of the Fe clump in Tycho. The central region of the W7 model is composed of stable <sup>56</sup>Fe and <sup>54</sup>Fe, but this gas is at too low a velocity to be compatible with the Tycho clump. Thus, the synthesis of fast material dominated by non-radioactive Fe remains to be demonstrated.
## 8 CONCLUSIONS
We used 2-D hydrodynamical simulations in spherical polar coordinates to investigate instabilities and clumpiness in a Type Ia supernova remnant. We adopted an exponential density profile for the smooth supernova ejecta and a constant density for the surrounding (interstellar) medium. The exponential density profile gives the best general approximation for the supernova ejecta in detailed explosion models, as opposed to a power law density profile for which self-similar solutions are available. The region between the forward shock in the interstellar medium and the reverse shock in the ejecta is Rayleigh-Taylor unstable. An initial perturbation grows to nonlinear strength, develops an increasing characteristic wavelength, and fades as the remnant evolves to the Sedov blast wave regime. The characteristic Rayleigh-Taylor mushroom caps are replaced by vortex rings from the Kelvin-Helmholtz instability, and turbulence appearing in vorticies remains dominant throughout the rest of the evolution. In these late phases, the depletion of the kinetic energy in the inner ejecta as the ejecta density profile continually weakens the acceleration needed to power the Rayleigh-Taylor instability. A sequence of simulations with increasing grid resolution shows the Rayleigh-Taylor fingers become increasingly slender and sharp and tend to bend, limited only by the grid resolution.
The structure that is formed as a result of the Rayleigh-Taylor instability is decelerated and is confined to a region $`85`$% of the remnant radius for a remnant with the dynamical age of Tychoโs remnant. These properties are not consistent with the knots observed in Tychoโs remnant. We thus considered a clump in the diffuse supernova ejecta with the aim of reproducing the properties of the the SE X-ray knots and the related protrusion on the forward shock outline. In the clump-shock interaction, the density contrast, clump size, and time of initiation of the interaction are the major factors. Ram pressure stripping gives rise to a core-plume structure and a strong Kelvin-Helmholtz instability in the downstream region. The Rayleigh-Taylor instability develops on the clumpโs upstream side, which facilitates fragmentation of the clump. The clump causes a bulge on the remnant outline as the ram pressure pushes material ahead. The bulge eventually disappears as the clump is fragmented and swept back. Punching through the forward shock does not occur for our similution with conditions like those in Tychoโs remnant, even for the case of a clump with an initial density contrast of 100. Clump interaction in the early phases of the remnant evolution are more likely to distort the forward shock front because of the larger density contrast between the ejecta and the interstellar medium.
The two X-ray clumps to the SE in Tycho are probably not alone. In fact, we find that the large radial position of most of the ejecta emission suggests that much of the emission is from clumps. If it were not in clumps, it could not extend to the observed radius. We thus suggest that the undecelerated knots are at the high column density end of a spectrum of knot properties. The required knot compression factor, up to $`100`$, is not a natural result of existing Type Ia supernova models. We suggest that the expansion of a radioactive Nickel bubble is responsible for the compression of the knots. The compression factor is larger than that found in existing calculations, but it is possible that radiative transfer effects during the final stages of the Ni bubble expansion can lead to enhanced compression. The Ni bubble interpretation has the implication that Fe clumps, which are observed, cannot have been synthesized as <sup>56</sup>Ni. They must have a non-radiogenic origin, such as <sup>54</sup>Fe. The synthesis of high velocity matter in which <sup>54</sup>Fe is the dominant component has yet to be shown in computations of nucleosynthesis.
In addition to these suggested theoretical investigations, advances can be expected on the observational front. Spatially resolved X-ray spectra of Type Ia supernova remnants, as is possible with the Chandra and XMM observatories, will be valuable for showing whether our picture of ejecta clumps is valid. We predict that the outer ejecta emission in Tychoโs remnant has a clumpier structure than the inner emission observed in the Fe K line.
C.-Y. Wang thanks John Blondin and John Hawley for their assistance with the numerical simulations, and Fred Lo and Academia Sinica for providing their facilities. R.A.C. is grateful to Vikram Dwarkadas, Peter Hรถflich, Phil Pinto, and Craig Wheeler for useful discussions and correspondence. The computations were carried out on the Cray T90 of NPACI, and the IBM SP2 at University of Virginia. Support for this work was provided in part by NASA grant NAG5-8232.
Series of density images illustrating the time evolution of the dynamical instability. The perturbation was imposed at $`t^{}=0.00054`$ with 1% amplitude and an $`l=100`$ mode of the spherical harmonic function. The grid has 500 radial zones by 2000 angular zones in 1/2 of a quadrant. The plots at early stages display the central 1000 angular zones in 1/4 of a quadrant, while the plots shown for later stages display 1/2 of a quadrant. The contour levels are exponentially spaced between the lowest value and the highest one sampled in the grid domain.
Series of density contours with varying resolutions for four epochs: $`t^{}=0.0016`$, $`0.0027`$, $`0.011`$, and $`0.11`$. The lowest resolution uses 200 angular zones on 1/2 of a quadrant and the highest uses 2000 zones on 1/2 of a quadrant. The top two low resolution plots exhibit the same dominant modes on 1/2 of a quadrant, growing from the initial perturbation modes $`l=100`$ and $`l=50`$. The other plots are shown on 1/4 of a quadrant. All of the plots use a 10% perturbing amplitude.
Snapshots of a clump in the ejecta expanding into the shocked region. The clump was initiated at $`t^{}=0.207`$, with an initial density contrast $`\chi =100`$ and a size $`a_0=1/8`$. The clump becomes crescent shaped and causes a protrusion in the forward shock front. The grid has 500 radial by 500 angular zones on 1/2 of a quadrant.
Snapshots of a clump with an initial density contrast $`\chi =100`$ and a size $`a_0=1/6`$ in the ejecta expanding into the shocked region. The interaction was initiated at $`t^{}=0.866`$. The Rayleigh-Taylor instability can be seen on the clumpโs downstream side. The grid has 500 radial by 500 angular zones on 1/2 of a quadrant.
Snapshots of multiple clumps with a initial density contrast $`\chi =50`$ and a size $`a_0=0.15`$ in the ejecta expanding into the shocked region. The interaction was initiated at $`t^{}=0.866`$. The grid has 500 radial by 500 angular zones on 1/2 of a quadrant. The top plot has one row of clumps and the bottom has two rows.
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# On the space of Fredholm operators
## Introduction
The work of Atiyah and Singer on the index of elliptic operators on manifolds has singled out the role of the space of bounded Fredholm operators in topology. It is a classifying space for a very useful functor, the topological $`K`$-theory. This means that a continuous family $`(L_x)_{xX}`$ of elliptic pseudo differential operators parameterized by a compact $`CW`$-complex $`X`$ naturally defines an element in the group $`K(X)`$, the index of the family.
In most examples, the elliptic operators are not bounded operators and thus the notion of continuity has to be defined carefully. The operator theorists have come up with a quick fix. The family $`xL_x`$ of Fredholm operators is called Riesz continuous if and only if the families of bounded operators
$$xL_x(1+L_x^{}L_x)^{1/2},xL_x^{}(1+L_xL_x^{})^{1/2}$$
are continuous with respect to the operator norm. In concrete applications this approach can be a nuisance. For example, consider as in a Floer family of elliptic boundary value problems (parameterized by $`sS^1`$)
$$u(t):[0,1],s[0,2\pi ]\{\begin{array}{ccc}\frac{du(t)}{dt}+a(t)u(t)=0& \mathrm{if}& t(0,1)\\ & & \\ u(0)& \mathrm{and}& e^{๐ขs}u(1)\end{array}$$
($`BV_s`$)
where $`a:[0,1]`$ is a given smooth function. This family ought to be considered continuous but verifying the above definition can be quite demanding. The first technical goal of this paper is to elucidate this continuity issue.
As observed in , for $`K`$-theoretic purposes it suffices to investigate only (possibly $`_2`$-graded) selfadjoint operators (super-)commuting with some Clifford algebra action. For example, the space of Fredholm operators on a Hilbert space $`H`$ can be identified with the space of odd, selfadjoint Fredholm operators on the $`_2`$-graded space $`HH`$ via the correspondence
$$L\left[\begin{array}{cc}0& L^{}\\ L& 0\end{array}\right].$$
That is why we will focus exclusively on selfadjoint operators.
In we have argued that in many instances it is much more convenient to look at the graphs of Fredholm selfadjoint operators on a Hilbert space $`H`$. If $`T`$ is such an operator and $`\mathrm{\Gamma }_THH`$ is its graph, then $`\mathrm{\Gamma }_T`$ is a Lagrangian subspace of $`HH`$ (with respect to a natural symplectic structure) and moreover, the pair $`(H0,\mathrm{\Gamma }_T)`$ is Fredholm. As shown in , the space of Fredholm pairs of Lagrangian subspaces is a classifying space for $`KO^1`$. (A similar description is valid for all the functors $`KO^n`$; see .)
A natural question arises. Suppose that two families of subspaces determined by the graphs of two families of Fredholm operators are homotopic inside the larger space of Fredholm pairs of Lagrangian subspaces. Can we conclude that the corresponding families of Fredholm operators are also homotopic inside the smaller space of operators?
The is the second issue we want to address in this paper. We will consider various topologies on the space of closed, unbounded Fredholm operators and analyze when the above graph map $`T\mathrm{\Gamma }_T`$ from operators to subspaces is a homotopy equivalence. Surprisingly, to answer this question we only need to decide the continuity of Floer type families of boundary value problems. The symplectic reduction technique developed in coupled with the Bott periodicity will take care of the rest.
The paper consists of three sections. In Section 1 we compare two topologies on the space of unbounded Fredholm operators: the gap topology, given by the gap distance between the graphs, and the Riesz topology, described above. In the second section we prove a general criterion (Proposition 2.1) for recognizing when a family of boundary value problems, such as ($`BV_s`$), is continuous with respect to the Riesz topology. In the last section we address the connections with $`K`$-theory.
Acknowledgments This paper addresses some subtle omissions in . I am grateful to Bernhelm Booss-Bavnbek for his warm reception at Roskilde University, where the ideas in this paper were born, and for the lively discussions concerning the results of .
## 1 Topologies on the space of selfadjoint operators
Let $`H`$ be a separable real Hilbert space. Denote by $`๐ฎ`$ the space of densely defined, selfadjoint operators on $`H`$ and by $`๐ฎ`$ the space of bounded selfadjoint operators $`T:HH`$. Set
$$[๐ฎ]:=\{T๐ฎ;T<1\}.$$
The Riesz map is the bijection
$$\mathrm{\Psi }:๐ฎ[๐ฎ],AA(1+A^2)^{1/2}.$$
There are two natural metrics on $`๐ฎ`$: the gap metric
$$\gamma (A_0,A_1):=(๐ข+A_0)^1(๐ข+A_1)^1+(๐ขA_0)^1(๐ขA_1)^1,$$
and the Riesz metric
$$\rho (A_0,A_1):=\mathrm{\Psi }(A_0)\mathrm{\Psi }(A_1).$$
###### Remark 1.1.
According to \[4, Thm. IV.2.23\] we have $`\gamma (A_n,A)0`$ if and only if
$$\delta (\mathrm{\Gamma }_{A_n},\mathrm{\Gamma }_A)0$$
where $`\mathrm{\Gamma }_T`$ denotes the graph of the linear operator $`T`$ and $`\delta `$ denotes the gap between two closed subspaces.
###### Lemma 1.2.
The identity map $`(๐ฎ,\rho )(๐ฎ,\gamma )`$ is continuous.
Proof Observe that for every $`A๐ฎ`$ we have
$$\frac{1}{๐ข\pm A}=\frac{A๐ข}{1+A^2}=\frac{A}{1+A^2}\frac{1}{1+A^2}=\frac{1}{(1+A^2)^{1/2}}\mathrm{\Psi }(A)๐ข\frac{1}{1+A^2}$$
and
$$\frac{1}{1+A^2}=1\mathrm{\Psi }(A)^2$$
so that $`\mathrm{\Psi }(A_n)\mathrm{\Psi }(A)0`$ implies $`(๐ข\pm A_n)^1(๐ข\pm A)^10`$. $`\mathrm{}`$
Denote by $`๐`$ the $`C^{}`$-algebra of continuous functions $`f:`$ such that the limits
$$f(\pm \mathrm{}):=\underset{\lambda \pm \mathrm{}}{lim}f(\lambda )$$
exist. Denote by $`๐_0`$ the subalgebra defined by the condition
$$f๐_0f(\mathrm{})=f(\mathrm{}).$$
Define $`P_0,P_\pm ๐_0`$ by
$$P_0(\lambda )1,P_\pm (\lambda )=(\lambda \pm ๐ข)^1.$$
The Stone-Weierstrass approximation theorem shows that the algebra $`๐ซ`$ generated by $`P_0,P_\pm `$ is dense in $`๐_0`$.
The functional calculus for selfadjoint operators show that any $`A๐ฎ`$ defines a continuous morphism of $`C^{}`$-algebras
$$๐๐ฎ,ff(A).$$
###### Proposition 1.3.
The following statements are equivalent.
(i) $`\gamma (A_n,A)0`$.
(ii) $`f(A_n)f(A)0`$, $`f๐_0`$.
Proof Clearly (ii) $``$ (i) since $`P_\pm ๐_0`$ and
$$\gamma (A_n,A)=P_{}(A_n)P_{}(A)+P_+(A_n)P_+(A).$$
To prove (i) $``$ (ii) we use an idea in \[7, Chap. VIII\]. Clearly if $`\gamma (A_n,A)0`$ then
$$P(A_n)P(A)0,P๐ซ.$$
Fix $`f๐_0`$. Since $`๐ซ`$ is dense in $`๐_0`$, for every $`\epsilon >0`$ we can find $`P๐ซ`$ such that $`fP\epsilon /3`$ and then $`n(\epsilon )>0`$ such that, $`nn(\epsilon )`$ such that
$$P(A_n)P(A)\epsilon /3.$$
Then, $`nn(\epsilon )`$ we have
$$f(A_n)f(A)f(A_n)P(A_n)+P(A_n)P(A)+P(A)f(A)\epsilon .\mathrm{}$$
###### Proposition 1.4.
Fix a function $`\alpha ๐`$ such that $`\alpha (\lambda )1`$ for $`\lambda 1`$ and $`\alpha (\lambda )0`$ if $`\lambda 1`$. Then the following statements are equivalent.
(i) $`\rho (A_n,A)0`$
(ii) $`f(A_n)f(A)0`$, $`f๐`$.
(iii) $`\gamma (A_n,A)0`$ and $`\alpha (A_n)\alpha (A)0`$.
Proof Define $`r๐`$ by
$$r(\lambda ):=\frac{\lambda }{(1+\lambda ^2)^{1/2}}.$$
The equivalence (i) $``$ (ii) follows exactly as in the proof of Proposition 1.3 using Lemma 1.2 and the fact that the subalgebra spanned by $`๐_0`$ and $`r`$ is dense in $`A`$. The equivalence (ii) $``$ (iii) relies on Proposition 1.3 and the fact that the algebra spanned by $`๐_0`$ and $`\alpha `$ is dense in $`๐`$. $`\mathrm{}`$
###### Remark 1.5.
(B. Fuglede) The topological spaces $`(๐ฎ,\rho )`$ and $`(๐ฎ,\gamma )`$ are not homeomorphic. Using Proposition 1.4 it is easy to construct an example of a sequence $`A_n\stackrel{\gamma }{}A`$ such that $`A_n`$ does not converge to $`A`$ in the Riesz metric. More precisely consider the space
$$\mathrm{}^2=\{(x_j)_{n1};x_j,\underset{j}{}x_j^2<\mathrm{}\}$$
with canonical Hilbert basis $`๐_1,๐_2,\mathrm{}`$. For $`n=0,1,2,\mathrm{}`$ define
$$A_n:D(A_n)\mathrm{}^2\mathrm{}^2,D(A_n)=\{(x_j)_{j1}\mathrm{}^2;\underset{j1}{}j^2|x_j|^2<\mathrm{}\}$$
$$A_n๐_j=\{\begin{array}{cc}j๐_j,\hfill & \hfill jn\\ n๐_j,\hfill & \hfill j=n\end{array}$$
One can see that
$$(๐ข\pm A_n)^1(๐ข\pm A_0)^1=\left|\frac{1}{๐ข+n}\frac{1}{๐ขn}\right|0$$
so that $`\gamma (A_n,A_0)0`$. On the other hand, if $`\alpha ๐`$ is as in Proposition 1.4 then for all sufficiently large $`n`$ we have
$$\alpha (A_n)\alpha (A_0)=1.$$
We now want to present a simple criterion of $`\rho `$-convergence. For any closed densely defined operator we denote by $`(T)`$ its resolvent set.
###### Proposition 1.6.
Suppose $`A๐ฎ`$ such that $`(A)\mathrm{}`$. Suppose $`S_n`$ is a sequence of densely defined symmetric operators satisfying the following conditions.
(a) $`D(A)D(S_n)`$.
(b) There exists a sequence of positive numbers $`c_n0`$ such that
$$S_nuc_n(Au+u),uD(A).$$
Then $`A+S_n๐ฎ`$ for all $`n0`$ and
$$\rho (A+S_n,A)0.$$
Proof Set $`A_n:=A+S_n`$. According to \[4, Thm.IV.2.24\] we have
$$\gamma (A_n,A)0$$
while \[4, Thm. V.4.1\] implies $`A+S_n๐ฎ`$ for all sufficiently large $`n`$. Let $`\beta (A)`$ and consider a small closed interval $`I=[\beta \epsilon ,\beta +\epsilon ]`$ such that $`I(A)`$. Then, using \[4, Thm. VI.5.10\] we deduce that for $`n`$ sufficiently large we have
$$I(A_n),n0.$$
Pick now a function $`\alpha ๐`$ such that $`\alpha (\lambda )1`$ for $`\lambda \beta +\epsilon `$ and $`\alpha (\lambda )0`$ for $`\lambda \beta \epsilon `$. Using \[4, Thm. VI.5.12\] we deduce
$$\alpha (A_n)\alpha (A)0.$$
We can now invoke Proposition 1.4 to conclude that $`\rho (A_n,A)0`$. $`\mathrm{}`$
## 2 Families of boundary value problems
Consider now as in \[6, App. A\] the following data.
$``$ A compact, oriented Riemannian manifold $`(M,g)`$ with boundary $`N=M`$ such that a tubular neighborhood of $`NM`$ is isometric to the cylinder
$$([0,1]\times N,dt^2+g_N)$$
where $`g_N`$ is a Riemann metric on $`N`$ and $`t`$ denotes the outgoing longitudinal coordinate.
$``$ An Euclidean bundle of Clifford modules $`EM`$ with Clifford multiplication
$$๐:T^{}M\mathrm{End}(E).$$
( $`๐(\alpha )`$ is skew-symmetric for any real $`1`$-form $`\alpha `$.) Set $`E_0:=E_N`$
$``$ $`D:C^{\mathrm{}}(E)C^{\mathrm{}}(E)`$ a symmetric Dirac operator with principal symbol $`๐`$ such that near $`N`$ it has the form
$$D=J(_tD_0),J:=๐(dt)$$
where $`D_0:C^{\mathrm{}}(E_0)C^{\mathrm{}}(E_0)`$ is symmetric and independent of $`t`$.
$``$ A sequence of symmetric endomorphisms of $`E`$ independent of $`t`$ near $`N`$ such that
$$T_n_{C^2}0$$
and (near $`N`$) the endomorphism $`JA_n`$ is symmetric. Set $`D_n:=D+T_n`$. Observe that near $`N`$ $`D_n`$ has the form
$$D_n:=J(_tD_0JT_n).$$
Following , we consider the family $`๐ซ`$ of admissible boundary conditions. It consists of zero order, formally selfadjoint pseudodifferential projectors with the same principal symbol as the Calderon projector of $`D_0`$. The symbol of any $`P`$ in $`P`$ commutes with the symbol of $`D_0`$ so that the commutator $`[P,D_0]`$ is a zeroth order pseudodifferential operator. We define a metric $`\nu `$ on $`๐ซ`$ by setting
$$\nu (P,Q):=PQ+[PQ,D_0]$$
where $``$ denotes the norm on the space of bounded operators $`L^2(E_0)L^2(E_0)`$.
Suppose now that we are given a projector $`P๐ซ`$ and a sequence $`(P_n)๐ซ`$. As in , we can form the Fredholm selfadjoint operators
$$A_n:D(A_n)L^2(E)L^2(E),D(A_n)=\{uH^1(E);P_nu_N=0\}$$
$$A_nu=D_nu$$
and
$$A:D(A)L^2(E)L^2(E),D(A)=\{uH^1(E);Pu_N=0\}$$
$$Au=Du.$$
###### Proposition 2.1.
If
$$\underset{n\mathrm{}}{lim}\nu (P_n,P)=0$$
(2.1)
Then
$$\underset{n\mathrm{}}{lim}\rho (A_n,A)=0.$$
Proof The proof relies on the following technical result.
###### Lemma 2.2.
There exists a sequence of bounded, invertible operators $`U_n:L^2(E)L^2(E)`$ such that
(i) $`1U_n`$ and $`1U_n^{}`$ define bounded operators $`H^1(E)H^1(E)`$
(ii) $`(U_n1),(U_n1)^{}0`$ in the norm topology on the space of bounded operators $`H^s(E)H^s(E)`$, $`s=0,1`$.
(iii) $`D(A_n)=U_n^{}D(A)`$, $`n`$.
We will prove this lemma after we have finished the proof of Proposition 2.1. Set
$$B_n:=U_nA_nU_n^{}.$$
Observe that $`B_n๐ฎ`$ and $`D(B_n)=D(A)`$. Moreover
$$\rho (B_n,A_n)=\mathrm{\Psi }(U_nA_nU_n^{})\mathrm{\Psi }(A_n)=U_n\mathrm{\Psi }(A_n)U_n^{}\mathrm{\Psi }(A_n)$$
$$=((U_n1)+1)\mathrm{\Psi }(A_n)((U_n1)+1)^{}\mathrm{\Psi }(A_n)C(U_n1)_{L^2,L^2}\mathrm{\Psi }(A_n)0$$
Thus it suffices to show that
$$\rho (B_n,A)0.$$
Observe that for all $`uD(A)`$ we have
$$B_nuAu=U_n(D+T_n)U_n^{}DU_nD(U_n^{}uu)+U_nT_nU_n^{}u$$
$$U_n_{L^2,L^2}D(U_n^{}uu)_{L^2}+CT_n_{C^2}u_{L^2}C\left((U_n^{}1)u_{H^1}+T_n_{C^2}u_{L^2}\right)$$
$$C\left((U_n^{}1)_{H^1,H^1}u_{H^1}+T_n_{C^2}u_{L^2}\right)$$
(use the elliptic estimates in )
$$C\left\{(U_n^{}1)_{H^1,H^1}(Au_{L^2}+u_{L^2})+T_n_{C^2}u_{L^2}\right\}c_n(Au+u)$$
where $`c_n0`$. Thus, the operator $`S_n=B_nA`$ satisfies all the conditions in Proposition 1.6. On the other hand, $`A`$ has compact resolvent so that $`(A)\mathrm{}`$. We deduce
$$\rho (A,B_n)=\rho (A,A+S_n)0.\mathrm{}$$
Proof of Lemma 2.2 Following the constructions in \[4, I.ยง6.4\] define
$$\widehat{U}_n:L^2(E_0)L^2(E_0),\widehat{U}_n=P_nP+(1P_n)(1P)=2P_nP(P_n+P)+1$$
$$=2(P+R_n)P(2P+R_n)+1=R_n(2P1)+1.$$
$`\widehat{U}_n`$ is a pseudodifferential operator of order zero with principal symbol $`1`$. Observe that
$$\widehat{U}_n^{}=PP_n+(1P)(1P_n)$$
and, as explained in \[4, I.ยง6.4\], $`\widehat{U}_n^{}`$ is invertible and maps $`\mathrm{ker}P`$ onto $`\mathrm{ker}P_n`$. Observe moreover that
$$\widehat{U}_n1_{L^2,L^2}R_n_{L^2,L^2}(2P1)_{L^2,L^2}0.$$
(2.2)
Next, observe that
$$[D_0,\widehat{U}_n]=[D_0,R_n](2P1)+2R_n[D_0,P]$$
defines a bounded operator $`L^2(E_0)L^2(E_0)`$ and, using (2.1) we deduce
$$[D_0,\widehat{U}_n]_{L^2,L^2}0.$$
(2.3)
Observe that $`\widehat{U}_n`$ defines in an obvious fashion a bounded operator
$$\widehat{U}_n:L^2(E_{[0,1]\times N})L^2(E_{[0,1]\times N})$$
Consider now a smooth increasing function
$$\eta :[0,1][0,1]$$
such that $`\eta (t)0`$ for $`t<1/4`$ and $`\eta (t)1`$ for $`t>3/4`$. We can regard $`\eta `$ as a function on the tubular neighborhood of $`NM`$ and then extending it by $`0`$ we can regard it as a smooth function on $`M`$. Notice that if $`u`$ is a section of $`E`$ then we can regard $`\eta u`$ as a section of $`E_{[0,1]\times N}`$.
For any section of $`E`$ smooth up to the boundary define
$$U_nu=(1\eta )u+\widehat{U}_n(\eta u).$$
It is clear that $`U_nu`$ is smooth up to the boundary. Notice also that there exists a constant $`C>0`$ independent of $`n`$ such that
$$U_nu_L^2Cu_{L^2}$$
for any section $`u`$ smooth up to the boundary. Thus $`U_n`$ extends to a bounded operator $`L^2(E)L^2(E)`$. Using (2.2) we deduce that
$$(U_n1)_{L^2,L^2}0.$$
We want to show that $`U_n`$ induces a bounded operator $`H^1(E)H^1(E)`$ and then estimate the norm of $`(U_n1)`$ as a bounded operator $`H^1H^1`$.
First of all observe that the elliptic estimates for $`D_0`$ imply that there exists a positive constant $`C`$ such that if $`u`$ is smooth up to the boundary then
$$C^1u_{H^1([0,1]\times N)}_tu_{L^2([0,1]\times N)}+D_0u_{L^2([0,1]\times N)}Cu_{H^1([0,1]\times N)}$$
Observe that for any section $`u`$ smooth up to the boundary we have
$$U_nuu_{H^1(M)}=(1\eta )u+\widehat{U}_n(\eta u)u_{H^1(M)}$$
$$=\widehat{U}_n(\eta u)\eta u_{H^1(M)}=\widehat{U}_n(\eta u)(\eta u)_{H^1([0,1]\times N)}$$
$$\begin{array}{c}C(\widehat{U}_n(\eta u)(\eta u)_{L^2([0,1]\times N)}+_t\widehat{U}_n(\eta u)_t(\eta u)_{L^2([0,1]\times N)}\\ \\ +D_0\widehat{U}_n(\eta u)D_0(\eta u)_{L^2([0,1]\times N)})\end{array}$$
(2.4)
Using (2.2) we deduce
$$\widehat{U}_n(\eta u)(\eta u)_{L^2([0,1]\times N)}c_nu_{L^2(M)},c_n0.$$
To estimate the second term in (2.4) notice first that $`[_t,\widehat{U}_n]=0`$ so that we have
$$_t\widehat{U}_n(\eta u)_t(\eta u)_{L^2([0,1]\times N)}=\widehat{U}_n_t(\eta u)_t(\eta u)_{L^2([0,1]\times N)}$$
$$c_n_tu_{L^2([0,1]\times N)}c_nu_{H^1(M)},c_n0.$$
The estimate of the third term in (2.4) requires a bit more work. Observe that
$$D_0\widehat{U}_n(\eta u)D_0(\eta u)=[D_0,\widehat{U}_n](\eta u)+\widehat{U_n}(D_0\eta u)D_0(\eta u)$$
$$=\eta \left([D_0,\widehat{U}_n]u+\widehat{U}_n(D_0u)D_0u\right)$$
so that
$$D_0\widehat{U}_n(\eta u)D_0(\eta u)_{L^2([0,1]\times N)}[D_0,\widehat{U}_n]u_{L^2([0,1]\times N)}+\widehat{U}_n(D_0u)D_0u_{L^2([0,1]\times N)}$$
(use (2.2))
$$c_n(u_{L^2([0,1]\times N)}+D_0u_{L^2([0,1]\times N)})c_n^{}u_{H^1(M)},c_n^{}0.$$
We have thus found a sequence of positive numbers $`c_n0`$ such that
$$U_nuu_{H^1(M)}c_nu_{H^1(M)}$$
for every section $`u`$ smooth up to the boundary. This shows that $`U_n`$ induces a bounded operator $`H^1(M)H^1(M)`$ and moreover,
$$U_n1_{H^1,H^1}c_n0.$$
One can prove a similar statement concerning $`U_n^{}`$. Clearly $`U_n`$ is invertible being so close to $`1`$. Since $`\mathrm{ker}P_n=\widehat{U}_n^{}(\mathrm{ker}P)`$ we deduce that $`D(A_n)=U_n^{}D(A)`$. Lemma 2.2 is proved. $`\mathrm{}`$
## 3 Classifying spaces for $`K`$-theory
For clarity purposes we will consider only a special case, that of the functor $`KO^1`$. To discuss the other functors $`KO^n`$ one should use the bigraded Karoubi functors $`KO^{p,q}`$ as we did in . The proof is only notationally more complicate.
Denote by $`๐ฎ`$ (resp. $`๐ฎ`$, $`[][๐ฎ]`$) the subspace of selfadjoint Fredholm operators. $`[]`$ has three connected components. Two of them $`[_\pm ]`$, are contractible while the third, $`[_0]`$ is a classifying space for $`KO^1`$ (see ). We deduce that $`(,\rho )`$ consists of three components
$$_\pm :=\mathrm{\Psi }^1([_\pm ]),_0:=\mathrm{\Psi }^1([_0])$$
and $`(_0,\rho )`$ is a classifying space for $`KO^1`$.
Observe that $`HH`$ is a symplectic space with complex structure
$$J=\left[\begin{array}{cc}0& 1_H\\ 1_H& 0\end{array}\right]$$
and $`\mathrm{\Lambda }_0:=H0`$ is a Lagrangian subspace. Define $`_0`$ the set of Lagrangian subspaces $`\mathrm{\Lambda }HH`$ such that $`(\mathrm{\Lambda }_0,\mathrm{\Lambda })`$ is a Fredholm pair. We topologize $`_0`$ using the gap distance $`\delta `$. The space $`(_0,\delta )`$ is also a classifying space for $`KO^1`$ (see ).
There is a natural $`11`$ map
$$\mathrm{\Gamma }:_0_0,A\mathrm{\Gamma }_A.$$
According to Lemma 1.2 the map $`\mathrm{\Gamma }:(_0,\rho )(_0,\delta )`$ is continuous.
###### Theorem 3.1.
The map
$$\mathrm{\Gamma }:(_0,\rho )(_0,\delta )$$
is a weak homotopy equivalence.
Proof Fix $`A_0_0`$. We have to show that for every $`n>0`$ the induced map
$$\mathrm{\Gamma }_{}:\pi _n(_0,A_0)\pi _n(_0,\mathrm{\Gamma }_{A_0})$$
is an isomorphism. Observe first that, according to Bott periodicity,
$$\pi _n(_0,\mathrm{\Gamma }_{A_0})๐ข:=\{0,,_2\}.$$
The groups in the family $`๐ข`$ have a remarkable property. If $`G๐ข`$ and $`\phi :GG`$ is a surjective morphism then $`\phi `$ is an isomorphism.
In \[6, ยง5.3\], using the symplectic reduction morphism it is shown that the morphism $`\mathrm{\Gamma }_{}`$ is surjective provided the (general) Floer families are $`\rho `$-continuous. This continuity was established in Proposition 2.1. Theorem 3.1 is proved. $`\mathrm{}`$
###### Remark 3.2.
In we claimed that the map $`\mathrm{\Gamma }:(_0,\gamma )(_0,\delta )`$ is a weak homotopy equivalence when in fact the arguments there, detailed in this paper, prove this only for the stronger $`\rho `$-topology. This has no effect on the results of but one รฆsthetical question still lingers. Is the space $`_0`$ equipped with the gap topology a classifying space for $`KO^1`$? If the answer is yes (which we continue to belive to be the case) then our claim in is true.
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# A New Narrow-Line Seyfert 1 galaxy : RX J1236.9+2656
## 1 INTRODUCTION
Narrow-line Seyfert 1 (NLS1) galaxies are considered to be a special class of โnormalโ Seyfert 1 galaxies because of their peculiar properties that distinguish them from the latter class. They are characterized by their optical spectra having permitted lines that are narrower than those found in the normal Seyfert 1 galaxies, e.g., full width at half maximum (FWHM) of H$`\beta `$ line is $`2000\mathrm{km}\mathrm{s}^1`$, relatively weak forbidden lines, $`\frac{[OIII]\lambda 5007}{H\beta }<3`$ (Osterbrock & Pogge 1985), and strong Fe II emission. NLS1 galaxies also have distinctive soft X-ray properties as well. They show steep soft X-ray spectrum with little or no absorption above the Galactic value (Grupe et al. 1998). They often show rapid and large amplitude as well as long-term X-ray variability (Boller et al. 1993; Brandt, Pounds, & Fink 1995; Grupe et al. 1995a,b). In spite of the dominance of soft X-ray emission, soft X-ray luminosity of NLS1 galaxies are similar to those of normal Seyfert 1s. $`ASCA`$ observations show that the hard X-ray ($`210\mathrm{keV}`$) continua of NLS1s are also steeper than those of normal Seyfert 1s with broader H$`\beta `$ FWHM (Brandt, Mathur, & Elvis 1997; Leighly 1999b). NLS1 galaxies also show more variability in hard X-rays than the normal Seyfert 1s (Leighly, 1999a). The spectral energy distribution (SED) from far-infrared (FIR) to X-rays of NLS1 galaxies appears to be similar to that of broad-line Seyfert 1 galaxies. However, the UV luminosity of NLS1 galaxies tends to be smaller than those of Seyfert 1s (Rodriguez-Pascual, Mas-Hesse, & Santos-Lleรณ 1997).
Optical spectroscopy of Ultra-soft X-ray sources discovered with $`Einstein`$, and $`ROSAT`$ has been an efficient way to identify NLS1 galaxies (e.g, Puchnarewicz et al. 1992; Grupe et al. 1998). As part of our programme to optically identify and study in detail the counterparts of the ultra-soft sources in the catalogue of Singh et al. (1995), we have discovered a NLS1 galaxy RX J1236.9+2656. The basic parameters of RX J1236.9+2656 are given in Table 1.
Throughout this note, luminosities are calculated assuming an isotropic emission, a Hubble constant of $`H_0=75\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ and a deceleration parameter of $`q_0=0`$ unless otherwise specified.
## 2 X-ray Observation, Analysis & Results
The region of the sky containing the source, RX J1236.9+2656, was observed seven times with the $`ROSAT`$ (Truemper et al. 1983) Position Sensitive Proportional Counter (PSPC) during 1991โ1993 and twice with the High Resolution Imager (HRI) (Pfeffermann et al. 1987) in 1996 JuneโJuly. The exposure times were in the range $`1420\mathrm{s}5422\mathrm{s}`$ for the PSPC observations while the two HRI observations were carried out with longer exposure times ($`14545\mathrm{s}`$ and $`16738\mathrm{s}`$). The offset of the source from the field center was $`17\mathrm{}`$ for each of the PSPC and HRI observations.
The X-ray source, RX J1236.9+2656, was identified by overlaying the contours of high resolution X-ray images obtained from $`ROSAT`$ HRI observations onto optical images obtained from the Digital Sky Survey (DSS). No other X-ray source, within the angular spread comparable to the point spread function of $`ROSAT`$ HRI, was seen in the overlays. Therefore, X-ray emission from RX J1236.9+2656 is not contaminated by emission from any other source. HRI count rates for RX J1236.9+2656 were obtained using a circle of radius $`50\mathrm{}`$ for the source and an annulus of inner circle radius $`60\mathrm{}`$ and width $`60\mathrm{}`$ for background. The HRI count rates thus estimated are ($`1.14\pm 0.14)\times 10^2`$ and ($`1.23\pm 0.13)\times 10^2\mathrm{count}\mathrm{s}^1`$ for the two observations. The spatial resolution with the PSPC at an offset of $`16.5\mathrm{}`$ is $`40\mathrm{}`$ (half power radius) (Hasinger et al. 1993). Therefore, PSPC count rates for RX J1236.9+2656 were obtained using a circle of radius $`2.5\mathrm{}`$ for the source and 7 nearby circular regions of radii $`2.25\mathrm{}`$ for background. The PSPC count rates are ($`3.53\pm 0.46)\times 10^2`$, ($`3.19\pm 0.57)\times 10^2`$, ($`3.19\pm 0.60)\times 10^2`$, ($`3.54\pm 0.57)\times 10^2`$, ($`2.23\pm 0.87)\times 10^2`$, ($`3.26\pm 0.34)\times 10^2`$, and ($`3.15\pm 0.40)\times 10^2\mathrm{count}\mathrm{s}^1`$ for the 7 PSPC observations.
In order to investigate the time variability of soft X-ray emission from RX J1236.9+2656, we extracted light curves from the $`ROSAT`$ PSPC observations using the same source and the background regions as described above and in the PSPC energy band of 0.1โ2.4 keV containing all the X-ray photons. The background subtractions were carried out after appropriately scaling the background light curves to have the same area as the source extraction area. The light curves of RX J1236.9+2656 do not show short-term variability during the PSPC observations. However, on a longer time scale of months to years, variability is clearly detected in the X-ray flux measurements plotted in Figure 2. The observed flux in the energy band $`0.12.0\mathrm{keV}`$, estimated from the best fit spectral model (see below), increased by about a factor of 2 within $`3\mathrm{yr}`$.
For analyzing the X-ray spectra of RX J1236.9+2656, we choose 3 PSPC spectra corresponding to those observations for which the exposure times were greater than $`3000\mathrm{s}`$. These observations were carried out on 1991 December 15, 1992 June 30, and 1993 June 17. Photon energy spectra of RX J1236.9+2656 were accumulated from their PSPC observations using the same source and background regions as stated above. $`ROSAT`$ PSPC pulse height data were appropriately re-grouped to improve the statistics.
We used the XSPEC spectral analysis package to fit the data with spectral models. Appropriate ancillary response file was used to account for the off-axis position of the source. An appropriate response matrix was used to define the energy response of the PSPC.
Each of the 3 PSPC spectra was first fitted with a redshifted power-law model absorbed by an intervening medium with absorption cross-sections as given by Balucinska-Church and McCammon (1992) and using the method of $`\chi ^2`$-minimization. The photon index ($`\mathrm{\Gamma }_X`$) was found to be very steep in all the cases. It was found that the absorbing column density in each case is similar within errors to the Galactic value ($`\mathrm{N}_\mathrm{H}\text{ }=1.33\times 10^{20}\mathrm{cm}^2`$) measured from 21-cm radio observations (Dickey & Lockman 1990) along the direction of the source, indicating that all the X-ray absorption is due to matter in our own Galaxy. Therefore, we have fitted the power-law models to these spectra after fixing the neutral hydrogen column density to the Galactic value. The best-fit minimum $`\chi _\nu ^2`$ and the power-law index do not change significantly from those obtained while varying the $`\mathrm{N}_\mathrm{H}`$ . The photon indices obtained for fixed $`\mathrm{N}_\mathrm{H}`$ are, however, better constrained. The values for $`\mathrm{\Gamma }_X`$ are $`3.0_{0.4}^{+0.5}`$, $`4.2_{0.7}^{+1.2}`$, and $`4.2_{0.7}^{+2.0}`$ for the spectra observed on 1991 Dec 15, 1992 June 30, and 1993 June 17, respectively, and are quite similar within the errors for all three spectra. The errors quoted, here and below, were calculated at the $`90\%`$ confidence level based on $`\chi _{\mathrm{min}}^2`$+2.71. In order to better constrain the model parameters, we have fitted the above model to the three spectra jointly. The best-fit photon index is now $`4.0_{0.5}^{+0.5}`$. The observed flux, based on the best-fit model parameters, is estimated to be $`1.4\times 10^{13}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ in the energy band of $`0.12.0\mathrm{keV}`$. In the same energy band the intrinsic soft X-ray luminosity, corrected for Galactic absorption, of RX J1236.9+2656 is calculated to be $`1.6\times 10^{43}\mathrm{erg}\mathrm{s}^1`$.
We have also analyzed all the 7 PSPC spectra jointly. The best-fit phton index is $`3.7_{0.5}^{+0.3}`$ which is similar to that obtained for the 3 PSPC observations above. Thus, it is clear that all available PSPC spectral data of RX J1236.9+2656 are well represented by a power-law of photon index $`3.7_{0.5}^{+0.3}`$ absorbed by the matter in our own Galaxy. We have also fitted redshifted blackbody models, absorbed by an intervening medium, to the three PSPC spectra (exposure times $`>`$ 3000 s) of RX J1236.9+2656 taken jointly as well as all the 7 data taken jointly. The absorbing column density derived from the model fit was lower than the Galactic value in that direction, indicating that the blackbody is not a suitable model. We then fixed the absorbing column to the Galactic value and carried out the joint fitting. The temperature thus obtained, $`kT=79_{24}^{+13}\mathrm{eV}`$, reflects the ultra-soft nature of RX J1236.9+2656.
We have calculated the $`ROSAT`$ HRI flux of RX J1236.9+2656 using the best-fit model parameters obtained from best-fit to all the 7 spectral data ($`\mathrm{\Gamma }_X=3.7`$, $`\mathrm{N}_\mathrm{H}\text{ }=1.33\times 10^{19}\mathrm{cm}^2`$). The observed HRI fluxes are $`2.7\times 10^{13}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^2`$ and $`2.9\times 10^{13}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^2`$ in the energy band of $`0.12.0\mathrm{keV}`$ for the two HRI observations. The HRI fluxes are about a factor of two higher than the values obtained from the PSPC observations.
## 3 Optical Spectroscopy
Low resolution optical spectroscopic observations of RX J1236.9+2656 were carried out at the 3.5m telescope of the Apache Point Observatory (APO) on the nights of 1995 April 25, 30 and 2000 January 6. The integrations ranged from 4-15 min, with the best spectrum obtained during the last observation. The Double Imaging Spectrograph (DIS) was used in low resolution mode with a $`1.5\mathrm{}`$ slit to obtain simultaneous blue and red spectra covering $`38005300\mathrm{\AA }`$ in the blue and $`56009900\mathrm{\AA }`$ in the red with a resolution of 14ร
.
The optical spectra were reduced using routines available within the IRAF <sup>1</sup><sup>1</sup>1IRAF (Image Reduction and Analysis Facility) is distributed by the National Optical Astronomy Observatories, which are operated by AURA, Inc., under cooperative agreement with the National Science Foundation. software. This included correcting the images using bias and flat fields, extracting sky-subtracted one-dimensional spectra and calibrating the wavelengths from HeNeAr lamps and the fluxes from standard stars. In the spectrum of RX J1236.9+2656, strong emission lines of Balmer H$`\alpha `$, H$`\beta `$, and the forbidden line \[O III\]$`\lambda 5007`$ are readily observed. Using the peak wavelengths and the rest wavelengths of these lines, a redshift of 0.225$`\pm 0.001`$ was derived from the 2000 January spectrum. The observed spectrum was then corrected for this redshift and is shown in Figure 2. The signal-to-noise ratio of the spectrum is $`7.6`$ measured from the dispersion in the continuum region $`6900\mathrm{\AA }7200\mathrm{\AA }`$. The continuum is observed to rise towards the blue end of the spectrum. Apart from the emission lines mentioned above, we have also identified Fe II emission between $`50705600\mathrm{\AA }`$ and a forbidden line of \[Fe VII\]$`\lambda 6087`$ (see Fig. 2). We have fitted Gaussian profiles to the strong emission lines using the profile fitting feature in the โsplotโ task within IRAF. Due to the poor resolution of the spectrum, it was not possible to deblend the H$`\alpha `$, and \[N II\]$`\lambda \lambda 6548,6584`$ lines. A single Gaussian profile is not well fitted to the core of the H$`\alpha `$ line and the FWHM is overestimated because of the presence of \[N II\]$`\lambda \lambda 6548,6584`$ lines. We also fitted a single Gaussian to the H$`\alpha `$ after excluding the wings. In this case, the FWHM can be considered as a lower limit. Thus we derive the FWHM of the H$`\alpha `$ line to be in the range of $`13151692\mathrm{km}\mathrm{s}^1`$. The widths have been corrected for instrumental broadening by subtracting, in quadrature, the instrumental broadening (FWHM = $`14\mathrm{\AA }`$) from the observed FWHM. The H$`\alpha `$+\[N II\] flux is estimated to be $`8.3\times 10^{15}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ by integrating the flux over the H$`\alpha `$+\[N II\] profile. The profile of the Balmer line H$`\beta `$ is well fitted by a Lorentzian profile but poorly fitted with a Gaussian. The FWHMs of the best-fit Lorentzian and Gaussian profiles to the H$`\beta `$ line are $`1122\mathrm{km}\mathrm{s}^1`$ and $`1392\mathrm{km}\mathrm{s}^1`$, respectively. The observed widths of H$`\alpha `$ and H$`\beta `$, and the presence of Fe II, \[Fe VII\] emission are characteristic of NLS1 galaxies. The flux of the H$`\beta `$ line is estimated to be $`2.0\times 10^{15}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$. Similarly the flux of the \[O III\]$`\lambda 5007`$ line is calculated to be $`7.6\times 10^{16}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$. Thus the ratio, $`\frac{[OIII]\lambda 5007}{H\beta }`$ is only $`0.38`$.
## 4 Discussion
RX J1236.9+2656 is luminous in soft X-rays, with a rest frame intrinsic luminosity $`1.6\times 10^{43}\mathrm{erg}\mathrm{s}^1`$, in the energy band of 0.1โ2.0 keV. Seyfert 1 galaxies studied by Rush et al. (1996) span over 4 orders of magnitude in soft X-ray luminosity, from below $`10^{42}\mathrm{erg}\mathrm{s}^1`$ to above $`10^{46}\mathrm{erg}\mathrm{s}^1`$ in the energy band of 0.1โ2.4 keV. Thus, the X-ray luminosity of RX J1236.9+2656 is similar to that of a Seyfert 1 galaxy. Assuming that the soft X-ray luminosity of RX J1236.9+2656 is about $`10\%`$ of the bolometric luminosity, the lower limit to the mass of the central supermassive object or the Eddington mass is $`10^6\mathrm{M}\mathrm{}`$.
The galaxy RX J1236.9+2656 shows long-term variability โ a change in intensity by a factor of $`2`$ within $`3\mathrm{yr}`$, another characteristic of NLS1 galaxies. Short-term ($`1000100000\mathrm{s}`$) variability is not detected from RX J1236.9+2656 due to poor signal-to-noise ratio of the $`ROSAT`$ data.
The soft X-ray spectrum of RX J1236.9+2656 is steeper ($`\mathrm{\Gamma }_X3.7`$) than those of normal Seyfert 1s \[$`<\mathrm{\Gamma }_X>`$ ($`90\%`$ range)$`=2.02.7`$\], and similar to those of NLS1 galaxies \[$`\mathrm{\Gamma }_X(90\%\mathrm{range})=2.33.7`$\] (Grupe et al. 1998). Lack of intrinsic soft X-ray absorption over the Galactic value in RX J1236.9+2656 is similar to the results found in normal Seyfert 1s and NLS1 galaxies. The steeper power-law index and blackbody model fit to the PSPC spectra of RX J1236.9+2656 indicate an ultra-soft nature of this object. The derived temperature of the blackbody, $`kT75\mathrm{eV}`$, is similar to those found in NLS1 galaxies (Brandt & Boller 1998).
The optical spectrum of RX J1236.9+2656 appears to be typical of NLS1 galaxies. The FWHM of the H$`\beta `$ line (in the range of 1122โ1392 km s<sup>-1</sup>) is narrower than those found in normal Seyfert 1s and is similar to those found in NLS1 galaxies (FWHM$`{}_{H\beta }{}^{}2000\mathrm{km}\mathrm{s}^1`$). The ratio $`\frac{[OIII]\lambda 5007}{H\beta }`$ ($`0.38`$) for RX J1236.9+2656 indicates that forbidden lines are weak, similar to that observed from NLS1 galaxies ($`\frac{[OIII]\lambda 5007}{H\beta }<3.0`$). An Fe II multiplet between $`50705600\mathrm{\AA }`$ is also detected from RX J1236.9+2656. At blue wavelengths, there is an indication of the presence of an Fe II multiplet between $`44354700\mathrm{\AA }`$ although our observation did not cover the multiplet fully. Thus the optical emission line parameters strongly suggest that RX J1236.9+2656 is a NLS1 galaxy. Furthermore, the position of RX J1236.9+2656 on the $`\mathrm{\Gamma }_X`$โH$`\beta `$ line width plane (Fig. 8 of Boller et al. 1996) is consistent with other NLS1 galaxies.
## 5 Conclusions
A narrow-line Seyfert 1 galaxy, RX J1236.9+2656, has been discovered based on the following soft X-ray and optical emission line properties: (i) Steep soft X-ray spectrum ($`\mathrm{\Gamma }3.7`$), high soft X-ray luminosity ($`1.6\times 10^{43}\mathrm{erg}\mathrm{s}^1`$), the lack of intrinsic soft X-ray absorption, and X-ray variability. (ii) Narrow Balmer lines (FWHM$`<2000\mathrm{km}\mathrm{s}^1`$), weak \[O III\]$`\lambda 5007`$ emission, and presence of Fe II multiplets.
## 6 Acknowledgments
This research has made use of $`ROSAT`$ archival data obtained through the High Energy Astrophysics Science Archive Research Center, HEASARC, Online Service, provided by the NASA-Goddard Space Flight Center. P.S and D.W. Hoard acknowledge support from NASA LTSA grant NAG 53345. We thank an anonymous referee for his$`/`$her suggestions to reduce the size of this research note.
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# Type Ia Supernovae and their implications for cosmology
## 1 Introduction
During the past three years two groups (Perlmutter et al. 1997; Schmidt et al. 1998) have presented strong evidence that the expansion of the universe is accelerating rather than decelerating (Riess et al. 1998; Perlmutter et al. 1998, 1999; and see Livio 1999 for a perspective). This surprising result comes from distance measurements to more than fifty supernovae Type Ia in the redshift range $`z=0.1`$ to $`z=1`$. The results are consistent with the cosmological constant (or vacuum energy) contributing to the total energy density about 60โ70% of the critical density, which in turn, is consistent with recent measurements of the anisotropy of the cosmic microwave background (e.g. Miller et al. 1999; Wilson et al. 1999; Mauskopf et al. 1999).
This unexpected finding, as well as the use of supernovae Type Ia to measure the Hubble constant (e.g. Sandage et al. 1996; Saha et al. 1997), have focused the attention again on the frustrating fact that in spite of decades of research, the exact nature of the progenitors of supernovae Type Ia remains unknown. Until this problem is solved, one cannot be fully confident that supernovae at higher redshifts are not somehow different from their low redshift counterparts. In the present review I therefore examine critically models for supernovae Type Ia and their progenitors. Other recent reviews include Branch et al. (1995), Livio (1996a; 2000), Renzini (1996), Iben (1997), and see Hรถflich & Dominguez, these procedings.
## 2 SNe Ia characteristics and the basic model
The defining characteristics of supernovae Type Ia (SNe Ia) are both spectral: (i) the lack of lines of hydrogen, and (ii) the presence of a strong red Si II absorption feature ($`\lambda `$6355 shifted to $`6100`$ ร
).
Once defined as SNe Ia, the following are several of the important observational characteristics of the class which may help in the search for progenitors:
1. Homogeneity: Until very recently, it has generally been claimed that more than 80% of all SNe Ia form a homogeneous class (see however (2) below) in terms of their spectra (e.g. Branch, Fisher, & Nugent 1993), light curves, and peak absolute magnitudes. The latter are given by
$$M_\mathrm{B}M_\mathrm{V}19.30(\pm 0.03)+5\mathrm{log}(H_0/60\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1)$$
(2.1)
with a dispersion of $`\sigma (M_\mathrm{B})\sigma (M_\mathrm{V})0.2`$โ0.3 (Hamuy et al. 1996a; Tamman & Sandage 1995; and see Branch 1998 for a review).
2. Inhomogeneity: Some differences in the spectra and light curves have been known to exist for a while (e.g. Hamuy et al. 1996b). In terms of explosion strength, SNe Ia have traditionally been roughly ordered as follows: SNe Ia like SN 1991bg and SN 1992K represent the weakest events, followed by weak events like 1986G, followed by about 80% of all SNe Ia which are called โnormalsโ (or sometimes โBranch normalsโ), to the stronger than normal events like SN 1991T. In a very recent work, however, Li et al. (2000) find indications for a considerably higher peculiarity rate, a total of ($`39\pm 10`$)%; of which ($`19\pm 7`$)% and ($`21\pm 7`$)% are SN 1991bg-like and SN 1991T-like objects respectively.
3. The luminosity function of SNe Ia was found in earlier studies to decline very steeply on the bright side (e.g. Vaughan et al. 1995). Since selection effects cannot prevent the discovery of SNe which are brighter than the โnormalsโ (unless they occur preferentially in high extinction regions), this is usually taken to imply that the normals are essentially the brightest. The recent study of Li et al. (2000) seems to show, however, that the luminosity function is relatively flat at both the overluminous and underluminous ends.
4. Near maximum light, the spectra are characterized by high velocity (8000โ30,000 km s<sup>-1</sup>) intermediate mass elements (OโCa). In the late, nebular phase, the spectra are dominated by forbidden lines of iron (e.g. Kirshner et al. 1993; Wheeler et al. 1995; Ruiz-Lapuente et al. 1995; Gรณmez et al. 1996; Filippenko 1997).
5. Fairly young populations appear to be very efficient at producing SNe Ia (e.g. they tend to be associated with spiral arms in spirals; Della Valle & Livio 1994; Bartunov, Tsvetkov & Filimonova 1994), but relatively old populations ($`\tau \stackrel{>}{}4\times 10^9`$ yr) can also produce them. In particular, SNe Ia do occur in ellipticals (e.g. Turatto, Cappellaro & Benetti 1994). In fact, the rates of SNe Ia in ellipticals appear to be similar to those in spirals, $`0.18SNu`$ (where $`SNu=1SN(100`$ yr)$`{}_{}{}^{1}\left(10^{10}L_{}^B\right)_{}^{1}`$; Turatto, Cappellaro & Petrosian 1999). This immediately implies that SNe Ia are not caused by the core collapse of stars more massive than 8 M.
6. There exist a number of correlations between different pairs of observables (see e.g. Branch 1998 for a review). Of these, the most frequently used in the context of determinations of cosmological parameters is the correlation between the absolute magnitude and the shape of the light curve. Basically, brighter SNe Ia decline more slowly. A parameter commonly used to quantify the light curve shape is $`\mathrm{\Delta }m_{15}`$ (Phillips 1993), the decline in magnitudes in the $`B`$ band during the first 15 days after maximum light. Hamuy et al. (1996a) find slopes $`dM_\mathrm{B}/d\mathrm{\Delta }m_{15}=0.78\pm 0.17`$, $`dM_\mathrm{V}/d\mathrm{\Delta }m_{15}=0.71\pm 0.14`$, and $`dM_\mathrm{I}/d\mathrm{\Delta }m_{15}=0.58\pm 0.13`$. Using a stretch-factor $`s`$ (Perlmutter et al. 1997), one can write $`M_\mathrm{B}=M_\mathrm{B}(s=1)\alpha (s1)`$, with $`M_\mathrm{B}(s=1)=19.46`$ (e.g. Sandage et al. 1996), and $`\alpha =1.74`$ (Perlmutter et al. 1999). Sophisticated techniques for using the different correlations in distance determinations have been developed (e.g. Riess et al. 1996, 1998).
The above characteristics can be augmented by the following suggestive facts:
1. The energy per unit mass, $`1/2(10^4`$ km s$`{}_{}{}^{1})^2`$, is of the order of the one obtained from the conversion of carbon and oxygen to iron.
2. The fact that the event is explosive suggests that degeneracy may play a role.
3. The spectrum appears to contain no hydrogen.
4. The explosions can occur with long delays, after the cessation of star formation.
All the properties above have led to one agreed upon model: SNe Ia represent thermonuclear disruptions of mass accreting white dwarfs.
It is interesting that there exists a unanimous consensus on this model in spite of the fact that the essence of flame physics, burning front propagation, and the details of the (presumed) transition from deflagration to detonation (in particular the density at which the transition occurs), which are at the heart of the model, remain as major unsolved problems (e.g. Khokhlov, Oran & Wheeler 1997; Woosley 1997; Reinecke, Hillebrandt & Niemeyer 1998; and see Hรถflich & Dominguez, and Khokhlov, these proceedings). In fact, given these uncertainties, it is almost difficult to understand how the entire family of SNe Ia light curves can be fitted essentially with one parameter (e.g. Perlmutter et al. 1997), although it is possible that all SNe Ia explode at the same WD mass (see ยง4), and that the entire observed diversity stems from different <sup>56</sup>Ni masses.
## 3 Why is identifying the progenitors important?
The fact that we do not know yet what are the progenitor systems of some of the most dramatic explosions in the universe has become a major embarrassment and one of the key unsolved problems in stellar and binary star evolution. There are several important reasons why identifying the progenitors has become more crucial than ever:
1. The use of SNe Ia as one of the main ways to determine key cosmological parameters like $`H_0`$, and the contributions to the energy density (by matter and by the cosmological constant) $`\mathrm{\Omega }_\mathrm{M}`$, $`\mathrm{\Omega }_\mathrm{\Lambda }`$ requires an understanding of the evolution of the luminosity, and the SN rate with cosmic epoch. Both of these depend directly on the nature of the progenitors.
2. Galaxy evolution depends on the radiative, kinetic energy, and nucleosynthetic output of SNe Ia (e.g. Kauffmann, White & Guiderdoni 1993).
3. Due to the uncertainties that still exist in the explosion mechanism itself, a knowledge of the initial conditions and of the distribution of matter in the environment of the exploding star are essential for the understanding of the explosion.
4. An unambiguous identification of the progenitors, coupled with observationally determined SNe Ia rates can help to place meaningful constraints on the theory of binary star evolution (e.g. Livio 1996b; Li & van den Heuvel 1997; Yungelson & Livio 1998; Hachisu, Kato & Nomoto 1999). In particular, a semi-empirical determination of the elusive common-envelope-ejection efficiency parameter, $`\alpha _{\mathrm{CE}}`$, may be possible (e.g. Iben & Livio 1993).
## 4 Refinements to the basic model
The basic model for SNe Ia (that essentially all researchers in the field agree upon) is that of a thermonuclear disruption of an accreting white dwarf (WD). However, additional refinements to the model are possible on the basis of existing observational data and theoretical models. These refinements still do not involve the question of the progenitor systems. Rather, they address the question of the WD composition, and of its mass at the instant of explosion.
### 4.1 The composition of the exploding WD
In principle, the WD that accretes to the point of explosion could be composed of He, of CโO, or of OโNe. Let us examine these possibilities one by one.
1. He WDs: Helium WDs have typical masses that are smaller than $`0.45`$ M (e.g. Iben & Tutukov 1985). While if accreting, these He WDs can explode following central He ignition at $`0.7`$ M, the composition of the ejected matter in this case will be that of He, <sup>56</sup>Ni and decay products (e.g. Nomoto & Sugimoto 1977; Woosley, Taam & Weaver 1986). This is entirely inconsistent with observations (observational characteristic (4) in ยง2). Therefore, He WDs certainly do not produce the bulk of SNe Ia.
2. OโNe WDs: OxygenโNeon WDs form in binaries from main sequence stars of $`10`$ M, although the precise range which allows formation is somewhat uncertain (e.g. Iben & Tutukov 1985; Canal, Isern & Labay 1990; Dominguez, Tornambรฉ & Isern 1993). These systems are probably not numerous enough to constitute the main channel of SNe Ia (e.g. Livio & Truran 1992; Livio 1993). It is also generally expected that OโNe WDs that manage to accrete enough material to reach the Chandrasekhar limit will produce (via electron capture) preferentially accretion-induced collapses (to form neutron stars) rather than SNe Ia (e.g. Nomoto & Kondo 1991; Gutierrez et al. 1996). Accretion induced collapses do not eject enough nickel to match the light curves of normal SNe Ia, although they may be able to explain very subluminous events like SN 1991bg (e.g. Fryer et al. 1999). I should note that the existing calculations have been performed for WDs of OโNeโMg composition, while some recent calculations of the evolution of a 10 M star produce degenerate cores which are almost devoid of magnesium (Ritossa, Garcia-Berro & Iben 1996). Nevertheless, because of the above two points it is unlikely that OโNe WDs produce the bulk of SNe Ia.
3. CโO WDs: CarbonโOxygen WDs are formed in binaries from main sequence stars of up to $`10`$ M. They are therefore both relatively numerous, and they provide a significant โphase space volumeโ (masses in the range 0.8โ1.2 M; accretion rates in the range $`10^8`$$`10^6`$ M/yr) in which they are expected to produce SNe Ia (upon reaching the Chandrasekhar limit; e.g. Nomoto & Kondo 1991). Consequently, the accreting WDs that produce most of the SNe Ia are very probably of CโO composition!
### 4.2 At what mass does the WD explode and where and in what fuel does the ignition take place?
While there is virtually unanimous agreement about everything I said up to now, namely, that: SNe Ia are thermonuclear disruptions of accreting CโO WDs, the next step in the refinement to the model is more controversial. Two major classes of models have been considered, and they suggest entirely different answers to the questions posed by the title of this subsection. In one class, the WD explodes upon reaching the Chandresekhar mass, as carbon ignites at its center. In the second, the WD explodes at a sub-Chandresekhar mass, as helium ignites off-center. I will now review briefly each of these classes and point out their strengths and weaknesses.
#### 4.2.1 Chandrasekhar mass carbon ignitors
In this model, considered โstandard,โ the WD accretes until it approaches the Chandrasekhar mass. Carbon ignition (triggered by compressional heating) occurs at or very near the center and the burning front propagates outwards. Three types of flame propagation models have been considered in the past three decades: (i) detonation (e.g. Arnett 1969; Hansen & Wheeler 1969), (ii) deflagration (e.g. Nomoto, Sugimoto & Neo 1976) and iii) delayed detonation, in which the flame starts as a deflagration which transitions into a detonation at some transition density (e.g. Khokhlov 1991; Woosley & Weaver 1994). Models of the latter two types ((iii) in particular) have generally been quite successful in explaining the observations (see e.g. Hรถflich & Dominguez, these proceedings). The main strengths of this model (central carbon ignition at the Chandrasekhar mass) are (see e.g. Hรถflich & Khokhlov 1996; Nugent et al. 1997; Hรถflich & Dominguez these proceedings, for detailed modeling):
1. Some $`10^{51}`$ ergs of kinetic energy are deposited into the ejecta by nuclear energy.
2. <sup>56</sup>Ni decay powers the lightcurve.
3. The density and composition as a function of the ejection of velocity (X<sub>i</sub>(V<sub>ej</sub>)) are consistent with the observed spectra.
4. The fact that the explosion occurs at the Chandrasekhar mass may explain the broad-brush homogeneity.
5. Spectra (e.g. of SNe 1994D, 1992A) can be fitted in great detail by theoretical models (e.g. Nugent et al. 1997).
The main weaknesses of the Chandrasekhar mass models are:
1. It has proven more difficult than originally thought for WDs to accrete up to the Chandrasekhar mass in sufficient numbers to account for the SNe Ia rate. The difficulty is associated with mass loss episodes in nova explosions, in helium shell flashes and in massive winds or common envelope phases. I will return to some of these problems when I discuss specific progenitor models.
2. For initial WD masses larger than $`1.2`$ M, (which can more easily, in principle, reach the Chandrasekhar mass) accretion-induced collapse is a more likely outcome than a SN Ia (e.g. Nomoto & Kondo 1991).
3. The late-time spectrum ($`300`$ days), and in particular the Fe III feature at $`4700`$ ร
does not agree well with Chandrasekhar mass models (Liu, Jeffrey & Schultz 1998).
4. The โstandardโ model has some difficulty in reproducing the observed (e.g. Riess et al. 1999a) $`\stackrel{>}{}20`$ days rise times.
My overall assessment of Chandrasekhar mass models is that the strengths significantly outweigh the weaknesses. The calculations of late-time, nebular spectra involve many uncertainties, and hence I do not regard weakness (3) above as fatal (although clearly more work will be required to explain it away). Both weaknesses (1) and (2) can be overcome if it can be demonstrated that SNe Ia statistics can be reproduced within the uncertainties that still plague the theoretical population synthesis models. As I will show in ยง5, this appears indeed to be the case. Weakness (4) can be overcome (in principle at least) by lower values of the C/O ratio, or by the presence of a 0.2โ0.4 M envelope (see e.g. Hรถflich, Wheeler & Thielemann 1998). This suggests to me that this is not a fundamental difficulty for the model.
#### 4.2.2 Sub-Chandrasekhar mass helium ignitors
In these models a CโO WD accumulates a helium layer of $`0.15`$ M while the total mass is sub-Chandrasekhar. The helium ignites off-center (at the bottom of the layer), resulting in an event known as โIndirect Double Detonationโ (IDD) or โEdge Lit Detonationโ (ELD). Basically, one detonation propagates outward (through the helium), while an inward propagating pressure wave compresses the CโO core which ignites off-center, followed by an outward detonation (e.g. Livne 1990; Livne & Glasner 1991; Woosley & Weaver 1994; Livne & Arnett 1995; Hรถflich & Khokhlov 1996; and Ruiz-Lapuente, talk presented at the Chicago meeting on Type Ia Supernovae: Theory and Cosmology, October 1998).
The main strengths of ELD (sub-Chandrasekhar) models are:
1. It is easier to achieve the required statistics, since less mass needs to be accreted, and the WD does not need to be extremely massive (e.g. Ruiz-Lapuente, Canal & Burkert 1997; Di Stefano et al. 1997; Yungelson & Livio 1998).
2. The late-time spectrum (in particular the Fe III feature at $`4700`$ ร
) agrees better with ELD models.
3. SNe Ia light curves can be reproduced adequately by ELD models (although the light curves rise somewhat faster than observed, due to <sup>56</sup>Ni heating; Hรถflich et al. 1997).
The main weaknesses of ELD models are:
1. The spectra that are produced by ELD models generally do not agree with observations (e.g. of SN 1994D; Nugent et al. 1997). In particular, the spectra are very blue (due to heating by radioactive Ni), and are dominated by Ni lines, while not showing a strong Si line. The agreement is somewhat better for the subluminous SNe Ia (e.g. SN 1991bg; Nugent et al. 1997; Ruiz-Lapuente, talk presented at the Chicago meeting on Supernovae, October 1998), but even there it is not very good.
2. The highest velocity ejecta have the wrong composition (<sup>56</sup>Ni and He moving at 11,000 to 14,000 km s<sup>-1</sup>, not intermediate mass elements; also no high velocity C; e.g. Livne & Arnett 1995). This is due to the fact that in these models, essentially by construction, the intermediate mass elements are sandwiched by Ni and He/Ni rich layers, at the inner and outer sides, respectively.
3. Since ELD models allow for a range of WD masses, and since more massive WDs produce brighter SNe, one might expect this model to produce a more gradual decline on the bright side of the luminosity function. While this is in contradiction to the observed sharp decline obtained for some of the earlier samples, it may not be in contradiction with the more recent observations of a relatively flat luminosity function (see ยง2 characteristic (3)).
My overall assessment of the sub-Chandrasekhar mass model is that the weaknesses (and in particular weaknesses (1) and (2) which appear almost inevitable) greatly outweigh the strengths in terms of this being a model for the bulk of SNe Ia. It is still possible that ELDs may correctly represent some subluminous SNe Ia (e.g. Ruiz-Lapuente, Canal, & Burkert 1997; Pinto, private communication).
### 4.3 The favored model
On the basis of the above discussion the basic model can now be further refined, and I tentatively conclude that: Most SNe Ia represent thermonuclear disruptions of mass accreting CโO white dwarfs, when these white dwarfs reach the Chandrasekhar limit and ignite carbon at their centers!
## 5 The two possible scenarios
The next step, in which we search for the progenitor systems of SNe Ia is even more controversial. Two possible scenarios have been proposed: (i) The double-degenerate scenario, in which two CO WDs in a binary system are brought together by the emission of gravitational radiation and coalesce (Webbink 1984; Iben & Tutukov 1984). (ii) The single-degenerate scenario, in which a CO WD accretes hydrogen-rich or helium-rich material from a non-degenerate companion (Whelan & Iben 1973; Nomoto 1982).
In the first scenario the progenitor systems are necessarily binary WD systems in which the total mass exceeds the Chandrasekhar mass, and which have binary periods shorter than about thirteen hours (to allow merger within a Hubble time).
In the second scenario the progenitors could be systems like: (i) Recurrent novae (both of the type in which the WD accretes hydrogen from a giant like T CrB, RS Oph, and of the type in which the WD accretes helium rich material from a subgiant like U Sco, V394 CrA, and Nova LMC 1990#2), (ii) Symbiotic Systems (in which the WD accretes hydrogen-rich material from a low mass red giant or a Mira variable), or (iii) persistent Supersoft X-ray Sources (in which the WD accretes at a high rate $`\stackrel{>}{}10^7`$ M/yr from a subgiant companion).
I will now examine the strengths and weaknesses of each one of these scenarios.
### 5.1 The double-degenerate scenario
There is no question that close binary white dwarf systems in which the total mass exceeds the Chandrasekhar mass are an expected outcome of binary star evolution (e.g. Iben & Tutukov 1984; Iben & Livio 1993). Once the lighter WD (which has a larger radius) fills its Roche lobe, it is entirely dissipated within a few orbital periods, to form a massive disk around the primary (e.g. Rasio & Shapiro 1994; Benz, Thielemann & Hills 1989). The subsequent evolution of the system depends largely on the accretion rate through this disk (e.g. Mochkovitch & Livio 1990; see discussion below).
The main strengths of this scenario are the following:
1. The absence of hydrogen in the spectrum is naturally explained in a model which involves the merger of two CโO WDs. In fact, if hydrogen is ever detected in the spectrum of a SN Ia, this would deal a fatal blow to this model. Tentative evidence for circumstellar H$`\alpha `$ absorption is SN 1990M was presented by Polcaro and Viotti (1991). However, Della Valle, Benetti & Panagia (1996) demonstrated convincingly that the absorption was caused by the parent galaxy, rather than by the SN environment.
2. In spite of some impressions to the contrary, many double WD systems do exist. In a sample of 153 field WDs and subdwarf B stars, Saffer, Livio & Yungelson (1998) found 18 new double-degenerate candidates. Maxted & Marsh (1999) showed (from a radial velocity survey of 46 WDs) that there is a 95% probability that the fraction of double degenerates among DA WDs lies in the range 0.017โ0.19. There are currently eight known systems with orbital periods of less than half a day (and the subdwarf B stars PG 1432+159 and PG 2345+318, with orbital periods of 5.4 hr and 5.8 hr respectively may also have WD companions; Moran et al. 1999). While only one of all of these systems (KPD 0422+5421; Koen, Orosz & Wade (1998)) has a total mass which within the errors could be higher than the Chandrasekhar mass, the sample of confirmed short-period double-degenerates is still smaller than the number predicted to contain a massive system.
3. Population synthesis calculations predict the right statistics for mergers, about $`10^3`$ yr<sup>-1</sup> events for populations that are $`10^8`$ yr old and $`10^4`$ yr<sup>-1</sup> for populations that are $`10^{10}`$ yr old.
4. Since double WD systems were found to exist, mergers with some โinterestingโ consequences (either a SN Ia or an accretion-induced collapse) appear inevitable.
5. The explosion or collapse is expected to occur at (or near) the Chandrasekhar mass, which as I noted in ยง4.3, I regard as a property of the favored model.
The main weaknesses of the double-degenerate scenario are the following:
1. There are strong indications that WD mergers may lead to off-center carbon ignition, accompanied by the conversion of the CโO WD to an OโNeโMg composition, followed by an accretion-induced collapse rather than a SN Ia (e.g. Mochkovitch & Livio 1990; Saio & Nomoto 1985, 1998; Woosley & Weaver 1986).
2. Galactic chemical evolution results, and in particular the behavior of the \[O/Fe\] ratio as a function of metallicity (\[Fe/H\]) have been claimed to be inconsistent with WD mergers as the mechanism for SNe Ia (Kobayashi et al. 1998).
3. While the unusually high luminosity of SN 1991T and some of its other features have been tentatively attributed to a super-Chandrasekhar product of the merger of two WDs (Fisher et al. 1999), there is little evidence for example for the presence of unburned carbon (as might be expected from the disk formed in the merger process) in most SNe Ia.
Since we are now getting to the final stages in the identification of the progenitors, it is important to assess critically the severity of the above weaknesses. I will therefore discuss now each one of them in some detail.
#### 5.1.1 Constraints from Galactic chemical evolution
Supernovae Type II (SNe II) are explosions resulting from the core collapse of massive ($`\stackrel{>}{}8`$ M) stars. These supernovae produce relatively more oxygen and magnesium than iron (\[O/Fe\] $`>0`$). On the other hand, SNe Ia produce mostly iron and little oxygen. Generally, the impression is that metal poor stars (\[Fe/H\] $`1`$) have a nearly flat relation of \[O/Fe\] vs. \[Fe/H\], with a value of \[O/Fe\] $`0.45`$ (e.g. Nissen et al. 1994), while disk stars (\[Fe/H\]$`\stackrel{>}{}1`$) show a linearly decreasing \[O/Fe\] with increasing metallicity (e.g. Edvardsson et al. 1993; McWilliam 1997). The โobservedโ (but see below) break (from flat to linearly decreasing) near \[Fe/H\] $`1`$ is traditionally explained by the fact that the early heavy element production was done exclusively by SNe II, with the break occurring when the larger Fe production by SNe Ia kicks in (e.g. Matteucci & Greggio 1986).
Recently, Kobayashi et al. (1998) performed chemical evolution calculations for both the double-degenerate scenario and for the single-degenerate scenario. For the latter they used two types of progenitor systems: one with a red giant companion and an orbital period of tens to hundreds of days, and the other with a near main sequence companion and a period of a few tenths of a day to a few days.
They obtained for the double-degenerate scenario (for which they took a time delay to the explosion of $`0.1`$โ0.3 Gyr) a break at \[Fe/H\] $`2`$. For the single-degenerate scenario (with a delay caused by the main sequence lifetime of $`\stackrel{>}{}1`$ Gyr; including metallicity effects), they obtained a break at \[Fe/H\] $`1`$. Kobayoshi et al. (1998) thus concluded that the Galactic chemical evolution that results from the double-degenerate scenario is inconsistent with observations.
Personally, I am not too convinced by this apparent discrepancy, since Galactic chemical evolution calculations and observations are notoriously uncertain. For example, a recent determination of \[Ba/Fe\] as a function of \[Fe/H\] shows a break near \[Fe/H\] $`2`$, which would be consistent with the double-degenerate scenario prediction (Burris et al. 1999). In addition, recent Keck observations of oxygen in unevolved metal-poor stars appear to show no break in the \[O/Fe\] vs. \[Fe/H\] relation. Rather, oxygen is enhanced relative to iron over three orders of magnitude in \[Fe/H\] in a linear relation (Boesgaard et al. 1999; see also Israelian, Garcia Lopez & Rebolo 1998). While some reservations about these findings have been raised, in particular, a re-analysis of two of the stars of Israelian et al. shows \[O/Fe\] ratios which are discrepant with the results of Israelian et al. and of Boesgaard et al. (Fulbright & Kraft 1999), this in fact demonstrates the uncertainties involved in such determinations (see also Stephens 2000).
#### 5.1.2 Merger only applicable to relatively rare events?
As I noted above, it has been shown that if SN 1991T is at the same distance as SNe 1981B and 1960F, then its luminosity is too high to be explained in terms of a Chandrasekhar mass ejection (Fisher et al. 1999). Thus, it has been suggested that this SN resulted from the explosion of a super-Chandrasekhar object, indicating perhaps that WD mergers may be responsible for at least some SNe Ia.
However, events like SN 1991T, which seem to be associated with regions of active star formation, represent at most $`20`$% of the SNe Ia (Li et al. 2000), and therefore, even if they are the results of mergers this still does not mean that WD mergers are the main class of progenitors of SNe Ia. In addition, it is still far from clear whether mergers can lead to explosions at all (see ยง5.1.3 below). Incidentally, data for a cepheid distance to NGC 4527 (which will help determine the true intrinsic luminosity of SN 1991T) have been obtained with HST, and the analysis is in progress (Saha et al. 2000).
#### 5.1.3 SN Ia or accretion induced collapse?
Potentially the most serious (and possibly even fatal) weakness of the double-degenerate scenario comes from the fact that some estimates and calculations indicate that the coalescence of two CโO WDs may lead to an accretion-induced collapse rather than to a SN explosion (e.g. Mochkovitch & Livio 1990; Saio & Nomoto 1985, 1998; Kawai, Saio & Nomoto 1987; Timmes, Woosley & Taam 1994; Mochkovitch, Guerrero & Segretain 1997).
The point is the following: once the lighter WD fills its Roche lobe, it is dissipated within a few orbital periods (Benz et al. 1990; Rasio & Shapiro 1995; Guerrero 1994), and it forms a hot thick disk configuration around the more massive white dwarf. This disk is mainly rotationally supported and hence central carbon ignition does not take place immediately, but rather the subsequent evolution depends largely on the rate of angular momentum transport and removal, since they determine the accretion rate onto the primary WD. As long as the accretion rate is higher than about $`\dot{M}\stackrel{>}{}2.7\times 10^6`$ M yr<sup>-1</sup>, carbon is ignited off-center (at the core-disk boundary; this may happen during the merger itself; e.g. Segretain 1994). Under such conditions, the flame was found (in spherically symmetric calculations) to propagate all the way to the center within a few thousand years, thus burning the CโO into an OโNeโMg mixture with no explosion (i.e. before carbon is centrally ignited; e.g. Saio & Nomoto 1998). Such configurations are expected to collapse (following electron captures on <sup>24</sup>Mg) to form neutron stars (Nomoto & Kondo 1991; Canal 1997). The main questions are therefore:
1. What accretion rates can be expected from the initial WD-thick disk configuration?
2. May some aspects of the flame propagation be different given the fact that the real problem is three-dimensional while most of the existing calculations were performed using a spherically symmetric code? In particular, could the carbon burning be quenched before the transformation to OโNeโMg composition occurs?
3. Could the WDs ignite even prior to the merger due to tidal heating, and what would be the outcome of such pre-merger ignition?
The answers to all of these questions involve uncertainties, however some possibilities appear more likely than others. First, it appears very difficult to avoid high accretion rates. If the MHD turbulence that is expected to develop in accretion disks (e.g. Balbus & Hawley 1998) is operative, with a corresponding viscosity parameter of $`\alpha 0.01`$ (where the viscosity is given by $`\nu \alpha c_\mathrm{s}H`$, with $`H`$ being a vertical scaleheight in the disk and $`c_\mathrm{s}`$ the speed of sound; e.g. Balbus, Hawley & Stone 1996), then angular momentum can be removed in a matter of days! In such a case, even if the accretion rate is Eddington limited (at $`10^5`$ M/yr), off-center carbon ignition should still occur, with an eventual collapse rather than an explosion. Deviations from spherical symmetry can only hurt, since they may allow accretion to proceed at a super-Eddington rate. It is difficult to see why the dynamo-generated viscosity would be suppressed for the kind of shear and temperatures expected in the disk.
Concerning the burning itself, recent attempts at multi-dimensional calculations of the flame propagation and a more detailed analysis of some of the processes involved (Garcia-Senz, Bravo & Serichol 1998; Bravo & Garcia-Senz 1999) indicate that if anything, accretion induced collapses are an even more likely outcome than previously thought. This is due to the effects of electron captures in Nuclear Statistical Equilibrium which tend to stabilize the thermonuclear flame, and to Coulomb corrections to the equation of state. The latter has the effect of reducing the flame velocities and the electronic and ionic pressures, all of which result in a reduction in the critical density which separates explosions from collapses.
As two WDs approach merger, their interiors can be spun up by tidal forces. If these tides can bring (at least one of) the WDs into quasi synchronism between the spin period and the orbital period, high dissipation rates and heating will ensue (Rieutord & Bonazolla 1987). The obtained luminosities due to this tidal heating can reach values as high as $`\stackrel{>}{}10^{37}`$ erg s<sup>-1</sup> (Rieutord & Bonazzola 1987; Iben, Tutukov & Federova 1998). If such heating indeed occurs, it could have (in principle at least) two important effects: (i) it would increase the probability of detection of pre-merger WDs, due to the increased luminosity and the expected periodic variability (due to mutual occultations), (ii) heating could lead to carbon ignition prior to or during the merger.
ยฟFrom the point of view of the present discussion it is important to assess whether the latter possibility makes the merging WDs more viable progenitor systems. This does not appear to be the case for the following reasons:
1. It is not obvious that a WD can be brought to synchronous rotation. The normal viscosity of WD matter is very low (e.g. Durisen 1973), which can make the viscous timescale even longer than the systemโs lifetime. This problem however may be overcome if turbulence develops due to the strong shear.
2. It is not clear if carbon will be ignited even if tidal heating occurs. In fact, in the calculations of Iben et al. (1998) carbon failed to be ignited (although only by a relatively narrow margin).
3. Even if carbon is ignited, it is very likely that the ignition will occur off-center, making an accretion induced collapse a more likely outcome than a SN Ia (as explained above).
Finally, on the observational side there are also two points which argue at some level against WD mergers as the main SNe Ia progenitors.
1. Even if MHD viscosity could somehow be suppressed in the disk, and the disk surrounding the primary WD could cool down, so that angular momentum would be transported only via the viscosity of (partially) degenerate electrons, this would result in an accretion timescale of $`10^9`$ yrs (Mochkovitch & Livio 1990; Mochkovitch et al. 1997). The system prior to the explosion would have an absolute magnitude of $`M_\mathrm{V}\stackrel{<}{}10`$ (with much of the emission occurring in the UV). There is no evidence for the existence of some $`10^7`$ such objects in the Galaxy.
2. The existence of planets around the pulsars PSR 1257+12 and PSR 1620โ26 (Wolszczan 1997; Backer 1993; Thorsett, Arzoumanian & Taylor 1993) could be taken to mean (this is a model dependent statement) that mergers tend to produce accretion induced collapses rather than SNe Ia. In one of the leading models for the formation of such planets (Podsiadlowski; Pringle & Rees 1991; Livio, Pringle & Saffer 1992), the planets form in the following sequence of events. The lighter WD is dissipated (upon Roche lobe overflow) to form a disk around the primary. As material from this disk is accreted, matter at the outer edge of the disk has to absorb the angular momentum, thereby expanding the disk to a large radius. The planets form from this disk in the same way that they did in the solar system, while the central object collapses to form a neutron star.
#### 5.1.4 Overall assessment of the double-degenerate scenario
It has now been observationally demonstrated that many double-degenerate systems exist. The general agreement between the distribution of the observed properties (e.g. orbital periods, masses) and those predicted by population synthesis calculations (Saffer, Livio & Yungelson 1998), suggests that the fact that no clear candidate (short period) system with a total mass exceeding the Chandrasekhar mass has been found yet, may merely reflect the insufficient size of the observational sample. Thus, there is very little doubt in my mind that statistics is not a serious problem. The most disturbing uncertainty is related to the outcome of the merger process itself. The discussion in ยง5.1.3 suggests that collapse to a neutron star is more likely than a SN Ia (see also Mochovitch et al. 1997).
### 5.2 The single-degenerate scenario
The main strengths of the single degenerate scenario are:
1. A class of objects in which hydrogen is being transferred at such high rates that it burns steadily on the surface of the WD has been identifiedโthe Supersoft X-ray Sources (e.g. Greiner, Hasinger & Kahabka 1991; van den Heuvel et al. 1992; Southwell et al. 1996; Kahabka and van den Heuvel 1997). If the accreted matter can indeed be retained, this provides a natural path to an increase in the WD mass towards the Chandrasekhar mass (e.g. Di Stefano & Rappaport 1994; Livio 1995, 1996a; Yungelson et al. 1996).
2. Other candidate progenitor systems are known to exist, like symbiotic systems (e.g. Munari & Renzini 1992; Kenyon et al. 1993; Hachisu, Kato & Nomoto 1999) and recurrent novae (Hachisu et al. 1999a).
3. There have been claims that the single degenerate scenario fits better the results of Galactic chemical evolution (e.g. Kobayashi et al. 1998). However, as I have shown in ยง5.1.1, recent observations cast doubt on this assertion. Similarly, nucleosynthesis results show that in order to avoid unacceptably large ratios of <sup>54</sup>Cr/<sup>56</sup>Fe and <sup>50</sup>Ti/<sup>56</sup>Fe, the central density of the WD at the moment of thermonuclear runaway must be lower than $`2\times 10^9`$ g cm<sup>-3</sup> (Nomoto et al. 1997). Such low densities are realized for high accretion rates ($`\stackrel{>}{}10^7`$ M yr<sup>-1</sup>), which are typical for the Supersoft X-ray Sources. Nucleosynthesis results suffer too, however, from considerable uncertainties (e.g. Nagataki, Hashimoto & Sato 1998).
The main weaknesses of the single degenerate scenario are:
1. The upper limits on radio detection of hydrogen at 2 and 6 cm in SN 1986G, taken approximately one week before optical maximum (Eck et al. 1995), rule out a symbiotic system progenitor for this system with a wind mass loss rate of $`10^7\stackrel{<}{}\dot{M}_\mathrm{W}\stackrel{<}{}10^6`$ M yr<sup>-1</sup> (Boffi & Branch 1995). This in itself is not fatal, since SN 1986G is somewhat peculiar (e.g. Branch and van den Bergh 1993), and the upper limit on the mass loss rate is at the high end of observed symbiotic winds. An even less stringent upper limit from x-ray and H$`\alpha `$ observations exists for SN 1994D (Cumming et al. 1996).
2. There exists some uncertainty whether WDs can even reach the Chandrasekhar mass at all by the accretion of hydrogen (e.g. Cassisi, Iben & Tornambe 1998). Furthermore, even if they can, the question of whether they can produce the required SNe Ia statistics is highly controversial (e.g. Yungelson et al. 1995, 1996; Yungelson & Livio 1998, 1999; Hachisu, Kato & Nomoto 1999; Hachisu et al. 1999a).
I will now examine these weaknesses in some detail.
#### 5.2.1 Observational detection of hydrogen
Ultimately, the presence or total absence of hydrogen in SNe Ia will distinguish unambiguously between single-degenerate and double-degenerate models. To date, hydrogen has not been convincingly detected in any SN Ia. It is interesting to note that narrow $`\lambda 6300`$, $`\lambda 6363`$ \[OI\] lines were observed only in one SN Ia (SN 1937C; Minkowski 1939), but even in that case there was no hint of a narrow H$`\alpha `$ line. Hachisu, Kato & Nomoto (1999) estimate in one of their models (which involves stripping of material from the red giant; see below) a density measure of $`\dot{M}/v_{10}10^8`$ M yr<sup>-1</sup> (where $`v_{10}`$ is the wind velocity in units of 10 km s<sup>-1</sup>), while the most stringent radio upper limit existing currently (for SN 1986G) is $`\dot{M}/v_{10}10^7`$ M yr<sup>-1</sup> (Eck et al. 1995; for SN 1994D Cumming et al. (1996) find from H$`\alpha `$ an upper limit of $`\dot{M}1.5\times 10^5`$ M yr<sup>-1</sup> for a wind speed of 10 km s<sup>-1</sup>; for SN 1992A Schlegel & Petre (1993) find from X-ray observations an upper limit of $`\dot{M}/v_{10}=(23)\times 10^6`$ M yr<sup>-1</sup>). Thus, while it is impossible at present to rule out single-degenerate models on the basis of the apparent absence of hydrogen, the hope is that near future observations will be able to determine definitively whether this absence is real or if it merely represents the limitations of existing observations (an improvement by two orders of magnitude in the sensitivity will give a definitive answer). I should note that a narrow emission feature possibly corresponding to H$`\alpha `$ was detected in SN 1981b (in NGC 4536), however no trace of the emission was seen 5 days later (Branch et al. 1983).
#### 5.2.2 Statistics
Growing the WD to the Chandrasekhar mass is not easy. At accretion rates below $`10^8`$ M/yr WDs undergo repeated nova outbursts (e.g. Prialnik & Kovetz 1995), in which the WDs lose more mass than they accrete between outbursts (e.g. Livio & Truran 1992). For accretion rates in the range $`10^8`$โa few $`\times 10^7`$ M/yr, while helium can accumulate, the WDs experience mass loss due to helium shell flashes and due to the common envelope phase which results from the engulfing of the secondary star in the expanding envelope (with mass loss occurring due to drag energy deposition). At accretion rates above a few $`\times 10^7`$ M/yr, the WDs expand to red giant configurations and lose mass due to drag in the common envelope and due to winds (e.g. Cassisi et al. 1998). The net result of these constraints has been that population synthesis calculations which follow the evolution of all the binary systems in the Galaxy, tended until recently to conclude that single degenerate channels manage to bring WDs to the Chandrasekhar mass only at about 10% of the inferred SNe Ia frequency of $`4\times 10^3`$ yr<sup>-1</sup> (e.g. Yungelson et al. 1995, 1996; Yungelson & Livio 1998; Di Stefano et al. 1997; although see Li & van den Heuvel 1997).
Very recently, a few serious attempts have been made to investigate whether the statistics could be improved by increasing the โphase spaceโ for single degenerate scenarios, given the fact that population synthesis calculations involve many assumptions. These attempts resulted in the identification of three directions in which the phase space could (potentially) be increased.
1. The accumulation efficiency of helium has been recalculated using OPAL opacities (Kato & Hachisu 1999). These authors concluded that helium can accumulate much more efficiently than found by Cassisi et al. (1998), mainly because the latter authors used relatively low WD masses (0.516 M and 0.8 M) and old opacities in their calculations.
2. Hachisu et al. (1999a,b) claimed to have identified two evolutionary channels for single-degenerate systems previously overlooked in population synthesis calculations. In the first of these channels, the CโO WD is formed from a red giant with a helium core of 0.8โ2.0 M (rather than from an asymptotic giant branch star with a CโO core). The immediate progenitors in this case are expected to be either helium-rich Supersoft X-ray Sources or recurrent novae of the U Sco subclass (where the accreted material appears to be helium rich).
3. In the second channel, Hachisu et al. (1999b) considered very wide (initial separations as large as $`40,000`$ R) symbiotic systems, in which the components are brought together by the inclusion of new physical effects (see (iii)โ(3) below).
4. It has been suggested that the inclusion of a few additional physical effects, can increase substantially the phase space of the symbiotic channel (Hashisu, Kato & Nomoto 1996, 1999b). These new effects included:
1. The WD loses much of the transferred mass in a massive wind. This has the effect that the mass transfer process is stabilized for a wider range of mass ratios, up to $`q_{\mathrm{max}}m_2/m_1=1.15`$ instead of $`q_{\mathrm{max}}=0.79`$ without the massive wind.
2. It has been suggested that the wind from the WD strips the outer layers of the red giant at a high rate. This increases the allowed mass ratios (for stability) even above 1.15, essentially indefinitely.
3. It has been suggested that at large separations (up to $`40,000`$ R), when the orbital velocity is of the order of the wind velocity, the wind from the red giant acts like a common envelope to reduce the separation, thus allowing much wider initial separations to result in interaction.
There are many uncertainties associated with all of these attempts to increase the phase space. For example, the efficiency of mass stripping from the giant by the wind from the WD may be much smaller than assumed by Hachisu et al. (1999b), for the following reasons. At high accretion rates, much of the mass loss from the WD may be in the form of an outflow or a collimated jet, perpendicular to the accretion disk rather than in the direction of the giant. Evidence that this is the case is provided by the jet satellite lines to He II 4686, H$`\beta `$ and H$`\alpha `$ observed in the Supersoft X-ray Source RX J0513.9$``$6951 (Southwell et al. 1996). These jet lines are very similar to those seen in the prototypical jet source SS 433 (e.g. Vermeulen et al. 1992). Furthermore, even if some of the WD wind hits the surface of the giant, it is not clear how efficient it would be in stripping mass, since the rate of energy deposition per unit area by the wind is smaller by two orders of magnitude that the giantโs own intrinsic flux.
Similarly, the efficiency of helium accumulation is still highly uncertain, as the differences between the results of Kato & Hachisu (1999) and Cassisi et al. (1998) have shown.
Also, the particular form of the specific angular momentum in the wind used by Hachisu, Kato & Nomoto (1999b; point (3) above) may be realized only in relatively rare cases (Yungelson & Livio 2000).
Finally, all the new suggestions for the increase in phase space rely very heavily on the results of the wind solutions of Kato (1990; 1991), which involve a treatment of the radiation and hydronamics not nearly as sophisticated as that of more state of the art radiative transfer codes (e.g. Hauschildt et al. 1995, 1996).
#### 5.2.3 Overall assessment of the single-degenerate scenario
The above discussion suggests that probably not all the scenarios for increasing the โphase spaceโ of the single-degenerate channels work (if they did, we might have had the opposite problem of too high a frequency of SNe Ia!). However, these attempts serve to demonstrate that the input physics to population synthesis codes still involves many uncertainties. My feeling is therefore that given the many potential channels leading to SNe Ia, statistics should not be regarded as a serious problem.
Single-degenerate scenarios therefore appear quite promising, since unlike the situation a decade ago, a class of objects in which the WDs accrete hydrogen steadily (the Supersoft X-Ray Sources) has actually been identified. The main problem with single-degenerate scenarios remains the non-detection of hydrogen so far. While a difficult observational problem (see ยง6), the establishment of the presence or absence of hydrogen in SNe Ia should become a first priority for SNe observers.
### 5.3 What ifโฆ.?
Given the fact that there are still uncertainties involved in identifying the SNe Ia progenitors, and that WD mergers and some form of off-center helium ignitions almost certainly occur, it is instructive to pose a few โwhat ifโ questions. For example: What if WD mergers with a total mass exceeding Chandrasekhar do not produce SNe Ia, what do they produce then? The answer in this case will have to be that they almost certainly produce either neutron stars via accretion induced collapses, or single WDs, if the merger is accompanied by extensive mass loss from the system. Fryer et al. (1999) estimate (from nucleosynthesis constraints) that less than 0.1% of the total Galactic neutron star population is produced via accretion induced collapses.
What if off-center helium ignitions (ELDs) do not produce SNe Ia? In this case, if an explosive event indeed ensues, a population of โsuper novaeโ (with $`0.15`$ M of <sup>56</sup>Ni and He) is yet to be detected (maybe SN 1885A in M31 was such an event?). What if off-center helium ignitions do produce SNe Ia? What comes out of the systems with $`M_{\mathrm{WD}}\stackrel{>}{}1`$ M, which should be even brighter? If indeed $`20`$% of SNe Ia are SN 1991T-like (Li et al. 2000), then maybe these could be represented by such events. This is far from certain, however, since an analysis of the properties of SN 1991T showed that these properties would be very difficult to reproduce even with a nickel mass approaching the Chandrasekhar mass (Fisher et al. 1999). Thus, we see that off-center helium ignitions seem to present an observational problem both if they do and if they do not produce SNe Ia. To me this suggests that the physics of these events is not yet well understood (for example, maybe off-center helium ignition fails to ignite the CโO core after all).
## 6 How can we hope to unambiguously identify the progenitors?
There are several ways in which observations of both nearby and distant supernovae could solve the mystery of SNe Ia progenitors:
1. A combination of early high resolution optical spectroscopy, x-ray observations and radio observations of nearby SNe Ia can both provide limits on $`\dot{M}/v`$ from the progenitors and potentially detect the presence of circumstellar hydrogen (if it exists).
2. For example, narrow HI in emission or absorption could be detected either very early, or shortly after the ejecta become optically thin ($`100`$ days). The latter is true because the SN ejecta probably engulfs the companion at early times (e.g. Chugai 1986; Livne, Tuchman & Wheeler 1992). The interaction of the ejecta with the circumstellar medium can be observed either in the radio (e.g. Boffi & Branch 1995) or in x-rays (e.g. Schlegel 1995). The collision of the ejecta (with circumstellar matter) can also set up a forward and a reverse shock (e.g. Chevalier 1984; Fransson, Lundqvist & Chevalier 1996), and radiation from the latter can ionize the wind and produce H$`\alpha `$ emission (e.g. Cumming et al. 1996).
3. Early observations (again of nearby SNe Ia) of the gamma-ray light curve (or gamma-ray line profiles) could distinguish between carbon ignitors and sub-Chandrasekhar helium ignitor models (see ยง4.2.2) since the latter can be expected to result in a quicker rise of the gamma-ray light curve due to the presence of <sup>56</sup>Ni in the outer layers (and different gamma-ray line profiles; because of the high velocity <sup>56</sup>Ni).
4. Another important aspect of the single degenerate scenario which can be tested by observations of both nearby and very distant supernovae is the dependence on metallicity. The increase in the โphase spaceโ of single degenerate progenitors, which is required to make the statistics more compatible with observations (ยง5.2.2), relies heavily on the existence of an optically thick wind from the WD. For a low metallicity of the accreted matter (\[Fe/H\]$`\stackrel{<}{}1`$), the wind from the WD is strongly suppressed (since the wind is driven by a peak in the opacity, which is due to iron lines; e.g. Kobayashi et al. 1998). Consequently, it is expected that SNe Ia rates will be significantly lower in low iron abundance environments. Thus, determinations of relative rates in dwarf galaxies (and in the very outskirts of spiral galaxies) can help determine the viability of a key ingredient in the single degenerate scenario.
5. Similarly, a significant drop may be expected in the rate of SNe Ia at redshift $`z2`$ (again due to the decrease in metallicity).
6. In general, observations of very distant supernovae (at $`z2`$โ4) with the Next Generation Space Telescope (NGST) can help significantly in identifying the progenitors (e.g. Yungelson & Livio 1999, 2000; Nomoto et al. 2000). For example, the progenitors can be identified from the observed frequency of SNe Ia as a function of redshift (e.g. Yungelson & Livio 1998, 1999, 2000; Ruiz-Lapuente & Canal 1998; Madau, Della Valle & Panagia 1998; Nomoto et al. 2000; and see ยง8), since different progenitor models produce different redshift distributions. Personally, I think that it would be quite pathetic to have to resort to this possibility. Rather, one would like to be able to identify the progenitors independently, and then use the observations of supernovae at high $`z`$ to constrain models of cosmic star formation rates, and of cosmic evolution of SNe rates, luminosity, and input into galaxies.
## 7 Cosmological implications: could we be fooled?
One of the key questions that result from the uncertainties in the theoretical models and the fact that we do not know with certainty which systems are the progenitors of SNe Ia is clearly: is it possible that SNe Ia at higher redshifts are systematically dimmer than their low-redshift counterparts? In this respect it is important to remember that a systematic decrease in the brightness by $`0.25`$ magnitudes is sufficient to explain away the need for a cosmological constant. This question became particularly relevant when an analysis of the rise times of SNe Ia (which was based on preliminary estimates for the high-$`z`$ sample; by Goldhaber 1998 and Groom 1998) seemed to show that high-redshift SNe have shorter rise times by 2.5 days than the low-$`z`$ SNe (Riess et al. 1999b). A more recent analysis (which used more realistic error estimates than those used by Goldhaber 1998), however, found a better agreement (within 2$`\sigma `$) between the rise times of the low- and high-redshift SNe Ia (Aldering, Knop & Nugent 2000).
In a recent work, Yungelson & Livio (1999) calculated the expected ratio of the rate of SNe Ia to the rate of SNe from massive stars (Types II, Ia, Ic) as a function of redshift for several progenitor models. The possibility of having different classes of progenitors contributing to the total SNe Ia rate should definitely be considered, especially in view of the tentative finding by Li et al. (2000) that there is a relatively high rate ($`40`$%) of peculiar SNe Ia among the local sample. If confirmed, these findings suggest that homogeneity should no longer be considered a very strong constraint on progenitor models. I should note though that diversity among SNe Ia does not necessarily imply different progenitors, since even in the context of one progenitor model diversity may arise for example from changes in the carbon mass fraction of the WD, which in turn may depend on the environment (e.g. Nomoto et al. 2000). Yungelson and Livio (1999) showed (within the uncertainties of population synthesis models) that if different progenitor systems can contribute to the total SNe Ia rate (e.g. double-degenerates and single degenerates), then it is possible, in principle, that one class of progenitors (e.g. double degenerates) will dominate the rates of the local (low-$`z`$) sample, while a different progenitor class (e.g. single degenerate) will start to dominate at $`z\stackrel{<}{}1`$. This is a consequence of the fact that the SNe Ia rate from double degenerates is expected to decline quite steeply from $`z=0`$ to $`z1`$, while the rate from single degenerates is expected to stay relatively flat in this redshift interval. However, at least within the assumptions of their model calculations, it appears that such a transition is not very likely, because of the following reason: If the contribution from physically different channels (like double degenerates and single degenerates, both at the Chandrasekhar mass) was indeed significant, with a transition from dominance by double degenerates to single degenerates occurring at $`z\stackrel{<}{}1`$, one would have expected to observe this division more clearly in the local and distant samples. For example, the local sample should be dominated by the double degenerate progenitors (with a ratio of double to single degenerates which should be consistent with the results of Li et al. 2000). At the same time, however, the high-$`z`$ ($`z\stackrel{<}{}1`$) sample should be dominated by single degenerates, but with the contributions from the single and double degenerate channels not being vastly different (in particular, the contribution should be equal at the transition point). This is not consistent with the observations of the high-$`z`$ sample, the latter appearing to be (within the observational uncertainties) very homogeneous (Li et al. 2000). Consequently, I do not think it likely that the observed universal acceleration is an artifact of the observed SNe Ia sample being dominated by different progenitor classes at high- and low-$`z`$ (this view is supported by the measurements of the anisotropy of the microwave background).
I should note that the most surprising aspect in the results of Li et al. (2000) is the fact that although very bright SNe Ia (SN 1991T-like) constitute ($`21\pm 7`$)% of the local sample, these bright objects appear to be totally absent from the high-$`z`$ sample. One way in which one could (in principle) explain this fact is the following. Suppose that the SN 1991T-like events are caused by mergers (since a super-Chandrasekhar mass is possible in this case), while the โnormalsโ are caused by single degenerates. In this case the local sample would be dominated by single degenerates ($`60`$% of the events to agree with Li et al., 2000), while double degenerates would contribute $`20`$% of the events (I ignore here the weak events since they cannot be seen at high-$`z`$). Now, since the rate of events from single degenerates stays quite flat till $`z1`$, while the rate of events from double degenerates declines quite steeply towards $`z1`$, the bright objects will be missing from the high-$`z`$ sample. Note, however, that this potential explanation for the behavior of the diversity has no obvious implications for the finding of accelerating expansion, since the same class of progenitors dominates both the high-$`z`$ and low-$`z`$ samples. Nevertheless, a better understanding of the apparent diversity at low-$`z`$ and apparent lack thereof at high-$`z`$ is definitely needed.
Other evolutionary effects are still possible, in principle (e.g. Drell, Loredo & Wasserman, 1999; Hillebrandt 2000), however, as far as I am aware, only one that is physically meaningful, likely, and mimics accelerated expansion, has been identified so far.
This one potential evolutionary effect that certainly deserves more work is the effect of metallicity on the density at the point of carbon ignition. Generally, it is expected that a lower metallicity will result in a lower central density (e.g. Nomoto et al. 1997). This is because a lower metallicity results in a lower abundance of the Urca-active element <sup>21</sup>Ne, which in turn reduces the neutrino cooling and leads to an earlier ignition. A lower central density could (in principle at least) result in a more rapid light curve development (due to the lower WD binding energy), and a lower inferred maximum brightness.
It is important to note that in a recent work, Riess et al. (2000) have shown that it is highly unlikely that the dimming of distant SNe Ia is caused by dust opacity (Galactic-type dust was rejected at the 3.4$`\sigma `$ confidence level, and โgrayโ dust with grain size $`>0.1\mu `$m was rejected at the 2.3 to 2.6$`\sigma `$ confidence level).
## 8 Tentative conclusions and observational tests
On the basis of the analysis and discussion in the present work, the following tentative conclusions can be drawn:
1. SNe Ia are almost certainly thermonuclear disruptions of mass accreting CโO white dwarfs.
2. It is very likely that the explosion occurs at the Chandrasekhar mass, as carbon is ignited at (or very near) the WD center. The flame propagates either as a deflagration, or, more likely perhaps, starting as a deflagration which transitions into a detonation. Off-center ignition of helium at sub-Chandrasekhar masses may still be responsible for a subset of the SNe Ia which are subluminous, but this is less clear.
3. The immediate progenitor systems are still not known with certainty. From the discussion in ยง5 (see in particular ยง5.1.3 and 5.2.3) however, I conclude that presently single degenerate scenarios look more promising, with hydrogen or helium rich material being transferred from a subgiant or giant companion (systems like Supersoft X-Ray Sources and Symbiotics). It is still possible, however, in view of the apparent diversity in the local sample, that more than one progenitor class contributes to the total SNe Ia rate. In particular, a scenario in which single degenerates contribute $`\stackrel{>}{}60`$% of the events and double degenerates $`20`$%, appears to be consistent with the diversity of the $`z0`$ sample and the lack thereof in the distant sample.
4. Definitive answers concerning the nature of the progenitors can be obtained from observations taken as early as possible in: x-rays, radio, and high resolution optical spectroscopy. The establishment of the presence or absence of hydrogen in SNe Ia should be regarded as an extremely high priority goal for supernovae observers. If hydrogen will not be detected at interesting limits (corresponding to $`\dot{M}/v_{10}10^8`$ M yr<sup>-1</sup>), this will point clearly towards the double-degenerate scenario.
5. Observations of SNe Ia at high redshifts can help to test particular ingredients of the models which are directly related to the nature of the progenitors. For example, most of the models aiming at improving the statistics of the single-degenerate scenarios rely on a strong wind from the accreting WD. These models thus predict an โinhibitionโ of SNe Ia in low-metallicity environments, and in particular a significant decrease in the rate of SNe Ia in spirals at $`z2`$ (Kobayashi et al. 1998; Nomoto et al. 2000). Furthermore, if the inferred cosmic star formation rate is used (e.g. Pettini et al. 1998), then the SN Ia rate is expected to drop significantly at $`z1.6`$. At present, the detection of a very likely SN Ia at redshift $`z=1.32`$ (SN 1997ff; Gilliland, Nugent & Phillips 1999) in the Hubble Deep Field, and two more at redshifts 1.20 and 1.23 (Perlmutter et al., private communication and Tonry et al., private communication) appear to be at least mildly inconsistent with this prediction, but more observations will be required to give a more definitive answer.
6. The potential โinhibitionโ of SNe Ia due to low metallicity should also manifest itself in an absence (or at least a significant decline in the rate) of SNe Ia in dwarf galaxies and in the outer regions of spirals. The statistics necessary to test this prediction are starting to accumulate.
7. It is possible, in principle, that the local and high-$`z`$ SNe Ia samples are dominated by different progenitor classes, thus mimicking accelerated expansion, however, this is neither very likely nor consistent with the observations of diversity. With more detections of SNe Ia at redshifts $`z\stackrel{>}{}1`$ it will probably become possible to directly confirm the transition in the expansion of the universe from deceleration to acceleration. Such a transition would be difficult to mimic by systematic or evolutionary effects, and it would therefore confirm the accelerated expansion.
###### Acknowledgements.
This research has been supported in part by NASA Grant NAG5โ6857. I am grateful to Adam Riess for helpful discussions.
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# Identifying the pairing symmetry in the Sr2RuO4 superconductor
## I Introduction
The search for the superconducting pairing symmetry in the layered perovskite material Sr<sub>2</sub>RuO<sub>4</sub> (SrRuO), and its attempted theoretical predictions, show remarkable parallels to the heavy-fermion superconductor UPt<sub>3</sub>. In both systems, early specific-heat measurements showed a large residual value of $`C/T`$ at low temperatures and were interpreted in terms of a superconducting phase with a nonunitary $`p`$-wave order parameter. The observation of a strong $`T_c`$ suppression with nonmagnetic impurities was an additional indication of a superconducting phase with an unconventional order parameter. However, newer measurements on high-quality single crystals have shown that the most likely pairing state in UPt<sub>3</sub> is an $`f`$-wave state, or more precisely a spin-triplet state whose orbital basis function belongs to the $`E_{2u}`$ representation of the hexagonal crystallographic point group (D<sub>6h</sub>). The experience with UPt<sub>3</sub> suggests that the early identification of the pairing state, based on low-quality, inhomogeneous samples, is at best inconclusive (for a review on UPt<sub>3</sub> see, for example, Refs. ). However, with improving sample quality it becomes feasible to identify the pairing state by studying transport properties.
Here we analyze new heat capacity measurements on high-quality single crystals of SrRuO, as well as thermal conductivity data on dirty samples with a strong $`T_c`$ suppression, to show that the proposed $`p`$-wave model, $`๐ซ(๐ฉ_f)(p_x+ip_y)\widehat{๐ณ}`$, is inconsistent with the available data. Our conclusion is that the pairing state in SrRuO, most likely, has lines of nodes with gap nodes given by the $`d_{xy}`$ gap function. This can occur in either an $`f`$-wave state, i.e., a spin-triplet pairing state belonging to the $`E_u`$ representation of the tetragonal crystallographic point group (D<sub>4h</sub>) or in a $`d_{xy}`$ singlet state. We argue that the $`f`$-wave nodal state is consistent with measurements of the heat capacity, thermal conductivity, penetration depth, Andreev reflection, NMR, Knight shift, and $`\mu `$SR experiments.
Recent band-structure calculations by Mazin and Singh indicate that there is an increase in the spin susceptibility $`\chi (๐ช,\omega )`$ at four points in the Brillouin zone at approximately $`๐ช_0(\pm 2\pi /3,\pm 2\pi /3)`$ that occur due to strong nesting effects of quasi-one-dimensional bands ($`\xi `$ and $`\zeta `$). Nesting effects among these bands lead to the increased interaction between particles on the Fermi surface near $`๐ช_0`$, see Fig. 1. In recent neutron-scattering experiments the predicted four incommensurate peaks near $`๐ช_0`$ were indeed observed thus supporting that nesting effects near these points are important.
In this paper we propose: 1) To identify the regions at the Fermi surface near $`๐ช_0`$ with the ones that develop the largest gap. We use the neutron-scattering data as an indication that near the nesting regions the particle-particle (or particle-hole) interactions are dominant and that these are the regions that would benefit the most from opening a superconducting gap. 2) We suggest that regardless of the singlet or triplet nature of the pairing in SrRuO the gap function should be proportional to a $`d_{xy}`$ harmonics. Such an order parameter would lead to lines of nodes along the $`k_z`$-axis in the gap and to power-law behavior in the thermodynamic and transport properties. Line nodes on the Fermi surface lead in clean superconductors, and for scattering in the Born limit, to the well-known temperature dependences of the specific heat $`CT^2`$, the nuclear spin relaxation rate $`1/T_1T^3`$, the deviation of the penetration depth from its zero-temperature value $`\mathrm{}\lambda T`$, the thermal conductivity $`\kappa T`$, and the longitudinal sound attenuation $`\alpha _L\mathrm{const}.`$, as well as for the transverse attenuation $`\alpha _TT^2`$. 3) Based on the proposed line nodes in the gap we make predictions for ultrasound attenuation and thermal conductivity measurements that can further distinguish between the remaining possible basis functions. We propose complimentary longitudinal and transverse attenuation measurements that can help to locate the location of the nodal lines of the order parameter on the Fermi surface. Another crucial experiment is the thermal conductivity with an in-plane magnetic field. We expect the fourfold modulation of the thermal conductivity $`\kappa (\theta ,H)`$ as a function of the angle between the nodes of the gap \[along the (1,0) and (0,1) direction\] and the field directions. Thermal conductivity measures the unpaired quasiparticle heat transport and is therefore sensitive to the angular (field) dependence of the quasiparticle scattering rate, which โknowsโ about the angular dependence of the gap. We use the analogy with the suggested d-wave paring state in high-$`T_c`$ superconductors where this fourfold modulation has been observed.
## II Model
The gap function for even parity (spin-singlet) or odd parity (spin-triplet) representations is described by an order parameter of the form
$`\mathrm{\Delta }_{\alpha \beta }(๐ฉ_f)=\mathrm{\Delta }(๐ฉ_f)(i\sigma _y)_{\alpha \beta },`$ $`(\mathrm{singlet})`$ (1)
$`\mathrm{\Delta }_{\alpha \beta }(๐ฉ_f)=๐ซ(๐ฉ_f)(i๐\sigma _y)_{\alpha \beta },`$ $`(\mathrm{triplet})`$ (2)
with $`\sigma _\alpha `$ being Pauli matrices. Since nonunitary states, i.e., $`๐ซ\times ๐ซ^{}0`$, have been ruled out by the very small residual value of the specific heat $`C/T`$ at zero temperature, we restrict our study of spin-triplet states to unitary order parameters that factorize into a single spin vector and an orbital amplitude, i.e., $`๐ซ(๐ฉ_f)=๐\mathrm{\Delta }(๐ฉ_f)`$, where $`๐`$ is a real unit vector in spin space and $`\mathrm{\Delta }(๐ฉ_f)`$ is an odd-parity orbital function. The vector $`๐`$ defines the axis along which the Cooper pairs have zero spin projection, e.g., if $`๐||\widehat{๐ณ}`$, then $`\mathrm{\Delta }_{}=\mathrm{\Delta }_{}=0`$ and $`\mathrm{\Delta }_{}=\mathrm{\Delta }_{}=\mathrm{\Delta }(๐ฉ_f)`$.
Whether or not spin-orbit coupling is weak or strong in Sr<sub>2</sub>RuO<sub>4</sub> has important ramifications for both spin and orbital components of the order parameter that are allowed by symmetry. While spin-orbit coupling is believed to be strong in the heavy-fermion system UPt<sub>3</sub> there are no experimental indications that this is likewise true for Sr<sub>2</sub>RuO<sub>4</sub>. In the meantime we will use the classification of basis functions in terms of irreducible representations of the tetragonal point group ($`D_{4h}`$) listed in Table I, implying that spin-orbit coupling is strong. Since the band-structure calculations and de Haas - van Alphen measurements show very little dispersion along $`k_z`$, we will consider only two-dimensional (2D) basis functions on a more or less cylindrical Fermi surface. A similar list of possible basis functions was recently compiled by Hasegawa and co-workers for further investigations. The listed hybrid state (#3) of the direct product $`B_{2g}E_u=E_u`$ is a non-trivial realization of the $`E_u`$ representation (also referred to as $`f`$-wave state). So far Knight shift data with an in-plane magnetic field $`๐||[100]`$ show no change below $`T_c`$ and have been interpreted in terms of spin-triplet pairing with the spin vector $`๐`$ locked to the crystal $`๐`$-axis. On the other hand, muon spin rotation ($`\mu `$SR) experiments observed a spontaneous internal magnetic field on entering the superconducting state, consistent with a time-reversal symmetry breaking state belonging to the two-dimensional $`E_u`$ representation.
At this place a caveat is warranted because neither Knight shift data at high fields and for a single field orientation, nor $`\mu `$SR measurements in impure samples provide a clear-cut identification for spin-triplet pairing or broken time reversal symmetry states. For example, in UPt<sub>3</sub> early $`\mu `$SR measurements indicated broken time-reversal symmetry in the superconducting phase (probably due to impurities), while newer measurements on very clean samples fail to detect any effect at all. What makes the interpretation of the Knight shift data in SrRuO for magnetic fields parallel to the planes even more complicated, is, that (1) the experiment was not performed in the low field limit, but rather deep in the mixed phase, $`HH_{c2}/2`$, where contributions from the vortices may be important, and (2) nonlocal and surface effects may be relevant due to a small Ginzburg-Landau parameter for in-plane currents, $`\kappa _{||}=\lambda _{||}/\xi _{||}2.6`$.
## III Thermodynamic and transport properties
We calculate the specific heat and thermal conductivity for the order-parameter models listed in Table I and fit the results to existing experiments. This way, we can determine the model parameters and make predictions for sound attenuation measurements. It is important to point out that none of the here analyzed transport experiments can distinguish between a spin-singlet and a spin triplet order parameter. Thus we obtain identical results for the states #1 and #3.
The specific heat, $`C=TdS/dT`$, can easily be obtained from the entropy,
$$S=4_0^{\mathrm{}}๐ฯตN(ฯต)\left(\frac{ฯต}{T}f(ฯต)\mathrm{ln}[1f(ฯต)]\right),$$
(3)
by numerical differentiation. Here $`f(ฯต)=1/[1+\mathrm{exp}(ฯต/T)]`$ is the Fermi-Dirac function and $`N(ฯต)=(N_f/\pi )\mathrm{Im}๐๐ฉ_fg^R(๐ฉ_f,ฯต)`$ is the density of states per spin with $`N_f`$ being the normal-state density of states at the Fermi surface.
In the limit of Born (weak) or unitarity (strong) impurity scattering the in-plane thermal conductivity of unitary spin-triplet superconductors is given by
$`\kappa _{ii}`$ $`=`$ $`{\displaystyle \frac{N_fv_f^2}{8\pi ^3T^2}}{\displaystyle ๐ฯตฯต^2\mathrm{sech}^2\frac{ฯต}{2T}๐๐ฉ_f\widehat{๐ฏ}_{fi}^2๐ฆ(๐ฉ_f,ฯต)},`$ (4)
$`๐ฆ(๐ฉ_f,ฯต)={\displaystyle \frac{1}{\mathrm{Re}C^R}}\left[g^R(g^R)^{}๐^R(๐^R)^{}+\pi ^2\right],`$ (5)
with the unit vector of the Fermi velocity, $`\widehat{๐ฏ}_{fi}`$, and $`C^R=\frac{1}{\pi }\sqrt{|๐ซ|^2(\stackrel{~}{ฯต}^R)^2}`$. The quasiclassical equilibrium Green functions are $`g^R=\stackrel{~}{ฯต}^R/C^R`$ and $`๐^R=๐ซ/C^R`$. Within the $`t`$-matrix approximation for isotropic scattering the impurity renormalized quasiparticle energy is $`\stackrel{~}{ฯต}^R=ฯต\sigma _{imp}^R(ฯต)`$. For weak scattering $`\sigma _{imp}^R(ฯต)=(\mathrm{\Gamma }/\pi )๐๐ฉ_fg^R`$, and for strong scattering $`\sigma _{imp}^R(ฯต)=\mathrm{\Gamma }/(\pi d๐ฉ_fg^R`$), with the normal-state scattering rate $`\mathrm{\Gamma }=\mathrm{}/2\tau `$.
In the hydrodynamic regime, $`\omega \tau 1`$, and long wavelength limit, $`q\mathrm{}1`$, the absorption of ultrasound of polarization $`๐บ`$ propagating along direction $`๐ช`$ is related to the viscosity by
$$\alpha =\frac{\omega ^2}{\varrho c_s^3}\eta _{ij,kl}\widehat{๐บ}_i\widehat{๐ช}_j\widehat{๐บ}_k\widehat{๐ช}_l,$$
(6)
with the speed of sound $`c_s`$, the mass density $`\varrho `$, and the viscosity tensor evaluated at $`\omega 0`$,
$`\eta _{ij,kl}={\displaystyle \frac{N_fv_f^2p_f^2}{8\pi ^3T}}{\displaystyle ๐ฯต\mathrm{sech}^2\frac{ฯต}{2T}๐๐ฉ_f\pi _{ij}\pi _{kl}๐ฆ(๐ฉ_f,ฯต)},`$ (7)
where $`\pi _{ij}=\widehat{๐ฏ}_{fi}\widehat{๐ฉ}_{fj}\frac{1}{2}\delta _{ij}`$.
Here we confine our discussion to order parameters with vanishing averages, $`๐๐ฉ_f๐ซ(๐ฉ_f)=0`$, which satisfy the gap equation for triplet (singlet) pairing interactions,
$$๐ซ(๐ฉ_f)=\frac{dฯต}{2\pi }\mathrm{tanh}\frac{ฯต}{2T}๐๐ฉ_f^{}V(๐ฉ_f,๐ฉ_f^{})\mathrm{Im}๐^R(๐ฉ_f^{},ฯต).$$
(9)
Note that for spin-singlet pairing all vector functions get replaced by the corresponding scalar functions. In the weak-coupling spin-fluctuation model the pairing interaction is written as
$`V(๐ฉ_f,๐ฉ_f^{})V^{}(๐ฉ_f,๐ฉ_f^{})\chi (๐ฉ_f๐ฉ_f^{}),`$ (10)
$`\chi (๐ช)=\chi _0/\left[1+\xi ^2(๐ช๐ช_0)^2\right].`$ (11)
The detailed form of the effective pairing interaction $`V^{}(๐ฉ_f,๐ฉ_f^{})`$ depends on the form of the spin singlet or spin triplet pairing interaction. $`\chi _0`$ is the static spin susceptibility, $`\xi `$ is the antiferromagnetic correlation length, and the incommensurate wave vectors are $`๐ช_0(\pm 2\pi /3,\pm 2\pi /3)`$. The spin-fluctuation scenario proposed here is similar to the one studied by many authors in the context of the heavy-fermion systems, the high-$`T_c`$ cuprates, the quasi-two-dimensional organic superconductors, and even SrRuO. In contrast to the microscopic model calculations in Refs. , we propose the existence of either an attractive triplet $`f`$-wave or singlet $`d`$-wave pairing channel in order to describe the power-laws observed in thermodynamic and transport coefficients. Our approach is guided by neutron-scattering data of the spin susceptibility that can lead to a gap function that is gapped along the $`(\pm \pi ,\pm \pi )`$ directions and has nodes along $`(\pi ,0)`$ and $`(0,\pi )`$. The immediate consequence of the proposed state is that the superconducting gap on the holelike $`\beta `$ band develops nodes at $`(\pm 2\pi /3,\pm \pi )`$ and $`(\pm \pi ,\pm 2\pi /3)`$.
## IV Results And Discussions
In our analysis of the thermodynamic and transport properties we make the simplifying assumption that all three Fermi surfaces $`(\alpha ,\beta ,\gamma )`$ simultaneously go superconducting and can be described by one effective, cylindrical band. At the present time we cannot rule out any admixture of the $`p`$-wave state #2 to the $`f`$-wave state #3, since both gap functions belong to the same two-dimensional representation $`E_u`$. However, from a detailed analysis of the calculated heat capacity we find, rather conservatively, that the admixture of a nodeless $`p`$-wave state has to be less than 20% to be consistent with the experimental $`C(T)`$. Thus, we neglect the possibility of a $`p`$-wave admixture to the $`f`$-wave gap function in the remainder of this work. Impurity calculations for the $`p`$-wave state #2 also were performed by Maki and Puchkaryov, who reported reasonably good agreement between experiment and the calculated $`p`$-wave order parameter. Very recently, Dahm, Won, and Maki discarded the nodeless $`p`$-wave state and argued in favor of $`f`$-wave pairing.
In the temperature range $`T^{}TT_c`$, where $`T^{}`$ is the characteristic temperature of the impurity band width, and in the clean limit, $`\mathrm{\Gamma }\mathrm{\Delta }_0`$, the evaluation of the entropy and transport coefficients simplifies significantly. In the presence of line nodes on the Fermi surface the density of states is $`N(ฯต)(ฯต/\mathrm{\Delta }_0)N_f`$. Similarly, we obtain for the Fermi surface averaged integrand $`๐๐ฉ_f๐ฆ(๐ฉ_f,ฯต)ฯต\tau (ฯต)/\mathrm{\Delta }_0`$, because of $`๐ฆ\pi ^4/(\mathrm{Re}\stackrel{~}{ฯต}^R\mathrm{Im}\sigma _{imp}^R)\mathrm{Im}C^R`$, with $`\stackrel{~}{ฯต}^Rฯต+i0^+`$. Thus, the transport coefficients show the usual power laws of clean superconductors when using the approximate relations for the scattering self-energies, $`1/2\tau (ฯต)=(\pi )^1\mathrm{Im}\sigma _{imp}^R\mathrm{\Gamma }ฯต/\mathrm{\Delta }_0`$ in the Born limit, or $`1/2\tau (ฯต)\mathrm{\Gamma }\mathrm{\Delta }_0/ฯต`$ for unitarity scattering.
### A Specific heat
The states #1 and #3 with line nodes yield $`CN_fT^2/\mu \mathrm{\Delta }_0`$, in excellent agreement with experiments, while the gapped state #2 disagrees with the data. The proposed multiband order-parameter model by Agterberg and co-workers, which assumes that only one band $`(\gamma )`$ out of three possible bands goes superconducting at $`T_c`$, fails to describe the low-$`T`$ dependence (see Fig. 2). Our result for the $`p`$-wave state #2 is in agreement with calculations of the heat capacity by Agterberg. In the multiband model the density of states (DOS) of the $`\gamma `$ band is weighted with 57% of the total DOS, while the remaining $`\alpha `$ and $`\beta `$ bands account for 43% of the total DOS. It is the $`\gamma `$ band on which the $`p`$-wave state #2 has been proposed to nucleate. The $`\alpha `$ and $`\beta `$ bands remain normal. Here $`\mu `$ is the slope parameter of the gap function at the nodes, $`\mu =|d\mathrm{\Delta }(\varphi )/\mathrm{\Delta }_0d\varphi |_{node}`$. In our calculations we have used variational basis functions, $`๐ซ(๐ฉ_f)๐ซ(๐ฉ_f)_{A_{1g}}(๐ฉ_f;\mu )`$, where the variational function $`_{A_{1g}}`$ belongs to the $`A_{1g}`$ representation and remains invariant under all group transformations. The slope parameter $`\mu `$ allows us to adjust the opening of the gap function at the nodes, which is otherwise not determined by symmetry. This enables us to quantitatively describe the ground state of the superconducting order parameter as probed by low energetic quasiparticles. An approach that has been quite successful in describing the low energetic quasiparticle excitations in UPt<sub>3</sub>.
Assuming that pure SrRuO has an optimal transition temperature of $`T_{c0}1.51\mathrm{K}`$, we obtain an excellent fit for scattering in the Born limit with a scattering phase shift $`\delta _00`$ and a scattering rate $`\mathrm{\Gamma }/\pi T_{c0}=0.01`$. On the other hand, resonant scattering $`(\delta _0\pi /2)`$ with the same scattering rate gives a residual value of $`C/T`$ that is too large. If impurity scattering is indeed resonant, then a value of $`\mathrm{\Gamma }/\pi T_{c0}10^3`$ is required to account for the lowest measured values of the specific heat. Furthermore, it would imply that the optimal transition temperature is closer to $`T_{c0}1.48\mathrm{K}`$.
The $`T_c`$ transition of the two components of the triplet $`p`$-wave order parameter #2, or of the two components of the $`f`$-wave order parameter #3, is doubly degenerate. Similar to the multicomponent superconducting order parameter in UPt<sub>3</sub> uniaxial strain (pressure) in the plane would lift the degeneracy of the two-component order parameter. As a consequence the transition temperature will split into two. This is a crucial test of the multicomponent nature of the order parameter. Along the same line of arguments, a magnetic field in the plane should also split $`T_c`$, as was pointed out in Ref. .
### B Thermal conductivity
The in-plane thermal conductivity is isotropic for all order parameter models listed in Table I, assuming a cylindrical Fermi surface. In the clean limit, $`T^{}TT_c`$, and neglecting logarithmic corrections, $`\kappa _{ii}T`$ for weak scattering and $`T^3`$ for strong scattering. In the dirty limit, $`TT^{}T_c`$, the thermal conductivity is linear in temperature, $`\kappa _{ii}T`$, and independent of the scattering strength. Unfortunately the samples studied by Suderow et al. exhibit a very strong $`T_c`$ suppression. The reported resistive transitions for samples #2 and #4, $`T_c^\varrho (\mathrm{\#}2)0.81\mathrm{K}`$ and $`T_c^\varrho (\mathrm{\#}4)0.58\mathrm{K}`$, occurred significantly above the bulk superconducting transitions identified by the thermal conductivity, $`T_c^\kappa (\mathrm{\#}2)0.60\mathrm{K}`$ and $`T_c^\kappa (\mathrm{\#}4)0.47\mathrm{K}`$. Not only does this suggest that the samples are in the dirty limit but also that they are considerably inhomogeneous. Thus the standard scattering $`t`$-matrix analysis in terms of pointlike defects in the dilute limit will most likely fail to give a quantitative description. Nevertheless, combining the facts of the $`T_c`$ suppression and that the ratios of the residual resistivities and the normal-state thermal conductivities are related to the scattering rates, $`\mathrm{\Gamma }(\mathrm{\#}4)/\mathrm{\Gamma }(\mathrm{\#}2)\varrho _0(\mathrm{\#}4)/\varrho _0(\mathrm{\#}2)\kappa _N(\mathrm{\#}2)/\kappa _N(\mathrm{\#}4)1.25`$, we find that the normal-state scattering rates are approximately given by $`\mathrm{\Gamma }(\mathrm{\#}2)/\pi T_{c0}0.20`$ and $`\mathrm{\Gamma }(\mathrm{\#}4)/\pi T_{c0}0.25`$ (see Table II for the corresponding $`T_c`$ suppression).
In Figs. 3 and 4 we show the best fits of $`\kappa _{xx}`$ for samples #2 and #4 measured by Suderow et al.. Although we cannot obtain a quantitatively good fit for any of the pairing models, we are able to ascribe the large residual value of $`\kappa /T`$ to impurity scattering (see Fig. 3) without having to invoke a multiband order parameter model (see Fig. 4). A surprising result of these fits is that, generally, we find better agreement between theory and experiment for weak impurity scattering in the Born limit. Very recently, Tanatar et al. reported measurements of $`\kappa `$ on cleaner crystals ($`T_c1.4\mathrm{K}`$) that are in good quantitative agreement with the gapless states #1 or #3 and impurity scattering in the unitarity limit.
For the predicted pairing states #1 or #3, we expect to observe a fourfold oscillation of the thermal conductivity when a magnetic field is parallel to the layers and rotated within the layers. However, the amplitude of the oscillations depends on the scattering strength. It is appreciable for strong scattering (unitarity limit) and very small for weak scattering (Born limit). So far, no oscillations have been observed. Certainly the experimental and theoretical situation remains unresolved and requires more study. Indeed such magnetic oscillations have been reported in the cuprate YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub>, and are considered as additional proof in support of the $`d_{x^2y^2}`$ symmetry of the superconducting state.
### C Sound Attenuation
The longitudinal ($`๐ช๐บ[100]`$) and transverse ($`๐ช||[100]`$ and $`\epsilon ||[010]`$) sound attenuations are identical for the pairing state #2, i.e., $`\alpha _{xx}(T)/\alpha _{xx}(T_c)=\alpha _{xy}(T)/\alpha _{xy}(T_c)`$. This result also was reported in Ref. . Whereas for pairing states #1 and #3 the longitudinal attenuation with $`๐ช๐บ[100]`$ is the same as the transverse attenuation rotated by $`\pi /4`$ with $`๐ช||[110]`$ and $`๐บ||[\overline{1}10]`$. These relations follow directly from Eq. (7) and are a peculiarity of the 2D Fermi surface and the 2D basis functions of the order parameters. Inspecting the momentum dependent weighting factors in Eq. (7),
$$\pi _{xx}^2=\frac{1}{4}\mathrm{cos}^22\varphi ,\pi _{xy}^2=\frac{1}{4}\mathrm{sin}^22\varphi ,$$
(12)
it is clear that by rotating the crystal (or the transducer) by $`\pi /4`$ around the $`๐`$-axis one simply exchanges these functions, $`\pi _{xx}^2\pi _{xy}^2`$, and, thus swaps the expressions for the longitudinal and transverse attenuation. Since the integrand $`๐ฆ(๐ฉ_f,ฯต)`$ for the $`p`$-wave state (#2) is independent of $`๐ฉ_f`$, the longitudinal and transverse attenuations are identical (within an overall scaling factor due to differences in the speed of sound) for arbitrary temperature and impurity concentration. These predictions should be straightforward to check experimentally. In Fig. 5 we show the predicted transverse and longitudinal sound attenuations for the $`d`$-wave (#1) and $`f`$-wave (#3) order parameter models. Our results are similar to the ones discussed by Moreno and Coleman for the case of the $`d_{x^2y^2}`$-wave gap function in the high-$`T_c`$ cuprates.
## V Conclusions
We have proposed a spin-fluctuation model based on the measured spin susceptibility by neutron scattering that leads to nodes of the gap function on the Fermi surface. We demonstrated that the measured specific heat and thermal conductivity are consistent with a spin-singlet order parameter ($`d_{xy}`$-wave symmetry belonging to $`B_{2g}`$) or a spin-triplet order parameter ($`f`$-wave symmetry belonging to $`E_u`$), though inconsistent with a gapped spin-triplet state ($`p`$-wave symmetry belonging to $`E_u`$). Based on this analysis we proposed sound attenuation measurements and thermal conductivity measurements in a magnetic field to locate the nodes on the Fermi surface, as well as measurements of the specific heat subjected to a uniaxial strain field in the plane in order to split the superconducting transition. It is clear that more experiments are needed to investigate the nodal regions on the Fermi surface and the spin structure of the order parameter.
###### Acknowledgements.
We are indebted to J.A. Sauls and L. Taillefer for many insightful discussions and thank Y. Maeno, M. Sigrist, and D. Agterberg for discussions. We thank M. Tanatar and Y. Matsuda for sharing their data prior to publication. We acknowledge the Aspen Center for Physics for its hospitality. This work was supported by the Los Alamos National Laboratory under the auspices of the US Department of Energy.
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# On a theorem of Henri Cartan concerning the equivariant cohomology
## Introduction
The equivariant cohomology of a smooth manifold acted by a Lie group is a concept which crystallized in the works of A. Borel, H. Cartan, C. Chevalley, H. Hopf, L. Koszul and A. Weil in the late forties and early fifties.
The differential geometric approach to this subject was brilliantly described by Henri Cartan in the beautiful survey which continues to be the first source for anyone interested in learning the basic facts of this theory. Recently, it has been the focus of intense research in connection with many problems in differential geometry, representation theory and quantum field theory.
A central result of is Cartanโs theorem which states that if a compact Lie group $`G`$ acts freely on a smooth manifold $`M`$ then the $`G`$-equivariant cohomology of $`M`$ (as defined by Cartan) is naturally isomorphic to the DeRham cohomology of the quotient.
There are currently many proofs of this fact (e.g. , , ) but, in the authorโs view, they all suffer of the same รฆsthetic โdeficiencyโ. They involve a quite large amount of amazing combinatorics whose origin is somewhat obscure. Moreover, the resulting isomorphism is extremely difficult to figure out explicitely at cochain level.
The main goal of this paper is to provide a new, direct and more transparent proof of the following slight generalization of Cartanโs theorem.
Theorem If $`G`$ is a compact Lie group acting on the smooth manifold $`X`$ and $`N`$ is a closed normal subgroup of $`G`$ acting freely on $`X`$ then the $`G`$-equivariant cohomology of $`X`$ is naturally isomorphic with the $`G/N`$-equivariant cohomology of $`X/N`$.
We will actually establish a more general algebraic result (see Theorem 5.1). Moreover, relying on a recent result of Kalkman (which provides a very explicit isomorphism between the Weil model and the Cartan model of equivariant cohomology) we will offer an explicit description of this isomorphism (along the lines of ). In the course of the proof we will provide yet another interpretation for moment map and the equivariant characteristic classes described for example in or . The very simple functorial principle behind our proof is explained in Remark 5.4. While some of the computations involved may not look too eye pleasing, they are entirely routine and more importantly, their logical succesion is very natural.
There are two surprising aspects of this proof which make it so attractive. They can best be grasped by looking at a special example. Suppose $`PB`$ is a smooth principal $`G`$-bundle. Cartanโs theorem then states that $`H_G^{}(P)`$, the $`G`$-equivariant cohomology of $`P`$, coincides with the DeRham cohomology of the base $`B`$. Naturally, one tries to construct cochain homotopy equivalences between the complexes leading to the two cohomologies. A geometer might even attempt this using purely geometric operations on the smooth manifolds involved. This approach is doomed to fail. The method we propose is to embedd the two complexes in the same larger complex consisting of โidealโ elements and then show that the two embeddings are homotopic with this larger space. The homotopies are described by the Weil transgression between a genuine $`G`$-connection on $`P`$ and a certain โidealโ connection which has only a formal meaning !!! The second surprise is the amazing effectiveness of this method. Normally one expects that by โpushingโ the geometric situation into an โidealโ abstract framework the resulting formulรฆ will be more involved. To our surprize, Kalkmanโs isomorphism fits perfectly in such a framework. A bonus of this proof is that the isomorphism $`H_G^{}(P)H^{}(B)`$ can be explicitely described at the cochain level. More precisely, to any Cartan representative of an element in $`H_G^{}(P)`$ we associate a closed form on $`B`$. This correspondence descends to an isomorphism between the two cohomologies. This map is obtained naturally, as a by-product of our computations.
The paper is divided into five section of which four are devoted to surveying the basic โplayersโ in Cartanโs approach to equivariant theory. In the first part we introduce the notion of operation which captures the essential features of the DeRham algebra of a smooth manifold with a Lie group action.
In the first section we introduce the main object of study, that of operation. It formalizes the algebra of exterior forms on a smooth manifold equipped with a Lie group action. Section 2 introduces the Cartan and Weil models of equivariant cohomology while section 3 describes the Kalkman isomorphism between them. In section 4 we review the basics of the Weil transgression trick in the framework of operations. In the final section we prove Cartanโs theorem and discuss a few consequences.
In the sequel $`G`$ will always denote a compact, connected Lie group. The prefix โsโ will refer to โsuperโ (i.e. $`\text{}_2`$-graded) objects as in . The bracket $`[,]_s`$ denotes the super-commutator in a super-algebra. Also, we will use Einsteinโs summation convention (unless otherwise indicated).
## 1 Operations
As in we will consider Frechet algebras. These are associative -algebras such that their algebraic operations are continuous with respect to a Frechet topology. The standard example of Frechet algebra is that of the algebra of smooth functions on a smooth manifolds.
In this section we want to introduce the algebraic counterpart of the geometric notion of smooth manifold acted on by a Lie group. This object appears in literature with various names. We have chosen the terminology of which stays closer to the original motivation.
###### Definition 1.1
An operation consists of the following.
(a) A s-commutative -graded Frechet algebra
$$๐=_n\text{}๐^n$$
such that $`ab=(1)^{|a||b|}ba`$ for any two homogeneous elements. ($`๐`$ is naturally a s-algebra by $`๐=๐^{even}๐^{odd}`$).
(b) A continuous odd derivation
$$d:๐^{}๐^{+1}$$
such that $`d^2=0`$.
(c) A smooth action of the Lie group $`G`$ on $`๐`$ via algebra automorphisms commuting with $`d`$. We denote by $``$ the derivative of this action at $`\mathrm{๐}G`$. Thus $``$ defines a representation of g into the Lie algebra of even derivations of $`๐`$. ($`_X`$ $`(X\text{g})`$ is the Lie derivative along the automorphisms $`\mathrm{exp}(tX)`$ of $`๐`$.) Note that
$$[_X,d]_s=_Xdd_X=0.$$
(d) A continuous $`G`$-equivariant linear map $``$ from g (called contraction) to the space of odd derivations of $`๐`$ such that $`X,Y\text{g}`$
(d1) $`_X(๐^n)๐^{n1},n.`$
(d2) $`[_X,_Y]_s=_X_Y+_Y_X=0`$.
(d3) $`[_X,_Y]=_{[X,Y]}`$.
(d4) (Cartan formula) $`[_X,d]_s=_Xd+d_X=_X`$.
###### Remark 1.2
The above contraction $``$ extends to an algebra morphism $`:\mathrm{\Lambda }\text{g}\mathrm{End}(๐)`$ thus defining a $`\mathrm{\Lambda }\text{g}`$-module structure on $`๐`$.
###### Example 1.3
The right action of a Lie group $`G`$ on the smooth manifold $`M`$ defines a structure of operation on $`\mathrm{\Omega }^{}(M)`$. For each $`X\text{g}`$ we will denote by $`_X`$ the Lie derivative along the flow $`mm\mathrm{exp}(tX)`$ while $`_X`$ denotes the contraction along $`X^\mathrm{\#}`$-the infinitesimal generator of the above flow.
###### Example 1.4
(The Weil algebra) Consider the Lie group $`G`$ and set
$$W_G=S\text{g}^{}\mathrm{\Lambda }\text{g}^{}$$
where $`\mathrm{\Lambda }`$ and $`S`$ denote the exterior and respectively the symmetric algebra. Topologize $`W_G`$ in the obvious fashion (as a space of polynomials) and equip it with the -grading
$$\mathrm{deg}(\mathrm{\Lambda }^p\text{g}^{})=p,\mathrm{deg}(S^q\text{g}^{})=2q.$$
Denote by h the obvious isomorphism
$$\text{h}:\mathrm{\Lambda }^1\text{g}^{}S^1\text{g}^{}.$$
(1.1)
The usual derivations $`d`$, $`L_X`$ and $`i_X`$ on $`\mathrm{\Omega }^{}(G)`$ have an algebraic counterpart on $`\mathrm{\Lambda }\text{g}^{}`$ which we denote by $`๐`$, $`๐_X`$ and $`ฤฑ_X`$., It is convenient to describe these operations in โlocal coordinatesโ. Choose a basis $`(e_i)`$ of g and denote by $`(\theta ^i)`$ the dual basis of $`\text{g}^{}`$. Set $`\mathrm{\Omega }^i=\text{h}(\theta ^i)S^1\text{g}^{}`$. Denote by $`๐_i`$ and $`ฤฑ_j`$ the Lie derivative along $`e_i`$ and respectively the contraction by $`e_j`$. Then
$$๐_i\theta ^j=C_{ik}^j\theta ^k\mathrm{and}ฤฑ_j\theta ^k=\delta _j^k$$
where $`C_{jk}^i`$ denote the structural constants of the Lie algebra g
$$[e_j,e_k]=C_{jk}^ie_i.$$
To describe d in local coordinates we use the formula
$$(๐\omega )=\omega ([X,Y]),\omega \text{g}^{},X,Y\text{g}.$$
This yields
$$๐\theta ^i=\frac{1}{2}C_{jk}^i\theta ^j\theta ^k.$$
This implies immediately Koszulโs formula
$$๐=\frac{1}{2}\mu (\theta ^k)๐_k=\frac{1}{2}\theta ^k๐_k$$
(1.2)
where for any algebra $`A`$ and any $`aA`$ we denote by $`\mu (a)`$ the left multiplication by $`a`$. For simplicity we will very often omit the $`\mu `$ symbol when there is very little room for confusion.
The action of $`G`$ on $`W_G`$ induced by the coadjoint action defines a Lie derivative $`๐`$ extending the Lie derivative on $`\mathrm{\Lambda }\text{g}^{}`$ according to the prescription
$$๐_X\text{h}(\omega )=\text{h}(๐_X\omega ).$$
We can extend $`๐`$ to an odd derivation $`d_W`$ of $`W_G`$ uniquely determined by its action on the generators
$$d_W\theta ^i=๐\theta ^i+\text{h}\theta ^i=\frac{1}{2}C_{jk}^i\theta ^j\theta ^k+\mathrm{\Omega }^i=(\frac{1}{2}\theta ^j๐_j+\mathrm{\Omega }^jฤฑ_j)\theta ^i$$
(1.3)
$$d_W\mathrm{\Omega }^i=C_{jk}^i\theta ^j\mathrm{\Omega }^k=\left(\theta ^j๐_j\right)\mathrm{\Omega }^i.$$
(1.4)
Extend the contraction $`ฤฑ_X`$ to $`W_G`$ by imposing $`ฤฑ_X\text{h}(\omega )=0`$.
We leave the reader verify that these three derivations on $`W_G`$ do indeed define a structure of operation. $`\mathrm{}`$
Given two operations $`(๐_i,d_i,^i,^i)_{i=1,2}`$ we can define a structure of operation on their (Grothendiek) topological tensor product $`๐_1๐_2`$ (described e.g. in ) according to the rules $`(\epsilon _1`$ is the grading operator of the s-algebra $`๐_1`$)
$$d=d_1\mathrm{๐}+\epsilon _1d_2$$
$$=^1\mathrm{๐}+\mathrm{๐}^2$$
and
$$=^1\mathrm{๐}+\epsilon _1^2.$$
Given an operation $`(๐,d,,)`$ we can define three subalgebras
$$๐_{inv}=\mathrm{ker}=\{a๐;_Xa=0X\text{g}\}$$
$$๐_{hor}=\mathrm{ker}\mathrm{and}๐_{bas}=๐_{inv}๐_{hor}.$$
Since $`[d,]_s=0`$ we deduce $`d๐_{inv}๐_{inv}`$. Moreover, Cartan formula implies $`d๐_{bas}๐_{bas}`$. Thus we can define the cohomology groups
$$H_{inv}^{}(๐)=H^{}(๐_{inv},d)\mathrm{and}H_{bas}^{}(๐)=H^{}(๐_{bas},d).$$
###### Example 1.5
Consider a smooth principal $`G`$-bundle $`GPB`$. The right action of $`G`$ on $`P`$ induces a structure of operation on $`\mathrm{\Omega }^{}(P)`$. The basic subalgebra of this operation is then naturally isomorphic to $`\mathrm{\Omega }^{}(B)`$.
## 2 The Cartan-Weil descriptions
We will work in the more general setting of operations. We will define two notions of equivariant cohomology and in following subsection we will show they coincide.
Consider a $`G`$-operation $`(๐,d,,)`$.
The Weil description We define Weilโs equivariant cohomology of $`๐`$ by
$$WH_G^{}(๐)\stackrel{def}{=}H_{bas}^{}(W_G๐)$$
where $`W_G`$ denotes the Weil algebra introduced in the previous subsection.
The Cartan description Consider the algebra
$$=S\text{g}^{}๐.$$
$`S\text{g}^{}`$ is graded as usually by
$$\mathrm{deg}S^p\text{g}^{}=2p.$$
$`G`$ acts smoothly on $``$ and so we can form the subalgebra of invariant elements
$$_{inv}=(S\text{g}^{}๐)^G.$$
Now define the operator
$$\text{d}=\mathrm{๐}d\underset{i}{}\mu (\mathrm{\Omega }^k)_k=\mathrm{๐}d\mathrm{\Omega }^k_k.$$
(2.1)
If we regard $`\omega S\text{g}^{}๐`$ as a polynomial map $`\text{g}๐`$ then $`\text{d}\omega `$ is the polynomial map
$$Xd(\omega (X))_X(\omega (X)).$$
d satisfies the following conditions (see )
$$\text{d}^2=\mathrm{\Omega }^k_k.$$
(2.2)
$$[\text{d},๐\mathrm{๐}+\mathrm{๐}]=0.$$
(2.3)
The equality (2.3) shows that $`_{inv}`$ is d invariant. Moreover, on this subalgebra $`\text{d}^2=0`$. Indeed, on this subalgebra we have $`๐\mathrm{๐}=\mathrm{๐}`$ so that by (2.2) we have
$$\text{d}^2=\mathrm{\Omega }^i_i=\mathrm{\Omega }^i๐_i\mathrm{๐}.$$
Now it is not difficult to see that $`\mathrm{\Omega }^i๐_i0`$ on $`S\text{g}^{}`$ due to the skew symmetry of the structural constants. Thus $`(_{inv},\text{d})`$ is a cochain complex and we define the Cartan equivariant cohomology of $`๐`$ by
$$CH_G^{}(๐)\stackrel{def}{=}H^{}(_{inv},\text{d}).$$
When $`๐`$ is the algebra of differential forms on a smooth manifold $`M`$ on which $`G`$ acts smoothly we will use the notations $`WH_G^{}(M)`$ and $`CH_G^{}(M)`$ to denote the corresponding equivariant cohomologies.
## 3 Weil model $``$ Cartan model
Consider a $`G`$-operation $`(๐,d,,)`$. The main result of this subsection is the following.
###### Theorem 3.1
There exists a natural isomorphism
$$WH_G^{}(๐)CH_G^{}(๐).$$
We briefly describe the proof in . For a different but related approach we refer to .
Consider the algebra $`=W_G๐`$. It has a tensor product structure of $`G`$-operation with structural derivations $`D`$, $`L`$ and respectively $`I`$. For each $`U_i^j=\theta ^je_i\text{g}^{}\text{g}`$ define the following operators (on $``$)
$$\text{๐ธ}_i^j=\theta ^j_i$$
$$\text{๐}_i^j=\theta ^j_i\mathrm{\Omega }^j_i$$
In general for any $`T=t_j^iU_i^j\text{g}^{}\text{g}`$ set
$$\text{๐ธ}_T=t_j^i\text{๐ธ}_i^j,\text{๐}_T=t_j^i\text{๐}_i^j\mathrm{and}D_T=D+\text{๐}_T.$$
Note that $`(\text{๐ธ}_i^j)^2=0`$, $`i,j`$ and moreover $`\text{๐ธ}_T\text{๐ธ}_S=\text{๐ธ}_S\text{๐ธ}_T`$, $`S,T`$. Thus $`\mathrm{exp}(\text{๐ธ}_T)`$ is well defined and invertible. A simple computation shows that for all $`i,j`$ the operator $`\mathrm{exp}(\text{๐ธ}_i^j)=\mathrm{๐}+\text{๐ธ}_i^j`$ is an algebra automorphism of $``$ so that $`\mathrm{exp}(\text{๐ธ}_T)`$ is an automorphism of $``$ for all $`T\text{g}^{}\text{g}`$.
The key step in the proof of Theorem 3.1 is contained in the following result.
###### Lemma 3.2
For any $`T\text{g}^{}\text{g}`$ we have
$$\mathrm{exp}(\text{๐ธ}_T)D\mathrm{exp}(\text{๐ธ}_T)=D_T.$$
Proof of the lemma An elementary computation shows that for any $`i,j,k,\mathrm{}`$ we have the following โdifferential equationsโ
$$[D,\mathrm{exp}(\text{๐ธ}_j^i)]=\text{๐}_j^i\mathrm{exp}(\text{๐ธ}_j^i)$$
(3.1)
$$[\text{๐}_j^i,\mathrm{exp}(\text{๐ธ}_l^k)]=0.$$
(3.2)
The equality (3.1) can be rephrased as
$$\mathrm{exp}(\text{๐ธ}_j^i)D\mathrm{exp}(\text{๐ธ}_j^i)=D_{U_j^i}:=D+\text{๐}_j^i$$
(3.3)
while (3.2) is equivalent to
$$\mathrm{exp}(\text{๐ธ}_{\mathrm{}}^k)\text{๐}_j^i\mathrm{exp}(\text{๐ธ}_{\mathrm{}}^k)=\text{๐}_j^i.$$
(3.4)
Using (3.4) in (3.3) we deduce $`(\text{๐ธ}=\text{๐ธ}_j^i+\text{๐ธ}_{\mathrm{}}^k)`$
$$\mathrm{exp}(\text{๐ธ})D\mathrm{exp}(\text{๐ธ})=D_\text{๐ธ}.$$
Lemma 3.2 now follows by iterating the above procedure. $`\mathrm{}`$
We now have a whole family of $`G`$-structures on $``$ parameterized by $`\text{g}^{}\text{g}`$,
$$_T:=(,D_T,L_T=\mathrm{exp}(\text{๐ธ}_T)L\mathrm{exp}(\text{๐ธ}_T),I_T=\mathrm{exp}(\text{๐ธ}_T)I\mathrm{exp}(\text{๐ธ}_T)),T\text{g}^{}\text{g}.$$
Moreover an elementary computation shows that $`L_TL_0`$. All these structures are isomorphic with the canonical tensor product structure and in particular
$$H_{bas}^{}(_T)H_{bas}^{}(_0)WH^{}(๐).$$
An interesting special case arises when $`T=\mathrm{๐ข๐}=\theta ^ie_i`$. In this case the derivation $`D_{\mathrm{๐ข๐}}`$ is known as the BRST ($`=`$ Bechi-Rouet-Stora-Tyupin) operator and it arises in the quantization of classical gauge theories.
We leave the reader check that
$$(_{\mathrm{๐ข๐}})_{hor}=\mathrm{ker}I_{\mathrm{๐ข๐}}=S\text{g}^{}๐.$$
Hence
$$(_{\mathrm{๐ข๐}})_{bas}=(S\text{g}^{}๐)^G$$
and it is not difficult to see that on this subalgebra $`D_{\mathrm{๐ข๐}}=\text{d}`$. Thus
$$H_{bas}^{}(_{\mathrm{๐ข๐}})CH_G^{}(๐).$$
Theorem 3.1 is proved. $`\mathrm{}`$
###### Remark 3.3
Kalkmanโs isomorphism
$$\varphi =\varphi _G=\mathrm{exp}(\text{๐ธ}_{\mathrm{๐ข๐}}):_0_{\mathrm{๐ข๐}}$$
(3.5)
has a particularly nice form when restricted to $`(_0)_{bas}`$. It is uniquely determined by the correspondences $`\mathrm{\Omega }^i\mathrm{\Omega }^i`$ and $`\theta ^j0`$.
## 4 Algebraic connections
Among the possible actions of a Lie group on a smooth manifold the free ones play a special role. Consider for example the case of a smooth principal $`G`$-bundle $`GPB`$. Such actions admit connections. Recall (see ) that a connection on $`P`$ is an equivariant splitting
$$TP๐ฑPP$$
where $`๐ฑP`$ is the bundle spanned by the infinitesimal generators of the $`G`$ actions. In fact, for any $`pP`$ the correspondence
$$\text{g}XX_p^\mathrm{\#}$$
identifies the fiber $`๐ฑ_pP`$ with g.
Alternatively, a splitting as above can be defined by a vertical projector i.e. a g-valued 1-form $`\mathrm{\Theta }\text{g}\mathrm{\Omega }^1(P)`$ which is $`G`$-invariant and satisfies
$$i_{X^\mathrm{\#}}\mathrm{\Theta }=XX\text{g}.$$
(4.1)
We can regard this connection as a linear map
$$\mathrm{\Theta }:\text{g}^{}\theta ^i\stackrel{~}{\theta }^i$$
so that
$$\mathrm{\Theta }=e_i\stackrel{~}{\theta }^i.$$
The condition (4.1) reads
$$_j\stackrel{~}{\theta }^k=\delta _j^k$$
(4.2)
or equivalently,
$$_X\mathrm{\Theta }=X,X\text{g}.$$
(4.3)
The invariance implies
$$๐_ke_i\stackrel{~}{\theta }^i+e_i_k\stackrel{~}{\theta }^i=0$$
i.e.
$$_k\stackrel{~}{\theta }^i=C_{kj}^i\stackrel{~}{\theta }^j.$$
(4.4)
or equivalently,
$$\mathrm{\Theta }๐_X=_X\mathrm{\Theta },X\text{g}.$$
(4.5)
The conditions (4.2)-(4.5) are formulated using a language which involves only the structure of operation. Thus we can define an abstract notion of algebraic connection on any $`G`$-operation $`๐`$ as a $`G`$-equivariant linear map $`\stackrel{~}{\theta }:\text{g}^{}๐`$ satisfying (4.2)-(4.5).
###### Example 4.1
The inclusion $`\text{g}^{}W_G`$ defines an algebraic connection.
Consider now a $`G`$-operation $`(๐,d,,)`$ equipped with a connection
$$\mathrm{\Theta }:\theta ^i\stackrel{~}{\theta }^i.$$
Define
$$\stackrel{~}{\mathrm{\Omega }}^i=d\stackrel{~}{\theta }^i\frac{1}{2}C_{jk}^i\stackrel{~}{\theta }^j\stackrel{~}{\theta }^k.$$
The form
$$\stackrel{~}{\mathrm{\Omega }}=e_i\stackrel{~}{\mathrm{\Omega }}^i\text{g}๐^2$$
is independent of the basis $`(e_i)`$ and it is called the curvature of the connection. An elementary computation shows that
$$_k\stackrel{~}{\mathrm{\Omega }}^i=0i,k$$
i.e. $`\stackrel{~}{\mathrm{\Omega }}^i๐_{hor}^2`$, $`i`$.
The main algebraic implications of the existence of a connection derive from the following decomposition result.
###### Proposition 4.2
The connection induced map
$$\mathrm{\Lambda }\text{g}^{}๐_{hor}๐,\theta ^A\omega \stackrel{~}{\theta }^A\omega $$
($`\omega ๐_{hor}`$, $`A`$ is an ordered multi-index $`(a_1,a_2,\mathrm{})`$ and $`\theta ^A=\theta ^{a_1}\theta ^{a_2}\mathrm{}`$) is an isomorphism of graded algebras.
Idea of proof The map is clearly injective. The surjectivity follows from the following simple observation
$$\omega ๐,\omega \stackrel{~}{\theta }^k_k\omega \mathrm{ker}_k(\mathrm{no}\mathrm{summation}).\mathrm{}$$
Thus we can uniquely represent any element $`\omega ๐`$ as a polynomial
$$\omega =\stackrel{~}{\theta }^A\omega _A$$
where in the above sum $`A`$ runs through all ordered multi-indices. $`\omega `$ is said to be horizontally homogeneous if all the coefficients $`\omega _A๐_{hor}`$ have the same degree called the horizontal degree and denoted by $`\mathrm{deg}_h`$.
The component $`\omega _{\mathrm{}}๐_{hor}`$ of $`\omega ๐`$ is called the horizontal component, the map $`\omega \omega _{\mathrm{}}`$ will be denoted by $`h`$ and will be named the horizontal projection.
###### Remark 4.3
It is not difficult to see that the horizontal projection can be explicitly described by
$$h=\underset{k}{}\left(\mathrm{๐}\stackrel{~}{\theta }^k_k\right)=\mathrm{exp}(\stackrel{~}{\theta }^k_k).$$
We can now define the covariant derivative of the connection $`\mathrm{\Theta }`$ as the composition
$$=hd.$$
A simple computation shows that
$$\stackrel{~}{\theta }^i=\stackrel{~}{\mathrm{\Omega }}^i(\mathrm{Maurer}\mathrm{Cartan})$$
(4.6)
and
$$\stackrel{~}{\mathrm{\Omega }}^i=0(\mathrm{Bianchi})$$
(4.7)
Set
$$\text{cw}:W_G๐,\theta ^i\stackrel{~}{\theta }^i,\mathrm{\Omega }^i\stackrel{~}{\mathrm{\Omega }}^i.$$
This map is independent of the basis $`(e_i)`$ and it is called the Chern-Weil correspondence. The following result explains the universal role played by theโexoticโ structure of $`W_G`$.
###### Proposition 4.4
The Chern-Weil correspondence induced by a connection is a morphism of $`G`$-operations. Moreover, given two connections $`\stackrel{~}{\theta }_i,i=0,1`$ on $`๐`$ the corresponding Chern-Weil maps $`\text{cw}_i`$ are homotopic as morphisms of cochain complexes.
Proof The first part is left to the reader. To prove the second part we use a familiar trick from the theory of characteristic classes.
Form the algebra $`\widehat{๐}=\mathrm{\Omega }^{}(\text{})๐`$. (If $`๐`$ were the algebra of differential forms $`\mathrm{\Omega }^{}(M)`$ on a smooth manifold $`M`$ then $`\widehat{๐}\mathrm{\Omega }^{}(\text{}\times M)`$.) Clearly $`\widehat{๐}`$ is a $`G`$-operation and $`\widehat{\theta }=(1t)\stackrel{~}{\theta }_0+t\stackrel{~}{\theta }_1`$ defines a connection on $`\widehat{๐}`$. Denote by $`\mathrm{\Psi }_i`$ $`(i=0,1)`$ the the maps $`\mathrm{\Psi }_i:\widehat{๐}๐`$ defined by the localizations at $`t=1`$
$$\mathrm{\Omega }^{}(\text{})\text{},f(t)f(i),dt0.$$
We have a fiberwise integration morphism
$$_I:\widehat{๐}๐$$
defined by
$$_If(t)\omega =\{\begin{array}{ccc}\hfill 0& \mathrm{if}& f\mathrm{\Omega }^0(\text{})\\ \hfill \left(_0^1f(t)\right)\omega & \mathrm{if}& f\mathrm{\Omega }^1(\text{})\end{array}$$
The fundamental theorem of calculus implies immediately the following homotopy formula
$$\widehat{\omega }\widehat{๐}:\mathrm{\Psi }_1\widehat{\omega }\mathrm{\Psi }_0\widehat{\omega }=d_I\widehat{\omega }+_I\widehat{d}\widehat{\omega }$$
where $`\widehat{d}`$ is the exterior derivative in $`\widehat{๐}`$ defined by
$$\widehat{d}=dt\frac{}{t}\mathrm{๐}+ฯตd.$$
$`d`$ is the exterior derivative in $`๐`$ while $`ฯต`$ is the s-grading operator in $`\mathrm{\Omega }^{}(\text{})`$. From the equalities $`\text{cw}_i=\mathrm{\Psi }_i\widehat{\text{cw}}`$ we deduce Weilโs transgression formula
$$\text{cw}_1\text{cw}_0=d_I\widehat{cw}+_I\widehat{d}\widehat{\text{cw}}=d_I\widehat{\text{cw}}+_I\widehat{\text{cw}}d_W.$$
(4.8)
Thus the map
$$K=K(\stackrel{~}{\theta }_1,\stackrel{~}{\theta }_0)=_I\widehat{\text{cw}}:W_G^{}๐^1$$
is a cochain homotopy connecting $`\text{cw}_0`$ to $`\text{cw}_1`$. $`\mathrm{}`$
###### Remark 4.5
(a) It is instructive to compute $`K(\theta ^i)`$ and $`K(\mathrm{\Omega }^j)`$. We have
$$K(\theta ^i)=_I\widehat{\theta }^i=0.$$
To compute $`K(\mathrm{\Omega }^j)`$ we need to compute the curvature $`\widehat{\mathrm{\Omega }}`$ of $`\widehat{\theta }`$. Set $`\dot{\stackrel{~}{\theta }}=\stackrel{~}{\theta }_1\stackrel{~}{\theta }_0`$. We have
$$\widehat{\mathrm{\Omega }}=\widehat{d}\widehat{\theta }+\frac{1}{2}[\widehat{\theta },\widehat{\theta }]$$
$$=dt\dot{\stackrel{~}{\theta }}+d\stackrel{~}{\theta }_0+\frac{1}{2}[\stackrel{~}{\theta }_0,\stackrel{~}{\theta }_0]+td\dot{\stackrel{~}{\theta }}+t[\stackrel{~}{\theta }_0,\dot{\stackrel{~}{\theta }}]+\frac{t^2}{2}[\dot{\stackrel{~}{\theta }},\dot{\stackrel{~}{\theta }}]$$
$$=dt\dot{\stackrel{~}{\theta }}+\stackrel{~}{\mathrm{\Omega }}_0+t\left(d\stackrel{~}{\theta }_1+[\stackrel{~}{\theta }_1,\stackrel{~}{\theta }_0]\stackrel{~}{\mathrm{\Omega }}_0\frac{1}{2}[\stackrel{~}{\theta }_0,\stackrel{~}{\theta }_0]\right)+\frac{t^2}{2}[\dot{\stackrel{~}{\theta }},\dot{\stackrel{~}{\theta }}]$$
$$=dt\dot{\stackrel{~}{\theta }}+(1t)\stackrel{~}{\mathrm{\Omega }}_0+t\left(d\stackrel{~}{\theta }_1+[\stackrel{~}{\theta }_1,\stackrel{~}{\theta }_0]\frac{1}{2}[\stackrel{~}{\theta }_0,\stackrel{~}{\theta }_0]\right)+\frac{t^2}{2}[\dot{\stackrel{~}{\theta }},\dot{\stackrel{~}{\theta }}].$$
where $`\stackrel{~}{\mathrm{\Omega }}_0`$ is the curvature of $`\stackrel{~}{\theta }_0`$. Set
$$X(t)=X(t,\stackrel{~}{\theta }_1,\stackrel{~}{\theta }_0)=(1t)\stackrel{~}{\mathrm{\Omega }}_0+t\left(d\stackrel{~}{\theta }_1+[\stackrel{~}{\theta }_1,\stackrel{~}{\theta }_0]\frac{1}{2}[\stackrel{~}{\theta }_0,\stackrel{~}{\theta }_0]\right)+\frac{t^2}{2}[\dot{\stackrel{~}{\theta }},\dot{\stackrel{~}{\theta }}].$$
Thus
$$K(\mathrm{\Omega })=_I\left(X(t)+dt\dot{\stackrel{~}{\theta }}\right)=\stackrel{~}{\theta }_1\stackrel{~}{\theta }_0.$$
More generally if $`PS\text{g}^{}`$ and $`Q\mathrm{\Lambda }\text{g}^{}`$ we have
$$K(PQ)=_IP(\widehat{\mathrm{\Omega }})Q(\widehat{\theta })$$
$$=_I\left\{P(X(t)+dt\dot{\stackrel{~}{\theta }})Q(\widehat{\theta })\right\}.$$
Using the Taylor expansion of $`P`$ at $`X(t)`$ we get
$$=_0^1\left(\underset{i}{}(\stackrel{~}{\theta }_1^i\stackrel{~}{\theta }_0^i)\frac{P}{\mathrm{\Omega }^i}(X(t))Q(\widehat{\theta })\right)๐t.$$
(4.9)
In particular, this shows that (i) $`K`$ commutes with the $`G`$ action and (ii) $`K(W_{bas})๐_{bas}`$. We call such a homotopy a basic homotopy. We will use the symbol โ$`_b`$โ to denote the basic homotopy equivalence relation. In particular, we say that two $`G`$-operations $`๐,`$ are b-homotopic if there exist morphisms $`f:๐`$ and $`g:๐`$ such that $`fg_b\mathrm{๐ข๐}`$ and $`gf_b\mathrm{๐ข๐}`$. We write this as $`๐_b`$.
(b) The above proposition shows that we could define the notion of connection as a morphism of $`G`$-operations $`W_G๐`$. We see that $`W_G`$ is extremely rigid since for any $`G`$-operation the collection $`[W_G,๐]_b`$ of classes of morphisms of $`G`$-operations modulo basic homotopies is very small. It consists of at most one element. A $`G`$-operation $`๐ฒ`$ equipped with a $`G`$-connection satisfying the above rigidity condition (i.e. $`[๐ฒ,๐]_b`$ consists of at most one element for any $`G`$-operation $`๐`$) will be called an universal $`G`$-operation. Note also the similarity between this result and the topological one: two continuous maps $`f_1,f_2:BBG`$ which induce isomorphic principal $`G`$-bundles are homotopic.
(c) The proof of the above proposition continues to hold in the following more general form: any two equivariant morphisms $`\varphi _i:W_G๐`$ of graded differential algebras are homotopic as cochain maps. In particular, this shows that $`(W_G,d_W)`$ is acyclic i.e. $`H^k(W_G)=0`$ for $`k>0`$.
## 5 The basic cohomology of a $`G`$-operation with connection
As we explained in the previous subsection, the $`G`$-operations with connections represent the algebraic counterpart of a smooth manifold $`M`$ on which $`G`$ acts freely. In such a case, the (Borel) equivariant cohomology of $`M`$ is naturally isomorphic with the ordinary cohomology of the quotient
$$H_G^{}(M)H^{}(M/G).$$
In the subsection we will establish the algebraic counterpart of this result. In fact, we will deal with a more general situation.
Assume we are given the following collection of data.
$``$ A Lie group $`G`$ and a closed normal subgroup $`NG`$. Set $`Q=G/N`$. Since $`N`$ is invariant under the adjoint of $`G`$ there is an induced action on n-the Lie algebra of $`N`$ and in particular, n is a Lie algebra ideal of g.
$``$ A $`G`$-operation $`(๐,D,,)`$ equipped with a $`G`$-invariant $`N`$-connection i.e. a $`G`$-equivariant morphism of $`N`$-operations
$$\stackrel{~}{\theta }:W_N๐.$$
By regarding $`๐`$ as an $`N`$-operation we can form the subalgebra
$$=๐_{N,bas}=\{\omega ๐;\omega \mathrm{is}N\mathrm{invariant},_X\omega =0,X\text{n}\}.$$
The $`G`$-operation structure on $`๐`$ induces a residual $`G/N`$-operation structure on $``$. Note that we have an inclusion
$$j:W_QW_G๐$$
such that
$$j(W_Q)_{Q,bas}(W_G๐)_{G,bas}.$$
The geometric intuition behind this algebraic situation is that of a smooth $`G`$-space $`E`$ such that the action of $`N`$ is free. In Borel cohomology we have an isomorphism
$$H_G^{}(E)H_Q^{}(E/N).$$
We will establish the algebraic analogue of this result.
###### Theorem 5.1
(Cartan) The inclusion $`j`$ induces an isomorphism
$$WH_Q^{}()WH_G^{}(๐).$$
Proof Our proof will be a simple application of Weilโs transgression trick. For different approaches we refer to , Chap. VIII, , or .
We will construct a $`G`$-connection on $`W_G๐`$ starting from the $`G`$-equivariant $`N`$-connection $`\stackrel{~}{\theta }๐^1\text{n}`$.
Define the linear map
$$\mu :\text{g}๐^0\text{n},\mu (X)=_X\stackrel{~}{\theta }.$$
$`\mu `$ is called the moment map of the $`G`$-equivariant connection $`\stackrel{~}{\theta }`$. We can also regard it as an element of $`\text{g}^{}๐^0\text{n}`$.
###### Lemma 5.2
The moment map $`\mu :\text{g}๐^0\text{n}`$ is $`G`$-equivariant.
Proof Regard $`\stackrel{~}{\theta }`$ as a $`G`$-equivariant map
$$\text{g}^{}\text{n}^{}๐^1.$$
For each $`X\text{g}`$ regard $`\mu (X)`$ as a map
$$\text{g}^{}\text{n}^{}๐^0.$$
The equivariance of $`\mu `$ is equivalent to
$$\mu (Ad_gX)=g\mu (X)Ad_{g^1}^{}$$
where $`Ad^{}`$ denotes the coadjoint action of $`G`$. We have
$$\mu (Ad_gX)=_{Ad_gX}\stackrel{~}{\theta }=g_Xg^1\stackrel{~}{\theta }=g_X\stackrel{~}{\theta }Ad_{g^1}^{}=g\mu (X)Ad_{g^1}^{}.$$
(The second equality is the $`G`$-equivariance of $``$ while the third equality is the $`G`$-equivariance of $`\stackrel{~}{\theta }`$.) $`\mathrm{}`$
Define $`๐ช:\text{g}๐^0\text{n}`$ as
$$๐ช(X)=X+\mu (X).$$
Note that $`๐ช(X)=0`$ for $`X\text{n}`$ so that q descends to a map
$$๐ช=\text{q}=\text{g}/\text{n}๐^0\text{n}.$$
Set $`\mathrm{\Xi }=๐ช+\stackrel{~}{\theta }`$. Note that
$$๐ช\text{g}^{}๐^0\text{n}W_G^1๐^0\text{g}$$
and
$$\stackrel{~}{\theta }๐^1\text{n}๐^1\text{g}.$$
Thus $`\mathrm{\Xi }(W_G๐)^1`$.
###### Lemma 5.3
$`\mathrm{\Xi }`$ defines a $`G`$-connection on $`W_G๐`$
Proof For $`X\text{g}`$ denote by $`I_X`$ the contraction by $`X`$ in $`W_G๐`$. Then
$$I_X\mathrm{\Xi }=๐ช(X)+_X\stackrel{~}{\theta }=X.$$
The $`G`$-invariance of $`\mathrm{\Xi }`$ now follows from the $`G`$-invariance of $`๐ช`$ and $`\stackrel{~}{\theta }`$. $`\mathrm{}`$
The $`G`$-operation $`W_G๐`$ admits the tautological connection
$$\mathrm{๐}:W_GW_G๐,ww1.$$
Denote by $`K=K(\mathrm{๐},\mathrm{\Xi })`$ the Weil transgression
$$K:W_GW_G๐$$
so that for all $`wW_G`$
$$w\mathrm{\Xi }w=\delta Kw+Kd_Ww$$
where $`\delta `$ denotes the exterior derivation in $`W_G๐`$. Now define
$$T_0:W_G๐W_G๐,wa=(\mathrm{\Xi }w)a$$
(5.1)
$$T_1=\mathrm{๐ข๐}:W_G๐W_G๐$$
and
$$๐ฆ:W_G๐W_G๐,waKwa.$$
Then for all $`xW_G๐`$
$$xT_0x=T_1xT_0x=\delta ๐ฆx+๐ฆ\delta x.$$
Both $`T_0`$ and $`T_1`$ are morphisms of $`G`$-operations and $`๐ฆ`$ is a basic homotopy. Also note that
$$T_0(W_G๐)W_Q๐$$
and
$$T_0(W_G๐)_{G,bas}(W_Q)_{Q,bas}.$$
Moreover, along the basic subalgebras $`T_0j=\mathrm{๐ข๐}`$. In (the basic) cohomology $`T_0`$ is bijective since it is homotopic to the identity. This completes the proof of Theorem 5.1. $`\mathrm{}`$
###### Remark 5.4
The reduction theorem we have just proved generalizes as follows. Consider a $`G`$-operation $`๐ฒ`$. Then the transgression trick in the proof of Theorem 5.1 can be used to show the statements below are equivalent.
(i) $`๐ฒ`$ is universal (in the sense defined in Remark 4.5 (b)).
(ii) Any morphism of $`G`$-operations $`\phi :๐ฒ๐`$ induces a b-homotopy equivalence
$$๐ฒ๐_b๐.$$
Note in particular that if $`๐ฒ_0`$, $`๐ฒ_1`$ are two universal $`G`$-operations and $`\alpha :๐ฒ_0๐ฒ_1`$ is a morphism then any morphism $`\tau :๐ฒ_1๐ฒ_0`$ induces a b-homotopy equivalence
$$๐ฒ_0_b๐ฒ_1.$$
(5.2)
Indeed, by (ii) $`\tau `$ induces a b-homotopy equivalence
$$๐ฒ_0๐ฒ_1_b๐ฒ_1\left(๐ฒ_0\right)_b๐ฒ_0$$
(since $`๐ฒ_1`$ is universal) and on the other hand $`\alpha :๐ฒ_0๐ฒ_1`$ induces the b-homotopy equivalence
$$๐ฒ_0๐ฒ_1_b๐ฒ_1$$
since $`๐ฒ_0`$ is universal. If we take $`๐ฒ_0=W_{G/N}`$ and $`๐ฒ_1=W_G`$ (note that these Weil algebras are clearly universal $`G`$-operations) then the equivalence (5.2) is precisely the content of Theorem 5.1.
###### Corollary 5.5
Let $`E`$ be a smooth $`G`$-space. If $`N`$ acts freely on $`E`$ then
$$WH_G^{}(E)H_Q^{}(E/N).$$
It is instructive to describe the reduction isomorphism
$$WH_G^{}(E)\stackrel{}{}WH_Q^{}(E/N)$$
using the Cartan model. Denote by $`W_{\stackrel{~}{\theta }}`$ the Weil model description of the above reduction isomorphism (defined in (5.1)). We denote by $`C_{\stackrel{~}{\theta }}`$ its correspondent in the Cartan model. Denote by $`\varphi _G`$ (resp. $`\varphi _Q`$) the Kalkman isomorphism (cf. (3.5))
$$WH_G^{}CH_G^{}(\mathrm{resp}.WH_Q^{}CH_Q^{}).$$
We then have
$$C_{\stackrel{~}{\theta }}=\varphi _QW_{\stackrel{~}{\theta }}\varphi _G^1.$$
To get a better feeling on the structure of $`C_{\stackrel{~}{\theta }}`$ we will work in local coordinates. Choose a basis $`(e_i)`$ of n and then extend it to a basis $`\{e_i;f_a\}`$ of g. Via the natural projection $`\text{g}\text{g}/\text{n}`$ the collection $`(f_a)`$ induces a basis of $`\text{g}/\text{n}`$ which we continue to denote by the same symbols. Denote the dual basis of $`\{e_i;f_a\}`$ by $`\{\theta ^i;\phi ^a\}`$. We can regard $`(\theta ^i)`$ as a basis of $`\text{n}^{}`$ and $`(\phi ^a)`$ as a basis of $`\text{q}^{}`$. We denote the image of $`\theta ^i`$ in $`\text{n}^{}`$ by $`\mathrm{\Omega }^i`$ and the image of $`\phi ^a`$ in $`S\text{q}^{}`$ by $`\mathrm{\Psi }^a`$. Set $`\theta =\theta ^ie_i`$, $`\phi =\phi ^af_a`$, $`\mathrm{\Omega }=\mathrm{\Omega }^ie_i`$ and $`\mathrm{\Psi }=\mathrm{\Psi }^af_a`$.Then any element in $`(S\text{g}^{}\mathrm{\Omega }^{}(E))^G`$ is a $`G`$-equivariant polynomial map $`P:\text{g}\mathrm{\Omega }^{}(E)`$ which we will schematically describe it as $`P=P(\mathrm{\Omega }\mathrm{\Psi })`$. Then
$$\varphi _G^1P(\mathrm{\Omega }\mathrm{\Psi })=\mathrm{exp}(\phi ^a_a)\mathrm{exp}(\theta ^i_i)P(\mathrm{\Omega }\mathrm{\Psi }).$$
The map $`\mathrm{\Xi }:W_GW_G\mathrm{\Omega }^{}(E)`$ is determined from the assignments
$$\theta ^i\stackrel{~}{\theta }^i\stackrel{~}{\theta }^i(f_a)\phi ^a,\phi ^a\phi ^a.$$
If we define $`\mathrm{\Xi }(e_i)=e_i`$ and $`\mathrm{\Xi }(f_a)=f_a`$ then we can rewrite
$$\mathrm{\Xi }(\theta \phi )=\stackrel{~}{\theta }+\phi \stackrel{~}{\theta }(f_a)\phi ^a=\stackrel{~}{\theta }+\phi \stackrel{~}{\theta }\phi .$$
(5.3)
Moreover
$$\mathrm{\Xi }(\mathrm{\Omega }\mathrm{\Psi })=\delta \mathrm{\Xi }(\theta \phi )+\frac{1}{2}[\mathrm{\Xi }(\theta \phi ),\mathrm{\Xi }(\theta \phi )]$$
$$=d\stackrel{~}{\theta }+d_W\phi +\frac{1}{2}\left\{[\stackrel{~}{\theta }+\phi ,\stackrel{~}{\theta }+\phi ]+[\stackrel{~}{\theta }\phi ,\stackrel{~}{\theta }\phi ]\right\}[\stackrel{~}{\theta }+\phi ,\stackrel{~}{\theta }\phi ]\delta (\stackrel{~}{\theta }\phi )$$
$$=d\stackrel{~}{\theta }+\frac{1}{2}[\stackrel{~}{\theta },\stackrel{~}{\theta }]+d_W\phi +\frac{1}{2}[\phi ,\phi ]+[\stackrel{~}{\theta },\phi ]+\frac{1}{2}[\stackrel{~}{\theta }\phi ,\stackrel{~}{\theta }\phi ][\stackrel{~}{\theta }+\phi ,\stackrel{~}{\theta }\phi ]\delta (\stackrel{~}{\theta }\phi )$$
$$=\stackrel{~}{\mathrm{\Omega }}+\mathrm{\Psi }+[\stackrel{~}{\theta },\phi ]+\frac{1}{2}[\stackrel{~}{\theta }\phi ,\stackrel{~}{\theta }\phi ][\stackrel{~}{\theta }+\phi ,\stackrel{~}{\theta }\phi ]\delta (\stackrel{~}{\theta }\phi ).$$
(5.4)
On the other hand
$$W_{\stackrel{~}{\theta }}\varphi _G^1P(\mathrm{\Omega }\mathrm{\Psi })=\mathrm{exp}(\phi ^a_a)\mathrm{exp}(\mathrm{\Xi }(\theta ^j)_j)P(\mathrm{\Xi }(\mathrm{\Omega }\mathrm{\Psi })).$$
(5.5)
Since $`W_{\stackrel{~}{\theta }}\varphi _G^1P`$ is a $`Q`$-basic element of $`W_Q\mathrm{\Omega }^{}(E/N)`$ the action of $`\varphi _Q`$ on this element is determined according to Remark 3.3 by setting $`\phi =0`$ in (5.5). Using the equalities (5.3) and (5.4) we deduce
$$C_{\stackrel{~}{\theta }}P=\varphi _QW_{\stackrel{~}{\theta }}\varphi _G^1P=\mathrm{exp}(\stackrel{~}{\theta }^j_j)P(\stackrel{~}{\mathrm{\Omega }}+\mathrm{\Psi }\delta (\stackrel{~}{\theta }\phi )_{\phi =0}).$$
(5.6)
We need to understand the term $`\delta (\stackrel{~}{\theta }\phi )_{\phi =0}`$. We have
$$\delta (\stackrel{~}{\theta }^i(f_a)\phi ^a)=d(\stackrel{~}{\theta }^i(f_a))\phi ^a+\stackrel{~}{\theta }^id_W\phi ^a$$
$$=d(\stackrel{~}{\theta }^i(f_a))\phi ^a+\stackrel{~}{\theta }^i(f_a)\{\mathrm{\Psi }^a๐ฌ(\phi )\}$$
where $`๐ฌ`$ denotes a quadratic term in the $`\phi `$โs. Thus when setting $`\phi =0`$ we get
$$\delta (\stackrel{~}{\theta }^i(f_a)\phi ^a\phi )_{\phi =0}=\stackrel{~}{\theta }^i(f_a)\mathrm{\Psi }^a.$$
Symbolically
$$\delta (\stackrel{~}{\theta }\phi )_{\phi =0}=\stackrel{~}{\theta }(\mathrm{\Psi })=\mu (\mathrm{\Psi })(S^1\text{q}^{}\mathrm{\Omega }^0(E))\text{n}.$$
(Recall that $`\mu `$ denotes the moment map of the connection.) Note that the differential form components in $`\mu (\mathrm{\Psi })`$ are $`N`$-basic so they can be regarded as forms on the basis $`E/N`$. Substituting this back in (5.6) we get
$$C_{\stackrel{~}{\theta }}P=\mathrm{exp}(\stackrel{~}{\theta }_j)P(\mathrm{\Psi }+\mu (\mathrm{\Psi })+\stackrel{~}{\mathrm{\Omega }}).$$
The exponential factor is precisely the horizontal projection $`h_{\stackrel{~}{\theta }}:\mathrm{\Omega }^{}(E)\mathrm{\Omega }^{}(E/N)`$ defined by the $`N`$-connection $`\stackrel{~}{\theta }`$. On the other hand the term $`\stackrel{~}{\mathrm{\Omega }}+\mu (\mathrm{\Psi })(\mathrm{\Omega }^2(E/N)(S^1\text{q}^{}\mathrm{\Omega }^0(E/N))\text{n}`$ is already $`Q`$-basic. It is called the equivariant curvature of the connection $`\stackrel{~}{\theta }`$ and will be denoted by $`\stackrel{~}{\mathrm{\Omega }}_Q`$. Note that $`\stackrel{~}{\mathrm{\Omega }}_Q`$ is an element of degree 2 in $`\stackrel{~}{\mathrm{\Omega }}_Q(S\text{q}^{}\mathrm{\Omega }^{}(E/N))^Q\text{n}`$. Thus
$$(C_{\stackrel{~}{\theta }}P)(\mathrm{\Psi })=h_{\stackrel{~}{\theta }}P(\mathrm{\Psi }+\stackrel{~}{\mathrm{\Omega }}_Q).$$
(5.7)
We still need to give an accurate definition of the right-hand-side term above. For any $`X\text{g}`$ define $`P(X+\stackrel{~}{\mathrm{\Omega }}_Q)`$ imitating the Taylor expansion at $`X`$
$$P(X+\stackrel{~}{\mathrm{\Omega }}_Q)=\mathrm{exp}(\stackrel{~}{\mathrm{\Omega }}_Q^i_i)P(X)$$
where $`\stackrel{~}{\mathrm{\Omega }}_Q=\stackrel{~}{\mathrm{\Omega }}_Q^ie_i(S\text{q}^{}\mathrm{\Omega }^{}(E))^2\text{n}`$ while $`_i`$ denotes the partial derivative along the direction $`e_i\text{n}\text{g}`$. Note that $`P(X+\stackrel{~}{\mathrm{\Omega }}_Q)=P(X+\mu (X)+\stackrel{~}{\mathrm{\Omega }})=P(\stackrel{~}{\mathrm{\Omega }})`$ for all $`X\text{n}`$ so that $`P(X+\stackrel{~}{\mathrm{\Omega }}_Q)`$ descends to a well defined map $`\text{q}\mathrm{\Omega }^{}(E/N)`$. Thus the polynomial in the right-hand-side of (5.7) should be rather viewed as an $`\mathrm{\Omega }^{}(E/N)`$-valued polynomial on g which descends to a polynomial on $`\text{q}=\text{g}/\text{n}`$.
In particular, to any $`G`$-invariant polynomial $`PS\text{g}^{}`$ one can associate an equivariantly closed element
$$\mathrm{\Psi }P(\mathrm{\Psi }+\stackrel{~}{\mathrm{\Omega }}_Q)(S\text{q}^{}\mathrm{\Omega }^{}(E/N))^Q.$$
This element clearly depends on the connection but its equivariant cohomology class does not. It will be denoted by $`P(E)CH_Q^{}((E/N)`$ and will be called the equivariant characteristic class of $`EE/N`$ corresponding to $`P`$. Note that when $`G=N`$ this correspondence is none other than the traditional Chern-Weil construction of the characteristic classes of a principal $`G`$-bundle.
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# QCD analysis of the diffractive structure functions measured at HERA and factorisation breaking at Tevatron
## 1 Regge parametrization
The 1994 data are first investigated in the framework of a Regge phenomenological model . The 1994 data are subjected to a fit in which a single factorisable trajectory ($`IP`$) is exchanged such that:
$`F_2^{D(3)}(Q^2,\beta ,x_{IP})=f_{IP/p}(x_{IP})F_2^{IP}(Q^2,\beta ).`$ (1)
In this parameterization, $`F_2^{IP}`$ can be interpreted as the structure function of the pomeron . The value of $`F_2^{IP}`$ is treated as a free parameter at each point in $`\beta `$ and $`Q^2`$. The pomeron flux takes a Regge form with a linear trajectory $`\alpha _{IP}(t)=\alpha _{IP}(0)+\alpha _{IP}^{^{}}t`$, such that
$`f_{IP/p}(x_{IP})={\displaystyle _{t_{cut}}^{t_{min}}}{\displaystyle \frac{e^{B_{IP}t}}{x_{IP}^{2\alpha _{IP}(t)1}}}dt,`$ (2)
where $`|t_{min}|`$ is the minimum kinematically allowed value of $`|t|`$ and $`t_{cut}=1`$ GeV<sup>2</sup> is the limit of the measurement. The value of $`\alpha _{IP}(0)`$ is a free parameter and $`B_{IP}`$ and $`\alpha _{IP}^{^{}}`$ are taken from hadron-hadron data . The fit with a single trajectory does not give a good description of the data in the same way as it is observed at $`Q^2=0`$ that secondary trajectories in addition to the pomeron are required to describe diffractive $`ep`$ data.
A much better fit is obtained when both a leading ($`IP`$) and a sub-leading ($`IR`$) trajectory are considered in the same way as in formula (1), where the values of $`F_2^{IP}`$ and $`F_2^{IR}`$ are treated as free parameters at each point in $`\beta `$ and $`Q^2`$, $`\alpha _{IP}(0)`$ and $`\alpha _{IR}(0)`$ being two free parameters. The flux factor for the secondary trajectory takes the same form as equation (2), with $`B_{IR}`$, and $`\alpha _{IR}^{^{}}`$ again taken from hadron-hadron data . This fit yields to the following value of $`\alpha _{IP}(0)=1.203\pm 0.020(\mathrm{stat}.)\pm 0.013(\mathrm{syst}.)_{0.035}^{+0.030}(\mathrm{model})`$ and is significantly larger than values extracted from soft hadronic data ($`\alpha _{IP}1.08`$). The quality of the fit is similar if interference between the two trajectories is introduced.
## 2 QCD fits and the structure of the Pomeron
It has been suggested that the $`Q^2`$ evolution of the Pomeron structure function may be understood in terms of parton dynamics from perturbative QCD where parton densities are evolved according to DGLAP equations , using the GRV parametrization for $`F_2^{IR}`$ .
For the pomeron, a quark flavour singlet distribution ($`zS_q(z,Q^2)=u+\overline{u}+d+\overline{d}+s+\overline{s}`$) and a gluon distribution ($`zG(z,Q^2)`$) are parameterized in terms of coefficients $`C_j^{(S)}`$ and $`C_j^{(G)}`$ at $`Q_0^2=3`$ GeV<sup>2</sup> such that :
$`zS(z,Q^2=Q_0^2)\left[{\displaystyle \underset{j=1}{\overset{n}{}}}C_j^{(S)}P_j(2z1)\right]^2e^{\frac{a}{z1}}`$ (3)
$`zG(z,Q^2=Q_0^2)\left[{\displaystyle \underset{j=1}{\overset{n}{}}}C_j^{(G)}P_j(2z1)\right]^2e^{\frac{a}{z1}}`$ (4)
where $`z=x_{i/IP}`$ is the fractional momentum of the pomeron carried by the struck parton, $`P_j(\zeta )`$ is the $`j^{th}`$ member in a set of Chebyshev polynomials, which are chosen such that $`P_1=1`$, $`P_2=\zeta `$ and $`P_{j+1}(\zeta )=2\zeta P_j(\zeta )P_{j1}(\zeta )`$. Some details about the fits can be found in Reference .
A sum of $`n=3`$ orthonormal polynomials is used so that the input distributions are free to adopt a large range of forms for a given number of parameters. The exponential factor is needed to ensure a correct convergence close to $`z`$=1.
The trajectory intercepts are fixed to $`\alpha _{IP}=1.20`$ and $`\alpha _{IR}=0.62`$. Only data points of H1 with $`\beta 0.65`$, $`M_X>2`$ GeV and $`y0.45`$ are included in the fit in order to avoid large higher twist effects and the region that may be most strongly affected by a non zero value of $`R`$, the longitudinal to transverse cross-section ratio.
## 3 Results of the QCD fits
The resulting parton densities of the Pomeron are presented in figure 1. As it was noticed in the 1994 $`F_2^D`$ paper , we find two possible fits quoted here as fit 1 and fit 2. Each fit shows a large gluonic content. The quark contribution is quite similar for both fits, but the gluon distribution tends to be quite different at high values of $`z`$. This can be easily explained as no data above $`z=0.65`$ are included in the fits. Thus there is no constraint from the data at high $`z`$. The quark densities is on the contrary more constrained in this region with the DGLAP evolution. Both fits show similar $`\chi ^2`$ (the $`\chi ^2`$ per degree of freedom is about 1.2) <sup>1</sup><sup>1</sup>1Fit 2 is a bit disfavoured compared to fit 1 (its $`\chi ^2`$ by degree of freedom is 1.3 compared to 1.2 for fit 1) and is quite instable: changing a little the parameters modifies the gluon distribution at high $`z`$.. Adding the 1995 data points into the fits also allows to get a better constraint on initial parton densities at $`Q_0^2=3`$ GeV<sup>2</sup> compared to the fits performed with 1994 data points alone. For the gluon density presented in figure 1, we have determined that $`\frac{\delta G}{G}25\%`$ for $`z`$ below 0.6.
The result of the fit is presented in figure 2 together with the experimental values for 1994 data points ; we see on this figure the good agreement of the QCD prediction and the data points, which supports the validity of description of the Pomeron in terms of partons following a QCD dynamics.
We have also tried to extend the QCD fits to lower $`Q^2`$ (below $`3`$ GeV<sup>2</sup>) using the 1995 $`F_2^D`$ measurement. The $`\chi ^2`$ of the fit turns out to increase ($`\chi ^2/ndf=1.6`$, adding 35 low $`Q^2`$ points to the 171 points) . This can be illustrated in figure 2 of Reference where changes of slopes of scaling violations for $`Q^2`$ below and above $`3`$ GeV<sup>2</sup> can be seen. It may indicate that breaking of perturbative QCD has already occured in this region.
The idea would then to use these parton distributions and to compare with the measurements at Tevatron in order to study factorisation breaking. The roman pots which will be available in the D0 experiment at Run II will allow a direct comparison with the results obtained from the HERA parton distributions. It will be possible to know where factorisation breaking takes place at Tevatron, e.g. is it at low or high $`\beta `$?
## 4 Acknowledgments
The results described in the present contribution come from a fruitful collaboration with J.Bartels, H.Jung R.Peschanski and L.Schoeffel.
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# AharonovโBohm spectral features and coherence lengths in carbon nanotubes
## I Spectral properties of metallic and semiconducting CNs
Without a magnetic field and disorder, the electronic properties of nanotubes are known to be dependent on their chiral vector $`\stackrel{}{๐}_h=(n,m)`$, expressed in unit-vectors of the hexagonal lattice by $`|\stackrel{}{C}_h|=\sqrt{3}a_{\mathrm{C}\mathrm{C}}\sqrt{n^2+m^2+nm}`$ ($`a_{\mathrm{C}\mathrm{C}}=1.42`$ ร
). From a tight-binding description of the graphite $`\pi `$ bands, with only first-neighbor C-C interactions, the dispersion relations can be obtained by diagonalization of the Hamiltonian (with periodic boundary conditions), and for instance for the case of armchair $`(n,n)`$ nanotubes, we can write
$$\begin{array}{ccc}\epsilon _\pm (k_{},k_{})\hfill & =\hfill & \pm \gamma _0\{1\pm 4\mathrm{cos}\frac{k_{}a}{2}\mathrm{cos}\frac{\sqrt{3}k_{}a}{2}\hfill \\ & & +4\mathrm{cos}^2\frac{k_{}a}{2}\}^{1/2}\hfill \end{array}$$
(1)
where $`\gamma _0`$ is the energy overlap integral between carbon atoms, $`a=\sqrt{3}a_{\mathrm{C}\mathrm{C}}`$, is the graphite lattice constant, and $`k_{}`$ is the wavevector perpendicular to the nanotube axis $`k_{}=2\pi q/(\sqrt{3}na)`$ where $`q=1,\mathrm{},2n`$, giving the quantized values of the wavevector in the $`\stackrel{}{๐}_h`$ direction, whereas the wavevector $`k_{}`$ parallel to the nanotube axis is associated with the specification of the 1D Brillouin zone $`\pi /\sqrt{3}<k_{}a<\pi /\sqrt{3}`$ which defines the 1D-band dispersion. By developing dispersion relations around the $`\stackrel{}{K}`$-points (i.e., close to the Fermi energy), with small $`\stackrel{}{\delta k}=\stackrel{}{k}\stackrel{}{K}=(2\pi /|\stackrel{}{๐}_h|)(q\nu /3)\stackrel{}{k_{}}+\delta k_{}\stackrel{}{k_{}}`$ (with $`\stackrel{}{T}\stackrel{}{k_{}}=2\pi `$ and $`\stackrel{}{T}\stackrel{}{k_{}}=0`$, where $`\stackrel{}{T}`$ is the smallest translation vector along the tube axis), one finds
$$\epsilon _\pm (\delta k)=\pm \frac{\sqrt{3}\gamma _0}{2}\sqrt{\left\{\frac{2\pi a}{|\stackrel{}{๐}_h|}\left(q\frac{\nu }{3}\right)\right\}^2+\left|\delta k_{}\right|^2}$$
(2)
where the integer $`\nu `$ is related to $`m`$ and $`n`$ by $`nm=3p+\nu `$ where $`p=0,1,2,\mathrm{}`$ and $`\nu =0,\pm 1`$. Two bands may then cross at the Fermi level according to the value of $`\nu `$. If $`\nu =\pm 1`$ (i.e., $`nm=3p\pm 1`$), one gets $`\mathrm{\Delta }_g=\epsilon _+(\delta k)\epsilon _{}(\delta k)=2\pi a_{\mathrm{C}\mathrm{C}}\gamma _0/|\stackrel{}{๐}_h|`$ which defines the gap at the Fermi energy of a semiconducting nanotube. For $`\nu =0`$ (i.e., $`mn=3p`$), the system is metallic (in the sense $`\mathrm{\Delta }_g=0`$). Typically, the gap energy $`\mathrm{\Delta }_g`$ is, respectively, in the range $`1.65<\mathrm{\Delta }_g<0.27`$ eV for nanotube diameters $`d_t`$ in the range $`0.5<d_t<3`$nm. Predicted theoretically, these results have been confirmed experimentally by scanning tunneling spectroscopy (STM) measurements.
## II Magnetic field induced MIT, splitting and shifting of VHSs
To investigate AharonovโBohm phenomena, we start from the Hamiltonian $`_{\mathrm{๐ค๐ค}^{}}`$ for electrons moving on a nanotube under the influence of a magnetic field :
$$\begin{array}{cc}\hfill _{\mathrm{๐ค๐ค}^{}}=& \frac{1}{N}\underset{๐,๐^{}}{}e^{i(๐ค.๐๐ค^{}.๐^{}){\displaystyle \frac{ie}{\mathrm{}}}\mathrm{\Delta }\phi _{R,R^{}}}\hfill \\ \hfill \times & \psi (๐ซ๐)\left|\frac{๐ฉ^2}{2m}+V\right|\psi (๐ซ๐^{})\hfill \end{array}$$
(3)
where the phase factors are given by
$$\mathrm{\Delta }\phi _{๐,๐^{}}=_๐^๐^{}๐(\xi )๐\xi ,$$
(4)
and $`|\psi (๐ซ๐^{})`$ is the localized atomic orbital, and $`๐ฉ`$ and $`V`$ are, respectively, the momentum and disorder operators. As shown hereafter, different physics is found according to the orientation of the magnetic field with respect to the nanotube axis. In the former case, the vector potential is simply expressed as $`๐=(\varphi /|๐_h|,0)`$ in the two-dimensional $`\stackrel{}{๐}_h/|\stackrel{}{๐}_h|`$, $`\stackrel{}{T}/|\stackrel{}{T}|`$ coordinate system, and the phase factors become $`\mathrm{\Delta }\phi _{๐,๐^{}}=i(๐ณ๐ณ^{})\varphi /|๐_h|`$ for $`๐=(๐ณ,๐ด)`$. This yields new magnetic field-dependent dispersion relations $`\epsilon (\delta k,\varphi /\varphi _0)`$. Close to the Fermi energy, this energy dispersion relation is affected according to $`k_{}k_{}+2\pi \varphi /(\varphi _0|๐_h|)`$ which leads to a $`\varphi _0`$-periodic variation of the energy gap $`\mathrm{\Delta }_g`$. Such patterns are reported in Fig. 1 (top) where the total density of states (TDOS) at the Fermi level is plotted as a function of magnetic field strength $`\varphi /\varphi _0`$ threading the nanotube for the metallic (9,0) and the semiconducting (10,0) nanotubes. Note that a finite DOS is found in Fig. 1, for semiconducting and metallic nanotubes as a function of magnetic field, since we consider that the Greenโs function has a finite imaginary part, that we use to calculate the DOS by a recursion method .
In the semiconducting case, the oscillations in the DOS correspond to the following variations of the gap widths :
$$\mathrm{\Delta }_g=\{\begin{array}{ccc}\mathrm{\Delta }_0\left|13\frac{\varphi }{\varphi _0}\right|\hfill & \text{if}\hfill & 0\varphi \frac{\varphi _0}{2}\hfill \\ \mathrm{\Delta }_0\left|23\frac{\varphi }{\varphi _0}\right|\hfill & \text{if}\hfill & \frac{\varphi _0}{2}\varphi \varphi _0\hfill \end{array}$$
(5)
where $`\mathrm{\Delta }_0=2\pi a_{\mathrm{C}\mathrm{C}}\gamma _0/|\stackrel{}{๐}_h|`$ is a characteristic energy associated with the nanotube. It turns out that at $`\varphi `$ values of $`\varphi _0/3`$ and $`2/3\varphi _0`$, in accordance with the values of $`\nu =\pm 1`$, there is a local gap-closing in the vicinity of either the $`K`$ or $`K^{}`$ points in the Brillouin zone. This can be seen simply by considering the coefficients of the general wavefunction in the vicinity of the $`K`$ and $`K^{}`$ points, which can be written as $`\mathrm{\Psi }_{\stackrel{}{K}+\delta \stackrel{}{k}_K}(\stackrel{}{r}+\stackrel{}{๐}_h)`$. Since periodic boundary conditions apply in the $`\stackrel{}{๐}_h`$ direction, one can write $`\text{exp}[i(\stackrel{}{K}+\delta \stackrel{}{k}_K)\stackrel{}{๐}_h]=1`$. For $`\nu =+1`$, we write $`\delta \stackrel{}{k}_K=(2\pi /|๐_h|)(q1/3)\stackrel{}{k_{}}+\delta k_{}\stackrel{}{k_{}}`$ whereas $`\delta \stackrel{}{k}_K^{}=(2\pi /|๐_h|)(q+1/3)\stackrel{}{k_{}}+\delta k_{}\stackrel{}{k_{}}`$. When $`KK^{}`$, then $`\pm 1/31/3`$ in the above expressions, as $`\nu `$ goes from $`+11`$, which makes the situation between $`K`$ and $`K^{}`$ symmetrical. Similarly for metallic nanotubes, the gap-width $`\mathrm{\Delta }_g`$ is expressed by
$$\mathrm{\Delta }_g=\{\begin{array}{ccc}3\mathrm{\Delta }_0\frac{\varphi }{\varphi _0}& \text{if}& 0\varphi \frac{\varphi _0}{2}\\ 3\mathrm{\Delta }_0\left|1\frac{\varphi }{\varphi _0}\right|& \text{if}& \frac{\varphi _0}{2}\varphi \varphi _0\end{array}$$
(6)
Here we show that the magnetic field also affects the positions of the VHSs over the entire spectrum, which may be experimentally observed. As an illustration, the TDOS for a (10,10) tube without a magnetic field, shown as curve (a) in Fig. 2, is compared to the magnetic field case with $`\varphi /\varphi _0=0.125,0.25,0.375`$ and $`0.5`$ \[shown in Fig. 2 as curves (b), (c), (d) and (e), respectively\] for energies up to $`3`$ eV in which we have taken $`\gamma _0=2.9`$ eV, which is suitable for comparison with experiments. To understand such effects, one calculates the TDOS from
$$\text{Tr}[\delta (E)]=\frac{2}{2\pi }\underset{n}{}๐k\delta (E\epsilon _{nk}),$$
(7)
where the sum is taken over the $`n`$ energy bands $`\epsilon _{nk}`$ and we use
$$\left|\frac{\epsilon (\delta k,\frac{\varphi }{\varphi _0})}{k_{}}\right|^1=\frac{2}{\sqrt{3}\gamma _0a}\frac{\left|\epsilon (\delta k,\frac{\varphi }{\varphi _0})\right|}{\sqrt{\epsilon ^2(\delta k,\frac{\varphi }{\varphi _0})\epsilon _q^2}},$$
(8)
where $`\epsilon _q`$ indicates the position (energy) of the VHSs. We can rewrite the DOS as
$$\rho (E)=\frac{2a}{\pi \gamma _0\left|\stackrel{}{๐}_h\right|}\underset{q=1}{\overset{2n}{}}\delta _q(E,\epsilon _q),$$
(9)
where $`\delta _q(E,\epsilon _q)=|E|/\sqrt{E^2\epsilon _q^2}`$ for $`|E|>|\epsilon _q|`$ and zero otherwise.
The $`\epsilon _q`$ denote the energy-positions of the VHSs, and for armchair nanotubes $`(n,n)`$, one finds that $`\epsilon _q=\pm \gamma _0\sqrt{1+3\mathrm{cos}^2q\pi /n}`$ and $`\epsilon _q=\pm \gamma _0\sqrt{1\mathrm{cos}^2q\pi /n}`$ define the whole set of VHSs ($`q=1,\mathrm{},5`$). The magnetic field induces a shift of $`k_q`$ by a factor $`2\pi \varphi /(\varphi _0|\stackrel{}{๐}_h|)`$, which results in a new expression for the quantized values of the wavevector in the $`\stackrel{}{๐}_h`$ direction, which read $`2\pi (q+\varphi /\varphi _0)/|\stackrel{}{๐}_h|`$. For the metallic (10,10) nanotube, an energy gap thus opens at the Fermi level, whose width $`\mathrm{\Delta }_g`$ is proportional to the magnetic strength for $`\varphi /\varphi _01/2`$. Given that $`\epsilon _q`$ is a function of $`k_q`$, a shift of the energy positions of the VHSs follows, as well as the breaking of the degeneracy of a pair of $`K`$ and $`K^{}`$ points contributing to a given VHS, as explained below. In the (10,10) nanotube, the first five VHSs (which each have a degeneracy of 4) are simply given by $`E_{\mathrm{VHS}}=\gamma _0\mathrm{sin}(\pi q/10),(q=1,\mathrm{},5)`$ which leads to $`E_{\mathrm{VHS}}^{q=1}=0.896`$ eV, $`E_{\mathrm{VHS}}^{q=2}=1.705`$ eV, $`E_{\mathrm{VHS}}^{q=3}=2.346`$ eV, $`E_{\mathrm{VHS}}^{q=4}=2.758`$ eV, and $`E_{\mathrm{VHS}}^{q=5}=2.900`$ eV. Then phase shifts induced by the magnetic field correspond simply to
$$E_{\mathrm{VHS}}^q=\gamma _0\mathrm{sin}\left\{\frac{\pi }{10}\left(q\pm \frac{\varphi }{\varphi _0}\right)\right\}.$$
(10)
This is illustrated in Fig. 2, which gives the evolution of the TDOS of a (10,10) nanotube close to some VHS.
At low field, the degeneracy of the $`K`$ and $`K^{}`$ points is broken as illustrated in Fig. 3, which shows the Brillouin zone for 2D graphite along with the equi-energy contours around $`K`$ and $`K^{}`$ points, so that in addition to the upshift of one component of the degenerate zero-field VHSs, the other magnetically-split component is correspondingly downshifted. Both VHSs move apart up to $`\varphi /\varphi _0=1/2`$, where the two originally different VHSs merge into one. The positions of such merged VHSs are related to a $`\mathrm{cos}k_qa/2`$ factor (that is, through a $`\mathrm{cos}2\pi \varphi /\varphi _0`$ factor), so that a further increase of the magnetic field strength yields an opposite shifting of the VHS positions, and at $`\varphi /\varphi _0=1`$, the initial positions of the VHSs are recovered, along with their degeneracy (those given at $`\varphi /\varphi _0=0`$), and a gap closing thus results. At low magnetic field since $`E_{\mathrm{VHS}}^q\varphi /\varphi _0`$, the shift of VHSs position for a given magnetic field is proportional to $`\varphi /\varphi _0`$ for all VHSs. In the inset of Fig. 2, the positions of the VHSs originally at $`E=0.896`$ eV (for zero field), is given as a function of the normalized magnetic flux. The splitting is illustrated and a strong shift of the lowest VHS, by $`0.4`$ eV is predicted at half a quantum flux unit. Note that the upshifted and downshifted components are symmetrical with respect to the $`\varphi /\varphi _0`$-axis.
The TDOS of the semiconducting $`(18,2)`$ nanotube is also investigated and displays similar features, as exemplified by Fig. 4. The inset gives the shifting of the energy position of the three VHSs, closest to the Fermi level as a function of normalized magnetic flux. An interesting $`\varphi _0/2`$ symmetry is clearly seen for semiconducting tubes, as in the former case (Fig. 2) for metallic nanotubes. These interesting effects of the magnetic field on the spectral structure has, to our knowledge, not been explicitly investigated experimentally up to now, but recent progress in the observation of room temperature resonant Raman spectra for purified single-wall nanotubes can disclose fine structure in the DOS, through the enhancement of the intensity of optical absorption at the energies associated with VHSs. Magneto-optical techniques could then provide a possible way to test the above predictions. To that end, we give here the predicted energy position shifts for the VHSs closest to the Fermi energy, which are related to the phase shift induced by magnetic flux variations from 0 to $`\varphi _0/2`$ for nanotubes with diameters $`d_1=1.8`$ nm (the largest SWNT made up to now), $`d_2=6`$ nm and $`d_3=8`$ nm (corresponding to typical MWNT outer diameters) and corresponding to $`(n,n)`$ armchair tubes. The predicted energy position shift is given by
$$\begin{array}{cc}\hfill \mathrm{\Delta }E_{\mathrm{VHS}}(0,\frac{1}{2})=& E_{\mathrm{VHS}}^{q=1}\left(\frac{\varphi _0}{2}\right)E_{\mathrm{VHS}}^{q=1}(0)\hfill \\ \hfill =& \gamma _0\left(\mathrm{sin}\frac{3\pi }{2n}\mathrm{sin}\frac{\pi }{n}\right),\hfill \end{array}$$
(11)
which implies $`\mathrm{\Delta }E_{VHS,1}=173`$ meV, $`\mathrm{\Delta }E_{VHS,2}=52.3`$ meV and $`\mathrm{\Delta }E_{VHS,3}=39.2`$ meV, and the corresponding magnetic field strengths at $`\varphi _0/2`$ are $`B_1200`$ T $`B_218`$ T, $`B_310`$ T, which are within the scope of present experimental capabilities (we take $`B=(2\pi \varphi _0)/3(na)^2`$, and $`\varphi _0=4.1356\times 10^{15}`$ Tm<sup>2</sup>). Note that the effect of the magnetic field should be observed for individual MWNTs, since the inter-tube coupling for MWNTs is believed to affect the electronic spectrum weakly. Whereas the splitting is universal for all tubes (chiral and achiral), in the most general case, the van Hove singularity shift will depend on the chirality, and on the applied magnetic field strength (not the position of the van Hove singularity). To estimate the corresponding shift, one may proceed in the following manner: given a chiral vector $`\stackrel{}{๐}_h=(n,m)`$, one calculates the associated vectors $`\stackrel{}{k}_{}`$ and $`\stackrel{}{k}_{}`$ which are drawn in the Brillouin zone.
This is here illustrated for the $`(10,10)`$ armchair for which $`\stackrel{}{k}_{}=(t_2\stackrel{}{b}_1+t_1\stackrel{}{b}_2)/N_{(10,10)}=(\stackrel{}{b}_1+\stackrel{}{b}_2)/N_{(10,10)}`$ with $`N_{(10,10)}=40`$, the number of hexagons within the unit cell of a (10,10) nanotube, and $`\stackrel{}{b}_i`$ the basis vector in reciprocal space (see Ref. ), where the direction $`\stackrel{}{k}_{}`$ is found to be perpendicular to the $`KK^{}`$-axis (Fig.3), and the spacing between lines is equidistant to a given $`K`$-point. For the (18,2) chiral nanotube with $`N_{(18,2)}=364`$, we find $`\stackrel{}{k}_{}=(11\stackrel{}{b}_1+19\stackrel{}{b}_2)/N_{(18,2)}`$ which defines another direction and spacing. No lines cross at $`K`$ and $`K^{}`$ points, and the spacings between two consecutive lines are not equidistant to the $`K`$-location (in fact the $`K`$ point always appears to be one-third of the distance between the two lines near the Fermi energy). This affects the splitting, for instance at the field corresponding to half a quantum flux unit, the pattern for the $`(18,2)`$ nanotube is much less simple than that for the $`(10,10)`$ nanotube but it can be evaluated systematically.
PERPENDICULAR CASE
For a magnetic field perpendicular to the nanotube axis, the situation is more cumbersome, due to a site-position dependence of the vector potential. No apparent symmetries for the AharonovโBohm interferences on the nanotube are found in this case, and indeed, even if semiconducting nanotubes can become metallic with increasing magnetic field strength, non-periodic-oscillations are found, as described below. For $`\stackrel{}{B}`$ normal to the tube axis, one starts from a vector potential, given in the two-dimensional coordinate system $`(\stackrel{}{๐}_h,\stackrel{}{T}`$), by
$$๐=(0,\frac{B\left|\stackrel{}{๐}_h\right|}{2\pi }\mathrm{sin}\frac{2\pi }{\left|\stackrel{}{๐}_h\right|}๐ณ,0).$$
(12)
The effect of the magnetic field is driven by the phase factors introduced into the hopping integrals between two sites $`๐_๐ข`$ and $`๐_๐ฃ`$ (with $`๐_๐ข=(๐ณ_i,๐ด_i)`$, and the phase factor can be deduced from the Peierls substitution as follows, $`\mathrm{\Delta }\phi _{๐,๐^{}}`$:
$$\mathrm{\Delta }\phi _{๐,๐^{}}=\{\begin{array}{cc}\multicolumn{2}{c}{\left(\frac{|\stackrel{}{๐}_h|}{2\pi }\right)^2B\frac{\mathrm{\Delta }๐ณ}{\mathrm{\Delta }๐ด}[\mathrm{cos}\frac{2\pi }{|\stackrel{}{๐}_h|}๐ณ}\\ \hfill \mathrm{cos}\frac{2\pi }{|\stackrel{}{๐}_h|}(๐ณ+\mathrm{\Delta }๐ณ)]& \hfill (\mathrm{\Delta }๐ณ0)\\ \frac{|\stackrel{}{๐}_h|}{2\pi }B\mathrm{\Delta }๐ด\mathrm{sin}\frac{2\pi }{|\stackrel{}{๐}_h|}๐ณ\hfill & \hfill (\mathrm{\Delta }๐ณ=0)\end{array}$$
(13)
where $`\mathrm{\Delta }๐ณ=๐ณ_i๐ณ_j`$ and $`\mathrm{\Delta }๐ด=๐ด_i๐ด_j`$. In Fig. 1 (b), we show the total density of states (TDOS) at the Fermi level ($`E_F=0`$) as a function of the effective magnetic field defined by $`\nu =(L/2\pi \mathrm{}_m)`$, where $`\mathrm{}_m=\sqrt{\mathrm{}c/eB}`$ is the magnetic length. At low fields, the TDOS of metallic-(9,0) and semiconducting-(10,0) nanotubes at the Fermi level increases with the magnetic field strength. For higher values of magnetic field, our results are in agreement with previous results obtained by exact diagonalization. Also Landau bands are generated for values for which $`\nu 2`$. The aperiodic fluctuations of the TDOS are stronger at higher fields, with occasional low values of the DOS, reminiscent of a non-zero DOS for the semiconducting nanotubes at the zero-field value.
In Fig. 5, several values of $`\nu `$ are considered for an initially semiconducting nanotube. For $`\nu =1`$ the radius of the nanotube equals the magnetic length. Landau levels emerge whenever the magnetic length becomes smaller than the nanotube circumference length. Comparison of the case of $`\nu =3.5`$ in Fig. 5 with the zero-field limit is instructive, since the VHS partition of the spectra has been totally replaced by a Landau level spectrum (Here, one can recognize square-root singularities for the VHSs, and a Lorentzian-shape singularity for the Landau levels). This transition from the VHS-pattern to the Landau-level pattern is, however, more unlikely to be observed experimentally, since its observation requires a very high magnetic field.
## III Disorder and magnetic field effects
The effect of Anderson-type disorder on the electronic properties in addition to the presence of a magnetic field is now addressed in order to investigate how disorder alters the VHS-pattern and the metal-insulator transition (MIT), and how disorder qualitatively modifies the localization properties of the nanotubes when the system remains metallic. We consider here only the case of the metallic zigzag nanotube (9,0), but similar results are obtained for other metallic or semiconducting nanotubes. In zero magnetic field, this problem of localization in nanotubes has recently attracted a great deal of attention.
The effect of randomness is considered by taking the site energies of the tight-binding Hamiltonian at random in the interval $`[W/2,W/2]`$, with a uniform probability distribution. Accordingly, the strength of the disorder is measured by $`W`$. Significant fluctuations in the local density of states (LDOS) are found as shown in Fig. 6 (for $`\varphi /\varphi _0=1/2`$). The TDoS for $`W=0.25`$ up to $`W=2`$ ($`1/3`$ of the total bandwidth) shows that disorder does not modify the gap at the Fermi level, even when the confinement effects disappear (vanishing of VHSs). Disorder obviously leads to a mixing of energy levels which results in a vanishing of the VHS at $`W=1`$ in our simulations. The gap is more resistant to disorder, even when its strength is as large as $`1/3`$ the total bandwidth. Thus magnetic effects are not affected by low levels of disorder, which is an interesting result, since it is believed that many defects or sources of weak scattering should be present in real nanotubes. The estimation of the electron mean free path ($`\mathrm{}_e`$) in nanotubes is relevant for determining which is the most likely transport regime (ballistic versus diffusive), and further tells us whether or not the conductance should be quantized. While some experiments have observed conductance quantization, others suggest, in contrast, rather short electronic coherence lengths. A calculation of $`\mathrm{}_e`$ for armchair nanotubes has been performed using the special symmetries of these systems. Here, we show that a simple recourse to the relaxation time approximation (RTA) in the limit of weak scattering is sufficient to give a fair estimate of $`\mathrm{}_e`$ for both chiral or achiral tubes, in qualitative agreement with prior results, demonstrated for the armchair case. Within the RTA, one can write $`\mathrm{}_e=v_F\tau _e=\mathrm{}v_F/(2\mathrm{}m\mathrm{\Sigma }(E_F))`$, where $`v_F`$ is an average of the group velocity over the Fermi surface, $`\tau _e`$ is the mean free time, and $`\mathrm{\Sigma }(E_F)`$ is the self energy due to scattering events. The calculation of the electronic velocity is performed by a linearization of the dispersion relations in the vicinity of the Fermi level. For all metallic nanotubes, one finds to the lowest order of approximation that $`\sqrt{v_k^2}=v_F=\sqrt{3}a\gamma _0/2\mathrm{}`$ which leads typically to $`v_F=8.5\times 10^5`$ ms<sup>-1</sup>. On the other hand, in the vicinity of the Fermi level, the density of states is given by
$$\begin{array}{ccc}\hfill \rho (E)& =& \text{Tr}[\delta (E)]\hfill \\ & =& \frac{2}{\stackrel{~}{\mathrm{}}}\underset{n}{}๐k\delta (kk_n)\times \left|\frac{\epsilon _{nk}}{k}\right|^1\hfill \\ & =& \frac{2}{\pi \gamma _0\sqrt{n^2+m^2+nm}}\hfill \end{array}$$
(14)
where $`\stackrel{~}{\mathrm{}}=\mathrm{\Delta }k|\stackrel{}{๐}_h|/2\pi `$ is the volume of $`k`$-space per allowed $`k`$ value, divided by the spacing between lines for these allowed values. Equation (14) gives the TDOS per unit length of metallic CNs along the direction of the tube axis. Application of the Fermi golden rule yields $`\mathrm{}m\mathrm{\Sigma }(E_F)\epsilon _i^2G_0(i,i,E_F)\pi W^2\rho (E_F)`$, where $`W`$ is the disorder bandwidth and $`G_0(i,i,E_F)`$ is the average of the local on-site Green function elements. An estimate of the carrier mean free path (identical for metallic nanotubes with the same radius) is:
$$\mathrm{}_e\frac{3\sqrt{3}a\gamma _{0}^{}{}_{}{}^{2}}{2W^2}\sqrt{n^2+m^2+nm}.$$
(15)
The important result that $`\mathrm{}_e`$ is proportional to the nanotube diameter is relevant to the fact that the DOS at the Fermi energy decreases with increasing diameter. Thus, for a given disorder acting as a weak perturbation coupling some eigenstates to others close to Fermi surface, there are less available states to be scattered into, as the tube diameter decreases, yielding an enhancement of the mean free path. For instance for (9,0) and (10,10) CNs, with respective diameters 0.7 nm and 1.37 nm, the corresponding mean-free paths are estimated to be $`0.9\mu m`$ and $`1.8\mu m`$, which are much larger than the circumference length and are about the typical length of the systems themselves (for a reasonable disorder $`W0.2eV`$ suggested by Ref. ). In experiments on SWNTs, or MWNTs in the metallic-like regime (i.e., with an Ohmic temperature dependence of the resistivity), the resistivity should scale as $`\rho (d_{nt})1/d_{nt}`$ for a given temperature, where $`d_{nt}`$ is the tube diameter. It would be interesting to analyze the departure from this law as the diameter of the MWNTs increases, since such departures would indicate a participation of the inner tubes to the measured transport regime. In particular, some activation transport process of inner semiconducting tubes may contribute or not, depending on the temperature, and thus producing superimposed $`\rho (d_{nt})\mathrm{exp}(\alpha _0d_{nt})`$ factors (where $`\alpha _0`$ is a constant). One notes that the Fermi golden rule does not include the quantum interferences which lead to localization. As nanotubes are basically 1D-systems, from their estimated mean free paths, one can estimate the localization lengths at the Fermi level following the Thouless argument. In thin wires, Thouless argued that by writing $`R2\mathrm{}/e^2`$ for the resistance, there should be a transition to a localized state and an exponential increase of the resistance $`R`$. Consequently, from the resistance $`R=1/G=1/\sigma \times L^{2D}`$ ($`G`$ is the conductance, $`\sigma `$ the conductivity and $`D`$ the dimension of the system), assuming that the conductivity is calculated for a free electron gas in the wire ($`\sigma =e^2/\mathrm{}\times k_F\mathrm{}_e/3\pi ^2`$), it is straightforward to deduce the localization length $`\xi =(2Ak_F^2\mathrm{}_e)/3\pi ^2`$ (we take $`A=L^{D1}`$ as the cross section of the wire and $`L=\xi `$ as the wire length for which $`R=2\mathrm{}/e^2`$). Rewriting the last expression as $`\xi =(A/\lambda _F^2)\mathrm{}_e`$, the localization length $`\xi `$ is shown to be related to the approximate number of independent electrons (with spatial extension $`\lambda _F^2`$, where $`\lambda _F`$ is the Fermi wavelength) which can be accommodated through the cross-section of a single tube times the average mean free path. When confinement effects appear in the direction perpendicular to the cylinder axis, the number of allowed states reduces to the number of bands crossing at the Fermi level (conduction modes), since quantization prevails, and thus $`\xi 2N_{ch}\mathrm{}_e`$, with $`N_{ch}=2`$, the number of bands crossing the Fermi level for metallic nanotubes. Accordingly, for the $`(10,10)`$ nanotube, the localization length is estimated to be $`\xi 5`$$`\mu `$m so that it is typically larger or similar to a typical nanotube length and thus no special effects due to localization (insulating regime) are to be expected in SWNTs. With respect to the magnetic field, renormalization group arguments suggest that the localization lengths will be weakly affected by the magnetic field in quasi-1D systems. However, transport properties may be affected by the dephasing magnetic length $`\mathrm{}_m=\sqrt{\mathrm{}/eB}`$ which is $`\mathrm{}_m=25.66`$ nm or $`8`$ nm, respectively, for $`B=1`$Tesla or $`B=10`$Tesla. Two cases may be distinguished, whenever the mean free path is much larger than the nanotube circumference. Indeed, following Beenakker and Van Houten , and noticing that in the case of a nanotube, the cross section of the quantum wire reduces to $`\lambda _F/2`$, the magnetic phase relaxation time $`\tau _B`$ for weak and strong magnetic field limits can be written as follows (for $`\lambda _F\mathrm{}_e`$):
$`\tau _B\{\begin{array}{ccc}12\left({\displaystyle \frac{\mathrm{}_m}{\lambda _F}}\right)^2\tau _e\hfill & \mathrm{if}\hfill & \mathrm{}_m\sqrt{{\displaystyle \frac{\lambda _F\mathrm{}_e}{2}}}\hfill \\ 128\left({\displaystyle \frac{\mathrm{}_m^4}{\lambda _F^3\mathrm{}_e}}\right)\tau _e\hfill & \mathrm{if}\hfill & \mathrm{}_m\sqrt{{\displaystyle \frac{\lambda _F\mathrm{}_e}{2}}}\hfill \end{array}`$ (18)
So for the considered $`(9,0)`$ and $`(10,10)`$ tubes, we find that $`\lambda _F\mathrm{}_e/218.24`$nm and $`25.8`$nm respectively (we take $`\lambda _F=0.74`$nm and $`W=0.2078`$eV). Thus for $`B=1`$Tesla we conclude that $`\mathrm{}_m\sqrt{\lambda _F\mathrm{}_e/2}`$ which is an intermediate situation between low and strong magnetic field, and no analytical expression of $`\tau _B`$ can be deduced, whereas $`B=10`$Tesla is closer to the strong magnetic field limit where magnetic phase relaxation can be estimated analytically by
$$\tau _B=\frac{12\gamma _0\mathrm{}}{(\lambda _FW)^2}\mathrm{}_m^2\sqrt{n^2+m^2+nm}$$
(19)
which gives for the $`(10,10)`$ armchair nanotube a dephasing rate of $`\tau _B2.8\times 10^{11}s`$. This means that for an $`(10,10)`$ armchair tube, the electronic phase is randomized by a 10 Tesla magnetic field roughly every 30 ps, which indicates that in the strong field limit, the dephasing rate due to a magnetic field is just a few times the mean free times (we find $`\tau _B/\tau _e3.14`$ with the aforementioned parameters), so that $`\tau _B`$ should strongly contribute to damp the quantum interferences in the weak localization regime. Actually, both the mean free time and the magnetic dephasing rate scale linearly with diameter. Inelastic dephasing rates due to electron-phonon coupling have been recently evaluated by Suzuura and Ando, and an interesting chirality-dependent dominant inelastic backscattering mechanism (breathing, stretching versus twisting) was revealed. From their analytical estimate of the electron-phonon inelastic scattering rate ($`\tau _{elph}`$), in the $`(10,10)`$ armchair case, one finds at room temperature $`\tau _{elph}1.34\times 10^5s`$ which is much larger than former estimated coherence times (elastic mean free time and magnetic dephasing rate), so that one deduces that the main damping effect for the weak localization regime is likely to be dominated by the magnetic field strength.
Finally, with regard to the experimental results, one notices that if electronic transport is conveyed only by the outer shell of a metallic-like nanotube, as the magnetic field tends to decrease the TDoS at discrete values of the magnetic field corresponding to $`(2n+1)/2\varphi _0`$ with integer $`n`$, an increase of resistance may follow from increasing the magnetic field. However, the short electronic coherence lengths that are observed in a magnetic field, the negative magnetoresistance and the $`\varphi _0/2`$ AharonovโBohm oscillatory pattern are unlikely to be fully due to spectral effects (on the density of states), and a calculation of diffusion coefficients is mandatory to account for interferences between propagating electronic pathways in the nanotube structure. A recent study of the Kubo formula has been performed in Ref. and gives a geometrical explanation for the enhanced backscattering in MWNTs.
###### Acknowledgements.
S.R. acknowledges the European NAMITECH Network for financial support \[ERBFMRX-CT96-0067 (DG12-MITH)\]. Part of the work by R.S. is supported by a Grant-in Aid for Scientific Research (No. 11165216) from the Ministry of Education and Science of Japan and the Japan Society for the Promotion of Science for the international collaboration. The MIT authors acknowledge support under NSF grants DMR 98-04734 and INT 98-15744.
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# Easy collective polarization switching in ferroelectrics
\[
## Abstract
The actual mechanism of polarization switching in ferroelectrics remains a puzzle for many decades, since the usually estimated barrier for nucleation and growth is insurmountable (โparadox of the coercive fieldโ). To analyze the mechanisms of the nucleation we consider the exactly solvable case of a ferroelectric film with a โdeadโ layer at the interface with electrodes. The classical nucleation is easier in this case but still impossible, since the calculated barrier is huge. We have found that the interaction between the nuclei is, however, long range, hence one has to study an ensemble of the nuclei. We show that there are the ensembles of small (embryonic) nuclei that grow without the barrier. We submit that the interaction between nuclei is the key point for solving the paradox.
\]
The polarization switching in ferroelectrics (FEs) is most commonly used in applications (capacitors, memory elements), yet this process remains the least understood in spite of numerous experimental studies. As a rule, the ferroelectrics are switching in the field $`E_c`$ which is some order of magnitude lower than the so called โthermodynamic coercive fieldโ $`E_{c0}`$. The latter is the field at which the homogeneously polarized ferroelectric loses its stability in external field applied in opposite direction to the polarization. The fact that the switching takes place well before the point of instability is reached means that it proceeds by inhomogeneous nucleation and growth of domains of a new phase. But how the domains nucleate? The difficulty in answering this question has been emphasized by Landauer in late 1950s . His and later estimates showed that the energy barrier for creating a nucleus with reversed polarization is practically insurmountable, $`U^{}10^3k_BT`$ in the field of about $`100`$kV/cm. This problem, or a โparadox of a coercive fieldโ, has in fact been realized for domain nucleation in ferromagnets in 1938. The experimental coercive field for the bulk magnetic materials was known to be many times smaller than that suggested by the micromagnetics theory (โBrownโs paradoxโ) . The disagreement is especially striking for hard magnetic materials, where the situation is closer to the ferroelectrics considered in the present paper. There were numerous suggestions over the years that some defects can assist the nucleation and reduce the coercive field down to experimentally observed values. The situation in magnetics was summarized by Brown in 1965 who noted that the idea is plausible but โthere has been no strikingly successful calculation based on a completely realistic modelโ. Since then the situation has apparently remained the same, whereas a switching in fine magnetic particles seems to be fairly described by micromagnetics theory.
The goal of the present paper is to suggest a possible new mechanism for the polarization switching in ferroelectric materials. To this end we consider an exactly solvable model of a ferroelectric material with an extended inhomogeneity of its dielectric response at the ferroelectric-electrode interface, which is called a โpassiveโ or โdeadโ layer (for references see e.g.).
The main feature of the ferroelectric film with the dead layer is that the film exists in a polydomain state at any thickness of the dead layer . The monodomain state in this system, which can be produced by cooling in the field, would tend to transform into the equilibrium polydomain state. The transformation is favored by the fact that the polydomain state reduces the energy of the depolarizing field in spite of increased surface energy of the domain walls. This is contrary to the usual notion that the depolarizing field hinders the nucleation, yet the classical nucleation remains impossible.
First, we shall study the classical nucleation of the individual domains of the new phase. The nuclei will be assumed below to have a form of stripes or cylinders with the domain walls perpendicular to the plane of the film ($`c`$-domains), which is a reasonable approximation. We shall evaluate the barrier for their nucleation and show that for the individual nucleus it is practically insurmountable, although the dead layer helps to reduce it. We then abandon the classical approach and study the interaction between nuclei and find that it is long range.
It becomes clear that when an individual nucleus cannot grow the ensemble of nuclei may be able to. As an example of such an ensemble we consider a periodic array of nuclei and show that it indeed provides a path to the equilibrium state. We shall show that there is no energy barrier for its growth already when the nuclei become larger than the thickness of the domain wall $`W`$. After this โembryonicโ state has been passed the free energy of the system decreases monotonously as the nuclei grow. Obviously, the energy of embryonic nuclei is much smaller than the critical energy of the Landauerโs nucleus. This is simply related to the fact that the size of the critical Landauer needle-like nucleus is large, the radius of its base is $`r_c^\text{L}=WE_{c0}/E_{\text{ext}}W,`$ where $`E_{\text{ext}}`$ $`\left(E_{c0}\right)`$ is the external field, therefore its energy $`U_\text{L}^{}`$ is huge. In the present case of the dead layer $`E_{\text{ext}}`$ should be replaced by the (small) depolarizing field in the ferroelectric and the critical radius of the Landauerโs nucleus remains $`W.`$ Since for the collective nucleation $`r_c^{\text{coll}}Wr_c^\text{L},`$ the incurred energy barrier for a nuclei to grow beyond the embryonic state $`U_{\text{coll}}^{},`$ if any, should be much smaller than the standard barrier $`U_\text{L}^{}`$ for individual nuclei. The equilibrium density of the embryos is large (see below) and the nucleation of the macroscopic domain (and the coercive field) would be determined by the waiting time for optimal fluctuation and its dependence on the electric field for an ensemble which can grow without, or almost without, the barrier,
We shall mainly discuss the energy aspects of the nucleation. What the present analysis demonstrates is that the nucleation of ferroelectric domains is facilitated by an interaction between the nuclei. Note that no such interaction is taken into account within the Kolmogorov-Avrami theory which is widely used to treat switching in ferroelectrics.
We shall consider first the problem of the barrier for single stripe and then cylindrical domains, which appear to be insurmountable. We establish, however, that the interaction between the nuclei is long range and may give a clue to actual nucleation process. We then turn to exactly solvable case of an ensemble of stripe domains and show that the barrier for nucleation is actually zero already when the nuclei are at the embryonic stage.
The geometry of the present problem for ferroelectric with dead layer is illustrated in Fig. 1. For a short-circuited electrodes (zero bias voltage) the free energy of the system is $`\stackrel{~}{F}=F_0+U_{es},`$ where the electrostatic energy $`U_{es\text{ }}`$is :
$$U_{es}=\frac{1}{2}_{\mathrm{FE}}\text{d}๐\sigma \phi $$
(1)
where $`\sigma `$ is the bound charge due to the spontaneous polarization, $`\phi `$ is the electrostatic potential, while $`F_0`$ includes the surface energy of the domain walls, and the integration goes over the FE surface. The electrostastic potential $`\phi `$ is found from solving the Poisson equation for assumed domain structure . We obtain with the use of the Fourier transformation the total electrostatic energy for arbitrary one-dimensional domain structure as
$$U_{es}=2_{\mathrm{}}^{\mathrm{}}\frac{\text{d}k}{k}\frac{\sigma _k^2}{\sqrt{\epsilon _a\epsilon _c}\mathrm{coth}\left(\sqrt{\frac{\epsilon _a}{\epsilon _c}}\frac{kl}{2}\right)+\epsilon _g\mathrm{coth}\frac{kd}{2}},$$
(2)
where $`\sigma _k_{\mathrm{}}^{\mathrm{}}๐x\mathrm{exp}(ikx)\sigma (x,z=l/2)`$ is the Fourier component of the surface bound charge $`\sigma (x,z=l/2),`$ $`l`$ the thickness of the FE film with $`\epsilon _{c(a)}`$ the dielectric constants in $`c`$ $`(a)`$ direction, $`d`$ the thickness of the dead layer (Fig. 1).
We begin with the case of a nucleus in the center of the plate, i.e. with the following distribution of the bound charge
$`\sigma (x,z`$ $`=`$ $`+()l/2)=(+)P_s,|x|<a/2;`$ (3)
$`\sigma (x,z`$ $`=`$ $`+()l/2)=+()P_s,a/2<|x|<R,`$ (4)
where $`2R`$ is the width of the plate and $`a`$ is the width of the nucleus, and $`P_s`$ the spontaneous polarization. In this case $`\sigma _k=\frac{2P_s}{k}\left(2\mathrm{sin}\frac{ka}{2}+\mathrm{sin}kR\right).`$ The dependence of all physical quantities on $`R`$ disappears in a limit $`R\mathrm{}`$, as it should. It is handy to always subtract the (constant) electrostatic energy of the uniformly polarized sample, which is characterized by the Fourier transform of the bound charge $`\overline{\sigma }_k=(2P_s/k)\mathrm{sin}kR.`$ By using Eq.(4) one then finds the change of the electrostatic energy due to creation of a stripe nucleus
$`U_{es}^{\mathrm{stripe}}`$ $`=`$ $`8\pi \epsilon _g^1P_s^2da\left[1{\displaystyle \frac{a}{\pi d}}\mathrm{ln}\left({\displaystyle \frac{e^{3/2}d}{a}}\right)\right],ad;`$ (5)
$`=`$ $`8\epsilon _g^1P_s^2d^2\mathrm{ln}\left({\displaystyle \frac{e^{3/2}a}{d}}\right),ad;`$ (6)
for $`d<l\sqrt{\epsilon _c/\epsilon _a},`$ $`\epsilon _g=\sqrt{\epsilon _a\epsilon _c}`$. Note that this is the change in electrostatic energy with respect to uniformly polarized sample when the nucleus is present, so it does not apply to completely reversed sample. The electrostatic energy favors the nucleation, since the domains reduce the energy of stray field . The total energy of the stripe nucleus per unit length is
$$\stackrel{~}{F}_{\text{stripe}}(a)=2l\gamma +U_{es}2P_sE_{\text{ext}}al,$$
(7)
where $`\gamma =P_s^2\mathrm{\Delta }`$ is the surface energy of the domain wall with $`\mathrm{\Delta }`$ the temperature dependent characteristic length. We see that the gain in electrostatic energy eventually overwhelms the surface energy, and there appears an exponentially wide barrier for the nucleus in $`\stackrel{~}{F}(a)`$ when $`E_{\mathrm{ext}}=0.`$
We shall now consider the case of two nuclei of a new phase in a form of stripe domains. With the help of Eq.$`\left(\text{2}\right)`$ one can easily calculate the change of the electrostatic energy due to formation of two nuclei having the same width $`a`$, with the separation $`r`$ between their centers. We find for the energy of the interaction of two stripe domains per unit length for $`ra`$
$$U_{int}^{\mathrm{stripe}}\left(r\right)=2\epsilon _g^1P_s^2a^2d^2\frac{1}{r^2},$$
(8)
which corresponds to long range dipole-dipole interaction between two stripes . This observation indicates that the interactions in a system with an ensemble of nuclei would be very important. Note that the interaction of the stripe nucleus with the edge of the sample $`U_{\text{edge}}`$ is also long range, $`U_{\text{edge}}(x_0)=2\epsilon _g^1P_s^2ad^2/x_0,`$ and they are repelled from the edge.
Similar treatment can be repeated for a single and a pair of cylindrical nuclei. The expression for the electrostatic energy for cylindrical nucleus is similar to that for the stripe case (2)
$$U_{es}=2_0^{\mathrm{}}\text{d}k\frac{\sigma _k^2}{\sqrt{\epsilon _a\epsilon _c}\mathrm{coth}\left(\sqrt{\frac{\epsilon _a}{\epsilon _c}}\frac{kl}{2}\right)+\epsilon _g\mathrm{coth}\frac{kd}{2}},$$
(9)
where for the nucleus with the radius $`a`$ in the center of the slab with the radius $`R`$ we obtain $`\sigma _k2\pi _0^R`$d$`rrJ_0(kr)\sigma (r,z=l/2)=(2\pi P_s/k)\left[RJ_1(kR)2aJ_1(ka)\right],`$ while for the uniformly polarized sampe $`\overline{\sigma }_k=(2\pi P_s/k)RJ_1(kR),`$ with $`J_n(z)`$ the Bessel function. The integral in the expression (9) for the electrostatic energy of the cylindrical nucleus can be evaluated with the result
$`U_{es}^{\mathrm{cyl}}`$ $`=`$ $`8\pi ^2\epsilon _g^1P_s^2da^2,ad;`$ (10)
$`=`$ $`8\pi \epsilon _g^1P_s^2d^2a\mathrm{ln}{\displaystyle \frac{8a}{e^{1/2}d}},ad;`$ (11)
The free energy of one cylindrical nucleus is
$$\stackrel{~}{F}_{\text{cyl}}(a)=2\pi al\gamma +U_{es}2\pi P_sE_{\text{ext}}la^2.$$
(12)
It is obvious from Eqs. (11),(12) that the gain in the electrostatic energy eventually overwhelms the growth of the the surface energy of the domain wall with increase of the radius $`a`$ of the nucleus. The critical radius is exponentially large when $`E_{\mathrm{ext}}=0`$, it behaves roughly as $`a_cd\mathrm{exp}\left(0.3a_K^2/d^2\right),`$ where $`a_K=\left(0.3\stackrel{~}{\epsilon }\mathrm{\Delta }l\right)^{1/2}`$ is the Kittel period, $`\stackrel{~}{\epsilon }=\epsilon _g+\sqrt{\epsilon _a\epsilon _c}`$. The corresponding barrier height $`\stackrel{~}{F}_c`$ is very large, $`\stackrel{~}{F}_cP_s^2d_{at}^3E_{at},`$ where $`d_{at}`$ is the characteristic โatomicโ length (of the order of the lattice parameter) and the โatomicโ energy $`E_{at}`$ amounts to a few eV. Thus, in this case the barrier is also huge, comparable to the barrier for the Landauerโs nucleus, and its growth is prohibitively expensive.
The interaction between two cylindrical nuclei can be estimated for $`ra`$ as
$$U_{int}^{\mathrm{cyl}}\left(r\right)=16\pi ^2\epsilon _g^1P_s^2a^4d^2\frac{1}{r^3}$$
(13)
In this case it appears to also be long range of a dipole-dipole type .
The importance of these interactions lies in their long range rather than sign. Therefore, one has to accurately evaluate the total energy for a system of nuclei, since it does not reduce to a sum of asymptotic interactions (8) or (13). The long range interactions give us a clue to the mechanism of nucleation and growth of a new phase. To illustrate that it indeed may solve the problem, we shall consider a system of stripe-like domains. The results below will demonstrate that the electrostatic energy not only favors nucleation but eliminates an energy barrier for the growth of nuclei in an ensemble when their size is larger than the domain wall width, i.e. the switching proceeds collectively.
To study the collective nucleation we shall analyze a domain structure with stripe domains of opposite polarization of widths $`a_1`$ and $`a_2`$, the period $`T=a_1+a_2,`$ and the asymmetry parameter $`\delta =\frac{a_1a_2}{a_1+a_2},`$ which measures a net polarization of the film . For zero external bias voltage assumed throughout the present paper $`\delta =0`$ in the equilibrium (polydomain) state, whereas in the monodomain state $`\delta =1.`$ At $`\delta 1`$ the system consists of very narrow domains with the polarization opposite to the net polarization (Fig. 1, inset), i.e. it is a periodic ensemble of nuclei. The free energy $`\stackrel{~}{F}`$ of this system at zero external bias is given by
$`{\displaystyle \frac{\stackrel{~}{F}}{๐P_s^2}}`$ $`=`$ $`{\displaystyle \frac{2\pi d\delta ^2}{\epsilon _c\left(d/l\right)+\epsilon _g}}+{\displaystyle \frac{2\mathrm{\Delta }l}{T}}+{\displaystyle \frac{16T}{\pi ^2}}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{\left(2j+1\right)^3}}{\displaystyle \frac{1}{D_{2j+1}}}`$ (15)
$`+{\displaystyle \frac{8T}{\pi ^2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(1\right)^n{\displaystyle \frac{1\mathrm{cos}\pi n\delta }{n^3}}{\displaystyle \frac{1}{D_n}},`$
where $`D_n=\sqrt{\epsilon _a\epsilon _c}\mathrm{coth}\sqrt{\frac{\epsilon _a}{\epsilon _c}}\frac{\lambda _nl}{2}+\epsilon _g\mathrm{coth}\frac{\lambda _nd}{2}`$ and $`\lambda _n=2\pi n/T`$ . Fortunately, this free energy can be found analytically in the asymptotic case of narrow dead layer $`(dl)`$ and correspondingly wide period of the domain structure, $`Td,`$ when $`\epsilon _g=\sqrt{\epsilon _a\epsilon _c}.`$ We have found earlier that the period of the domain structure in this case is exponentially large, $`T=0.95d\mathrm{exp}(0.4a_K^2/d^2)`$. We note that in this case $`D_n=\epsilon _g\left(1+\mathrm{coth}\frac{\lambda _nd}{2}\right)`$ and summation in $`\left(\text{15}\right)`$ can be performed to yield
$`{\displaystyle \frac{\stackrel{~}{F}}{๐P_s^2}}`$ $`=`$ $`{\displaystyle \frac{2\pi d\delta ^2}{\epsilon _c\left(d/l\right)+\epsilon _g}}+{\displaystyle \frac{2\mathrm{\Delta }l}{T}}+{\displaystyle \frac{4T}{\pi ^2\epsilon _g}}[\zeta (3)2Li_3\left(e^b\right)`$ (18)
$`+{\displaystyle \frac{1}{4}}Li_3\left(e^{2b}\right)Li_3\left(e^b\right)`$
$`ReLi_3(e^{i\pi \delta })+ReLi_3(e^{i\pi \delta b})],`$
where $`b=4\pi d/T1,`$ and $`Li_n(z)_{k=1}^{\mathrm{}}z^k/k^n`$. For the case of an ensemble of narrow domains, $`\delta 1`$ the free energy can be found from the known asymptotic behavior of the $`Li_n(z)`$ function , yielding an approximate expression
$$\frac{\stackrel{~}{F}}{๐P_s^2}=\frac{2\pi d\delta ^2}{\epsilon _c\left(d/l\right)+\epsilon _g}+\frac{2\mathrm{\Delta }l}{T}+\frac{T}{\epsilon _g}\left(1\delta \right)^2\mathrm{ln}\frac{e^3b^2}{\pi ^2(1\delta )^2}.$$
(19)
Now everything depends on how the free energy $`\stackrel{~}{F}`$ behaves as a function of $`\delta `$ when the nuclei grow, i.e. when $`\delta `$ reduces from unity towards zero. One can easily see that at $`x1\delta 1`$ the free energy $`\stackrel{~}{F}/๐P_s^2`$ is given, with respect to a constant, by the function
$$f(x)\frac{4\pi dx}{\epsilon _c\left(d/l\right)+\epsilon _g}+\frac{T}{\epsilon _g}x^2\mathrm{ln}\frac{e^3b^2}{\pi ^2x^2},$$
(20)
where $`f(0)=0.`$ This function does not have a barrier as a function of $`x`$ when $`T`$ is kept constant. Indeed, the second term in (20) has an exponentially small maximum ($`=T\epsilon _g^1x_0^2`$) at $`x_0=4ed/T1.`$ It is, however, suppressed by the first term, which corresponds to the energy of homogeneous field created by the net polarization and is linear in $`x.`$ As a result, there appears to be no barrier for the growth of nuclei. Note that we have actually restricted the systemโs path for nucleation by constraining the domain pattern to a fixed period. Even under this constraint the growth proceeds without the energy barrier, and this would be even more so if we were to lift the constraint and allow the system to follow an optimal path to equilibrium. This behavior does not depend on the approximation we have made for evaluating the free energy, the exact calculation of the free energy (15) for all $`\delta `$ shows that the collective growth of nuclei we just described proceeds without the barrier (Fig. 1). As mentioned above, the smallest size of nuclei where the present analysis applies is of the order of the domain wall width $`W`$ and the barrier for the nucleation of such small embryonic nuclei is expected to be zero or much smaller than the usual estimates for individual nucleation.
For nucleation to proceed by the present mechanism the only condition is the presence of extended dielectric inhomogeneity in a sample. The external field does promote growth of the nuclei, but the nucleation in the present system occurs even without it. The likely requirement is that the lateral extent of this inhomogeneity should be much larger than the period of the equilibrium domain (Kittel) structure $`a_K.`$ One can estimate the rate of embryo nucleation in $`1\mathrm{c}\mathrm{m}^3`$ as $`(1/W^3\tau _{ph})\mathrm{exp}(U_{em}/kT),`$ where $`\tau _{ph}`$ is the characteristic (optical phonon) time, and the โatomicโ estimate of the energy of the embryo is $`U_{em}\gamma W^2E_{at}\sqrt{T_c/T_{at}}2310^3\mathrm{K}\text{[14]},`$ with the characteristic โatomicโ temperature $`T_{at}10^4\mathrm{K}`$. Taking a conservative estimate of the embryo lifetime as $`10\tau _{ph},`$ one finds the equilibrium density of the embryos $`10^{17}`$ cm$`^3.`$ At such high densities the embryos should โfeelโ the field of each other, and favorable ensembles should appear within a reasonable time, unlike in the case of the Landauerโs nucleus where the expectation time much exceeds the lifetime of the universe.
In conclusion, we have suggested a possible way to solve the โparadox of the coercive fieldโ by demonstrating a collective mode of domain growth past the embryonic stage with sizes about the domain wall thickness $`W`$, which proceeds without energy barrier. The origin of this cooperative phenomenon is a long-range interaction of electrostatic origin between the nuclei. The possible screening by free charges in ferroelectric does not seem to be important, since the conductivity of ferroelectrics is usually too low to have any effect . As a corollary, we note that the Kolmogorov-Avrami model (KA) is inapplicable to growth of domains in ferroelectrics, since the essential long-range interaction between nuclei is completely neglected in this approach (note that KA fully bypasses the question of how the domains were nucleated in the first place). The present results are general, and have the implication that in ferroelastic materials and possibly also magnetic materials the nucleation would be facilitated by the long-range interaction between nuclei of a new phase. This could alleviate the Brownโs paradox of the coercive field in relation to switching in bulk ferromagnets.
We acknowledge helpful discussions with A. Aharoni, V.V. Osipov, A.L. Roytburd, and A.K. Tagantsev.
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# Part I
## Part I
The most elementary proof of the Prime Numbers Theorem
These few lines are not part of the proof. They simply show the history which has led to the starting point of our proof. The two main steps of it involve only formal elementary algebra, with no recourse naturally to Functionsโ theory nor complex variables. Euler proved in 1747 , quite formally, that the prime numbers are linked with natural integers, by establishing what is universally known as โEuler identityโ and which is so famous that we do not recall it here.
From this identity, one can deduce straight forwardly an approximate formula<sup>1</sup><sup>1</sup>1Which, in fact, is the strict equivalent of Euler identities. giving
$$(\mathrm{I})\underset{p<\mathrm{\Lambda }}{}1/p=1\mathrm{log}(\mathrm{log}\mathrm{\Lambda })+\mathrm{}.$$
This result has been refined, mainly by Mertens with the aim to establish the value of constants possibly entering a much more exact expression of the sum in (I). This expression, we call it โEuler-Mertens identityโ and is the starting point of our proof, the formula 1 of our theorem.
Proof Consider the Euler-Mertens identity (I), exact when $`\mathrm{\Lambda }\mathrm{}`$. We introduce it under the following form (1) and apply the RG analysis
$$\overline{F}=(\mathrm{log}\mathrm{log}\mathrm{\Lambda })^_1^\mathrm{\Lambda }\frac{dn\overline{n}(n)}{n}C^{te}1,\mathrm{\Lambda }\mathrm{}$$
(1)
$`\overline{n}(n)dn`$ is a measure $`df(n)`$ of the Stieltjes type and $`\overline{n}(n)`$ can be considered as a density in the physical sense<sup>2</sup><sup>2</sup>2$`df(n)=\overline{n}(n)dn`$ and, according to the result (5) $`f(\mathrm{\Lambda })=^\mathrm{\Lambda }\frac{dn}{\mathrm{log}n}Li(\mathrm{\Lambda })`$..
Then , since
$`\mathrm{\Lambda }/\mathrm{\Lambda }\overline{F}`$ $`=`$ $`{\displaystyle \frac{}{\mathrm{log}\mathrm{\Lambda }}}(\mathrm{log}\mathrm{log}\mathrm{\Lambda })^1{\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{dn\overline{n}(n)}{n}}+{\displaystyle \frac{\overline{n}(\mathrm{\Lambda })}{\mathrm{log}\mathrm{log}\mathrm{\Lambda }}}`$
$`=`$ $`(\mathrm{log}\mathrm{\Lambda })^1(\mathrm{log}\mathrm{log}\mathrm{\Lambda })^2{\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{dn\overline{n}(n)}{n}}+{\displaystyle \frac{\overline{n}(\mathrm{\Lambda })}{\mathrm{log}\mathrm{log}\mathrm{\Lambda }}}`$
one deduces, by the RG method, (since $`d\overline{F}/d\mathrm{log}\mathrm{\Lambda }=0)`$
$`0`$ $`=`$ $`{\displaystyle \frac{(\mathrm{log}\mathrm{\Lambda })^1}{\mathrm{log}\mathrm{log}\mathrm{\Lambda }}}F+{\displaystyle \frac{\overline{n}(\mathrm{\Lambda })}{\mathrm{log}\mathrm{log}\mathrm{\Lambda }}}+\left[{\displaystyle \frac{\overline{n}(\mathrm{\Lambda })}{\mathrm{log}\mathrm{\Lambda }}}\right]{\displaystyle \frac{\delta }{\delta \overline{n}(\mathrm{\Lambda })}}\overline{F}`$ (2)
$`=`$ $`\left[(\mathrm{log}\mathrm{\Lambda })^1+\overline{n}(\mathrm{\Lambda })\right](\mathrm{log}\mathrm{log}\mathrm{\Lambda })^1+0\left({\displaystyle \frac{1}{\mathrm{\Lambda }\mathrm{log}^2\mathrm{\Lambda }\mathrm{log}\mathrm{log}\mathrm{\Lambda }}}\right).`$ (4)
So from 4, at this approximation one gets
$$\overline{n}(\mathrm{\Lambda })(\mathrm{log}\mathrm{\Lambda })^1,\mathrm{\Lambda }1$$
(5)
which is the prime Numbers theorem, since the density
$$\overline{n}(\mathrm{\Lambda })=\mathrm{\Lambda }^1\pi (\mathrm{\Lambda })$$
## Part II
The RG Equation for the Density of Prime Numbers<sup>3</sup><sup>3</sup>3The RG-method has been designed in order to know how the structure of a theory gets modified when the scale is changed.
The density of natural integers is scale invariant:
$$\lambda \frac{}{\lambda }\overline{d}(n_i\lambda ,\overline{d}(1))=0,$$
(6)
$`\overline{d}`$ being the density around $`n_i\lambda `$, and $`\lambda \frac{}{\lambda }`$ the generator of scale transformations of the natural integers.
If, on the other hand, $`n_i`$ is a prime and $`\lambda n_i`$ is around another prime, say $`n_i^{}`$, then
$$\overline{d}(n_i^{}\lambda n_i,\overline{d}(1))\overline{d}(n_i,\overline{d}(1)).$$
In such a case, instead of having an equation like (6), expressing the invariance for scale changes, one uses, as is customary, the so-called Renormalization Group (RG) equation (or better: renormalization transformation equation) which generally substitutes Eq. (6) when scale invariance is broken. The strategy is to compensate the broken invariance, for example in the present case, by a density $`\overline{d}`$, which this time depends upon $`\lambda `$ and is different from that of (6), namely $`\overline{d}(1)`$. The RG equation, as is well known , reads
$$\lambda \frac{}{\lambda }\overline{d}(n_i\lambda ,\overline{d}(\lambda ))+\left[\frac{\lambda }{\lambda }\overline{d}(\lambda )\right]\frac{}{\overline{d}(\lambda )}\overline{d}(n_i\lambda ,\overline{d}(\lambda ))=0.$$
(7)
Equation (7)<sup>4</sup><sup>4</sup>4$`n_i\lambda `$ prime, as $`n_i`$ is. introduces the quantity $`[\lambda \frac{}{\lambda }\overline{d}(\lambda )]\frac{}{\overline{d}(\lambda )}`$ which is a one-dimensional vector field $`\underset{ยฏ}{V}(\overline{d}(\lambda ))`$ on the axis of integer numbers.
Now, the problem is to solve (7) for $`\overline{d}`$.
This solution, as will be explained in the Appendix, when taken between two different numbers $`N_1`$ and $`N_2`$, turns out to be
$$\frac{1}{\overline{d}(N_1)}\frac{1}{\overline{d}(N_2)}=\mathrm{log}\frac{N_1}{N_2}.$$
(8)
But, perhaps, more instructively for the scope of this short note,
$$\overline{d}(t,\overline{d}(0,\overline{d}_0)=\frac{\overline{d}(0,\overline{d}_0)}{1+t\overline{d}(0,\overline{d}_0)}$$
(9)
with $`t=\mathrm{log}N`$, exhibiting the neperian logarithmic single power decrease as having its origin in the violation of the scale invariance symmetry.
Indeed, if in the region of large primes (say between $`10^{15}`$ and $`2\times 10^{17}`$), the numerical results obtained by the use of (8), formula (8) does not tell else that, for each number $`N_1`$ and $`N_2`$,
$$\frac{N}{\pi (N)}=\mathrm{log}N$$
which is a 100 years-old, over-demonstrated asymptotic result .
However, in this note, our aim is to look for the deep reason why the density of primes decreases with the single power of the natural logarithm. We hope that we have been able to shed some light on this fact: the breaking of a symmetry, namely that of scale invariance with generator $`\lambda \frac{}{\lambda }`$, is the very factor responsible for this specific decrease.
The coincidence of the results obtained is striking when compared to the formulas of the first non-trivial approximation of Quantum ChromoDynamics (mutatis mutandis, of course, the concepts between two such different fields).
But a main common feature emerges: in both cases the two fields are afflicted by the same broken symmetry, that of scale invariance.
APPENDIX
* For natural integers, scale invariance holds for the density $`\overline{d}(n)`$, i.e. when $`n\lambda n,\overline{d}(n)=\overline{d}(\lambda n)`$.
* For primes $`p_i`$,
$$p_i\lambda p_ip_j,\overline{d}(p_j)\overline{d}(p_i),$$
so that $`\overline{d}`$ becomes a function of $`\lambda `$.
One finds easily that a functional equation of the type
$$\overline{d}(\lambda n_i,\overline{d}(\lambda ))=\overline{d}(n_i,\overline{d}_0),$$
(A.1)
$`(\overline{d}_0`$ fixed, and $`n_i`$ representing primes as well as $`\lambda n_i`$ in the dose vicinity of $`n_j`$, prime itself) exists.
The RHS is $`\lambda `$-independent and one gets at once
$$\lambda \frac{}{\lambda }\overline{d}(\lambda n_in_j,\overline{d}(\lambda ))+\left[\lambda \frac{}{\lambda }\overline{d}(\lambda )\right]\frac{}{\overline{d}(\lambda )}\overline{d}(\lambda n_i,\overline{d}(\lambda ))=0.$$
(A.2)
Or else, calling $`t=\mathrm{log}\lambda `$ and passing to logarithmic variables
$$[\frac{}{t}+\underset{ยฏ}{V}(\overline{d}(t)]\overline{d}(t+\mathrm{log}n_i,d(t))=0$$
(A.3)
with
$$\underset{ยฏ}{V}(\overline{d}(t))=\left[\frac{}{t}\overline{d}(t)\right]\frac{}{\overline{d}(t)}.$$
To solve (A.2), one proceeds in the following way
$`\overline{d}(\delta t,\overline{d}(0,\overline{d}_0))`$ $`=`$ $`\overline{d}(0,\overline{d}(0,\overline{d}_0)`$
$`+\delta t\underset{ยฏ}{V}(\overline{d}(0,\overline{d}_0))\overline{d}(0,\overline{d}_0)+0(\delta t^2)`$
by Taylor-expanding $`\overline{d}`$ around $`\delta t=0`$.
According to the properties of flows of vector fields, one has $`\overline{d}(0,x)=x`$, that is, for example, $`\overline{d}(0,\overline{d}(0,\overline{d}_0))=\overline{d}(0,\overline{d}_0)=\overline{d}_0,\overline{d}_0`$ being a fixed arbitrary density.
One gets then
$$\overline{d}(\delta t,\overline{d}_0)\overline{d}_0=\delta t\underset{ยฏ}{V}(\overline{d}_0)\overline{d}_0$$
or<sup>5</sup><sup>5</sup>5In passing we recall a well-known property of the flows $`\overline{d}(t,x)`$; namely they satisfy the one-parameter Abelian group: $`\overline{d}(s+t,x)=\overline{d}(t,\overline{d}(s,x))`$, i.e. the composition law $`\overline{d}_{s+t}=\overline{d}_t\overline{d}_s`$. This group is trivially generated as the one-parameter group of diffeomorphisms by the vector field $`V`$ on the manifold considered. For details see Ref. a.)
$$\overline{d}(\delta t,\overline{d}_0)=(1+\delta t\underset{ยฏ}{V}(\overline{d}_0))\overline{d}_0$$
(A.4)
By theorems by Chebyshev and Mertens , $`V(x)`$ can be shown to be quadratic in its argument $`x`$.
It remains to exponentiate the RHS of (A.4):
$`\overline{d}(t,\overline{d}_0)`$ $`=`$ $`\left(1+t\underset{ยฏ}{V}(\overline{d}_0){\displaystyle \frac{}{\overline{d}_0}}+{\displaystyle \frac{t^2}{2!}}V(\overline{d}_0){\displaystyle \frac{}{\overline{d}_0}}V(\overline{d}_0){\displaystyle \frac{}{\overline{d}_0}}+\mathrm{}\right)\overline{d}_0`$ (A.5)
$`=`$ $`{\displaystyle \frac{\overline{d}_0}{1+t\overline{d}_0}};(V(\overline{d}_0)=\overline{d}_0^2.\text{see above}).`$ (A.7)
(A.7) is the formula (9) of the text and seems to us to be an explanation we were searching for, to explain the decrease of $`\overline{d}(t,d_0)`$ with a single power of the natural logarithm $`t=\mathrm{log}N`$ (Remember that $`\overline{d}(t,\mathrm{})\overline{d}(e^t,\mathrm{})`$, as (A.2) and (A.3) show without further comments.)
As a final remark, (8) follows straight forwardly from (9) by trivial algebra.
Take (9) with two different values for $`t`$: $`t_1=\mathrm{log}N_1`$, and $`t_2=\mathrm{log}N_2`$. It follows at once that
$$\overline{d}^1(t_1,\overline{d}_0)\overline{d}^1(t_2,\overline{d}_0)=\mathrm{log}N_1/N_2.$$
(Additionally it confirms the arbitrariness of $`\overline{d}_0`$ which might be chosen at will; $`\overline{d}_0=1`$, for example.)
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# Theorem 1.
Disjointly Strictly Singular Inclusions
of Symmetric Spaces
S.V.Astashkin
Samara State University, 443011 Samara 11, Russia
## Abstract
In this paper, the disjoint strict singularity of inclusions of symmetric spaces of functions on an interval is considered. A condition for the presence of a โgapโ between spaces sufficient for the inclusion of one of these spaces into the other to be disjointly strictly singular is found. The condition is stated in terms of fundamental functions of spaces and is exact in a certain sense. In parallel, necessary and sufficient conditions for an inclusion of Lorentz spaces to be disjointly strictly singular (and similar conditions for Marcinkiewicz spaces) are obtained and certain other assertions are proved.
Keywords: Banach space, disjointly strictly singular operator, inclusion operator, symmetric space, fundamental function, Lorentz space, Marcinkiewicz space, Orlicz space
Introduction
Recall that a bounded linear operator $`T`$ from a Banach space $`X`$ into a Banach space $`Y`$ is called strictly singular (or a Cato operator) if $`X`$ does not contain an infinite-dimensional subspace $`Z`$ such that the restriction of $`T`$ to $`Z`$ is an isomorphism.
In recent decades the class of strictly singular operators has been extensively studied (see the references in, e.g., the monograph ). One of the historically first results important for our purposes is the Grothendick theorem on strict singularity of the identity inclusion operator from $`L_{\mathrm{}}(\mathrm{\Omega },\mu )`$ into $`L_p(\mathrm{\Omega },\mu ),`$ where $`1p<\mathrm{}`$ and $`\mu `$ is a probability measure on $`\mathrm{\Omega }`$ (see or \[3, Theorem 5.2\]). However, as a rule, the identity inclusion operator from one symmetric space into another (the definition is given below) is not strictly singular because of the existence of โthroughโ subspaces (such as the subspace generated by the Rademacher functions ). In part because of this, the close notion of disjointly strictly singular operator was introduced in 1989 .
A bounded linear operator $`T`$ from a Banach lattice $`X`$ into a Banach space $`Y`$ is called disjointly strictly singular (or has the DSS property) if there exists no sequence of nonzero disjoint vectors $`\{x_n\}_{n=1}^{\mathrm{}}`$ in $`X`$ such that the restriction of $`T`$ to their closed linear hull $`[x_n]`$ is an isomorphism.
Clearly, any strictly singular operator is a DSS operator. A simple example shows that the converse is not true. For instanse, the identity inclusion operator $`I:L_p[0,1]L_q[0,1]`$ $`(1q<p\mathrm{})`$ has the DSS property, because the closed linear hull in $`L_r`$ of disjoint functions $`x_nL_r[0,1]`$ is isomorphic to $`l_r`$ $`(1r\mathrm{}):`$ $`[x_n]_rl_r.`$ However, if $`p<\mathrm{},`$ then Khintchineโs inequality implies that $`[r_n]_p[r_n]_ql_2`$ ($`r_n`$ are Rademacher functions), and therefore $`I`$ is not strictly singular. At the same time, it is easy to show that if $`X`$ has a Schauder basis of disjoint vectors, then the class of DSS operators on $`X`$ coincides with the class of strictly singular operators .
The notion of DSS operator proved important in studies of the geometric properties of function spaces. For example, the existence of operators without the DSS property makes it possible to construct complemented subspaces that admit โnonstandardโ projections onto them .
The goal of this paper is to study the following question: when does the identity inclusion operator (throughout, we denote it by $`I`$) from one symmetric space in another have the DSS property? The conditions are stated in terms of fundamental functions of these spaces.
If $`z=z(t)`$ is measurable on $`[0,1]`$ with respect to the Lebesgue measure $`\mu `$, then we call the function $`n_z(\tau )=\mu \{t:|z(t)|>\tau \}`$ $`(\tau >0)`$ the distribution function of $`z.`$ Two functions $`x(t)`$ and $`y(t)`$ are called equimeasurable if $`n_x(\tau )=n_y(\tau )`$ for $`\tau >0.`$
Recall that a Banach space $`E`$ of measurable functions on $`[0,1]`$ is called a symmetric space (briefly is an SS) if the following conditions hold:
(1) if $`yE`$ and $`|x(t)||y(t)|,`$ then $`xE`$ and $`xy;`$
(2) if $`yE`$ and functions $`x(t)`$ and $`y(t)`$ are equimeasurable, then $`xE`$ and $`x=y.`$
The fundamental function of an SS $`E`$ is defined by $`f_E(t)=\chi _{(0,t)}_E,`$ where, as usual, $`\chi _U(t)=1(tU),`$ $`\chi _U(t)=0(tU).`$ The function $`f_E(t)`$ is quasiconcave on $`(0,1]`$ \[8, p.137\], i.e., it is nonnegative, increasing, and $`f_E(t)/t`$ decreases. As is known (see, e.g., \[8, p.70\]), such a function is equivalent to its least concave majorant. Throughout, $`G`$ denotes the class of all positive increasing functions concave on $`(0,1].`$
An important example of an SS is an Orlicz space. Let $`N(t)`$ be an increasing convex function on $`[0,\mathrm{})`$ such that $`N(0)=0`$ and $`N(\mathrm{})=\mathrm{}.`$ The Orlicz space $`L_N`$ consists of all functions $`x=x(t)`$ measurable on $`[0,1]`$ and such that
$$_TN\left(\frac{|x(t)|}{u}\right)๐\mu <\mathrm{}$$
for some $`u>0;`$ the norm of this space is
$$x=inf\{u>0:_TN\left(\frac{|x(t)|}{u}\right)๐\mu 1\}.$$
Direct calculation shows that the fundamental function of the space $`L_N`$ is $`f_N(t)=1/N^1(1/t)`$ ($`N^1(u)`$ is the inverse of $`N(u)`$) .
In , the following disjoint strict singularity theorem for inclusions of $`L_N`$ into $`L_M`$ is proved.
Theorem. If $`L_NL_M,`$ then the following conditions are equivalent:
(1) the inclusion $`I:L_NL_M`$ is a DSS operator;
(2) for any $`n=1,2,..`$ and $`๐ฆ>0,`$ there exist $`1x_1<x_2<\mathrm{}<x_n`$ and $`c_1>0,..,c_n>0`$ such that
$$\underset{i=1}{\overset{n}{}}c_iN(tx_i)๐ฆ\underset{i=1}{\overset{n}{}}c_iM(tx_i)\text{for}t1.$$
Let us show that condition (2) follows from the relation
$$\underset{t+0}{lim}\frac{f_M(t)}{f_N(t)}=\mathrm{\hspace{0.17em}0},$$
where $`f_N`$ and $`f_M`$ are the fundamental functions of the respective Orlicz spaces. Indeed, this relation implies that $`N^1(t)hM^1(t)`$ for an arbitrary positive $`h1`$ and $`tt_0`$ and, since $`N(t)`$ is convex, $`M(t)N(ht)hN(t),`$ if $`tM^1(t_0).`$ Therefore,
$$\underset{t\mathrm{}}{lim}\frac{M(t)}{N(t)}=\mathrm{\hspace{0.17em}0},$$
and condition (2) holds.
Quite naturally, this observation leads us to the following general problem.
Suppose that functions $`\phi G`$ and $`\psi G`$ satisfy the condition
$`(A)lim_{t+0}\psi (t)/\phi (t)=\mathrm{\hspace{0.17em}0}`$
and are, respectively, the fundamental functions of symmetric spaces $`E`$ and $`F`$ such that $`EF.`$ Does the identity inclusion operator $`I:EF`$ have the DSS property?
In what follows, we show that the answer to this question for โclassicalโ symmetric spaces (such as the Lorentz and Marcinkiewicz spaces) as well as for Orlicz spaces is positive. Moreover, it is so for an inclusion of a Lorentz space into an arbitrary SS (and, vice versa, of an arbitrary SS into a Marcinkiewicz space).
However, in the general case, this is not true: this paper contains an example of two symmetric spaces $`E`$ and $`F`$ such that $`EF`$ and their fundamental functions satisfy condition (A), but $`I:EF`$ is not a DSS operator.
At the same time, it is possible to state a condition on fundamental functions stronger than (A) under which the answer to the stated question is positive for all symmetric spaces. First, recall the definition of the dilation function.
For a positive function $`f`$ on $`(0,1],`$ its dilation function $`_f(t)`$ is defined as
$$_f(t)=sup\{\frac{f(st)}{f(s)}:\mathrm{\hspace{0.33em}0}<s\mathrm{min}(1,\frac{1}{t})\}(t>0).$$
Since $`_f(t)`$ is semimultiplicative, there exist numbers
$$\gamma _f=\underset{t0}{lim}\frac{\mathrm{ln}_f(t)}{\mathrm{ln}t}\text{and}\delta _f=\underset{t\mathrm{}}{lim}\frac{\mathrm{ln}_f(t)}{\mathrm{ln}t},$$
which are called, respectively, the lower and upper dilation indices of function $`f`$. If $`\phi G`$, then we have $`0\gamma _\phi \delta _\phi 1`$ \[8, p.76\].
Let us introduce one more condition on the functions $`\phi `$ and $`\psi `$ from the class $`G:`$
$`(B)\gamma _{\psi /\phi }>\mathrm{\hspace{0.17em}0}.`$
The definition of lower dilation index readily implies that condition (A) follows from (B). The converse, of course, is not true: it suffices to take for $`\phi `$ and $`\psi `$ functions differing by a logarithmic factor (see also the proof of Theorem 3).
We shall show that, if (B) holds, then operator $`I:EF`$ will be a DSS operator for arbitrary symmetric spaces $`E`$ and $`F,`$ $`EF,`$ with fundamental functions $`\phi `$ and $`\psi ,`$ respectively. This result generalizes and simultaneously refines a similar theorem for the Orlicz spaces proved in . In parallel, we shall show that condition (A) is necessary and sufficient for the identity inclusion operator from one Lorentz space into another to have the DSS property (and similar assertion for Marcinkiewicz spaces). These results also supplement the theorem for Orlicz spaces proved in and cited above.
$`๐\mathrm{\hspace{0.17em}1}.`$ The inclusions $`\mathrm{\Lambda }(\phi )F`$ and $`EM(\theta )`$
For $`\phi G`$, the Lorentz space $`\mathrm{\Lambda }(\phi )`$ consists of all functions $`x=x(s)`$ measurable on $`[0,1]`$ and such that
$$x_{\mathrm{\Lambda }(\phi )}=_0^1x^{}(s)๐\phi (s)<\mathrm{},$$
where $`x^{}(s)`$ is a decreasing left-continuous rearrangement of the function $`|x(s)|`$ \[8, p.83\]. Clearly, the fundamental function of the Lorentz space $`\mathrm{\Lambda }(\phi )`$ is $`f_{\mathrm{\Lambda }(\phi )}(t)=\phi (t).`$
###### Theorem 1.
Suppose that the functions $`\phi G`$ and $`\psi G`$ satisfy (A) and $`F`$ is an SS on $`[0,1]`$ with fundamental function $`\psi (t).`$ Then $`\mathrm{\Lambda }(\phi )F`$ and the identity inclusion $`I:\mathrm{\Lambda }(\phi )F`$ is a DSS operator.
###### Proof.
Condition (A) and the continuity of the concave functions $`\phi `$ and $`\psi `$ at $`t>0`$ imply that $`\psi (t)C_1\phi (t)`$ for $`t[0,1].`$ By the definition of the norm of a Lorentz space, we have $`\mathrm{\Lambda }(\phi )\mathrm{\Lambda }(\psi ).`$ The Lorentz space with a given fundamental function is minimal among the symmetric spaces with the same fundamental function \[8, p.160\]; therefore, $`\mathrm{\Lambda }(\phi )F.`$
Suppose that $`I:\mathrm{\Lambda }(\phi )F`$ has not the DSS property. Then there exists a sequence of nonzero disjoint functions $`x_n0`$ such that
$$x_n_{\mathrm{\Lambda }(\phi )}C_2x_n_F\text{for}n=1,2,..$$
$`(1)`$
By condition (A), for any $`0<\epsilon <1,`$ there exists an $`h>0`$ such that
$$\psi (t)\epsilon \phi (t)$$
$`(2)`$
for all positive $`t<h.`$ Choose $`N`$ so that, for $`nN,`$ $`\mu (g_n)<h,`$ where $`g_n=\{t[0,1]:x_n(t)0\}.`$
The subset of finite-valued functions is dense in any Lorentz space on $`[0,1]`$ \[8, p.149\]. Therefore, for each $`nN,`$ there exists a function
$$y_n(t)=\underset{k=1}{\overset{m_n}{}}a_k^n\chi _{e_k^n},\text{where}a_k^n0,e_1^ne_2^n\mathrm{}e_{m_n}^n,\text{and}\mu (e_1^n)<h,$$
for which
$$\mathrm{max}(x_ny_n_{\mathrm{\Lambda }(\phi )},x_ny_n_{\mathrm{\Lambda }(\psi )})<\epsilon \mathrm{min}(x_n_{\mathrm{\Lambda }(\phi )},x_n_{\mathrm{\Lambda }(\psi )}).$$
$`(3)`$
Hence $`x_n_{\mathrm{\Lambda }(\psi )}y_n_{\mathrm{\Lambda }(\psi )}\epsilon x_n_{\mathrm{\Lambda }(\psi )},`$ and \[8, p.160\] and (2) imply that
$$\begin{array}{c}x_n_Fx_n_{\mathrm{\Lambda }(\psi )}1/(1\epsilon )y_n_{\mathrm{\Lambda }(\psi )}=\mathrm{\hspace{0.17em}1}/(1\epsilon )\underset{k=1}{\overset{m_n}{}}a_k^n\psi (\mu (e_k^n))\hfill \\ \hfill \epsilon /(1\epsilon )\underset{k=1}{\overset{m_n}{}}a_k^n\phi (\mu (e_k^n))=\epsilon /(1\epsilon )y_n_{\mathrm{\Lambda }(\phi )}.\end{array}$$
In addition, by (3),
$$y_n_{\mathrm{\Lambda }(\phi )}(1+\epsilon )x_n_{\mathrm{\Lambda }(\phi )}.$$
Thus,
$$x_n_F\frac{\epsilon (1+\epsilon )}{1\epsilon }x_n_{\mathrm{\Lambda }(\phi )}.$$
This inequality with an $`\epsilon >0`$ satisfying
$$\frac{1\epsilon }{\epsilon (1+\epsilon )}>C_2$$
contradicts (1). โ
###### Corollary 1.
Suppose that $`\phi G`$, $`\psi G`$, and $`\psi (t)C_1\phi (t)`$ for $`t(0,1].`$ The following conditions are equivalent:
(1) (A) holds;
(2) the inclusion $`I:\mathrm{\Lambda }(\phi )\mathrm{\Lambda }(\psi )`$ is a DSS operator;
(3) there exist no sequence of nonzero disjoint functions $`\{x_n\}`$ and no constant $`C_2>0`$ such that
$$x_n_{\mathrm{\Lambda }(\phi )}C_2x_n_{\mathrm{\Lambda }(\psi )}\text{for}n=1,2,\mathrm{}$$
$`(4)`$
###### Proof.
The implications $`(1)(2)`$ and $`(2)(3)`$ follow from Theorem 1 and the definition of a DSS operator, respectively.
Suppose that condition (A) does not hold, i.e., that
$$\underset{t0}{lim\; sup}\frac{\psi (t)}{\phi (t)}>\mathrm{\hspace{0.17em}0}.$$
Then there exist a sequence $`\{t_k\}(0,1]`$ and a constant $`C_2>0`$ such that we have $`_{n=1}^{\mathrm{}}t_n1`$ and $`\phi (t_n)C_2\psi (t_n)`$ for $`n=1,2,\mathrm{}`$ The functions $`x_n=\chi _{e_n},`$ where $`e_n[0,1]`$ are disjoint and $`\mu (e_n)=t_n,`$ satisfy (4). Therefore, (3) implies (1); this completes the proof of Corollary 1. โ
Let $`\theta `$ be a function from $`G.`$ The Marcinkiewicz space $`M(\theta )`$ consists of all functions $`x=x(s)`$ measurable on $`[0,1]`$ and such that
$$x_{M(\theta )}=\underset{0<t1}{sup}\frac{1}{\theta (t)}_0^tx^{}(s)๐s<\mathrm{}.$$
The fundamental function of the space $`M(\theta )`$ is $`f_{M(\theta )}(t)=\stackrel{~}{\theta }(t)=t/\theta (t).`$
###### Theorem 2.
Let the functions $`\phi G`$ and $`\psi G`$ satisfy condition (A). If $`E`$ is an SS on $`[0,1]`$ with fundamental function $`\phi (t),`$ then $`EM(\stackrel{~}{\psi })`$ and the identity inclusion $`I:EM(\stackrel{~}{\psi })`$ is a DSS operator.
###### Proof.
Since any Marcinkiewicz space is maximal among all symmetric spaces with the same fundamental function \[8, p.162\], we have $`EM(\stackrel{~}{\psi }).`$
By condition (A), $`\psi (t)C_1\phi (t)`$ for $`t(0,1];`$ therefore, by the definition of a Marcinkiewicz space, $`M(\stackrel{~}{\phi })M(\stackrel{~}{\psi }).`$ Hence, $`EM(\stackrel{~}{\psi }).`$
Suppose that $`I:EM(\stackrel{~}{\psi })`$ is not a DSS operator. Then, in particular, there exists a sequence of disjoint functions $`x_n`$ such that
$$x_n_{M(\stackrel{~}{\psi })}=1\text{and}x_n_{M(\stackrel{~}{\phi })}C_2\text{for}n=1,2,\mathrm{}$$
$`(5)`$
Choose a $`t_k(0,1]`$ for which
$$_0^{t_k}x_k^{}(s)ds\frac{1}{2}\stackrel{~}{\psi }(t_k)(k=1,2,..).$$
Since the functions $`x_k`$ are disjoint, we can assume that $`t_k0.`$ Therefore,
$$x_k_{M(\stackrel{~}{\phi })}\frac{1}{\stackrel{~}{\phi }(t_k)}_0^{t_k}x_k^{}(s)๐s\frac{\phi (t_k)}{2\psi (t_k)}.$$
By (A), $`x_k_{M(\stackrel{~}{\phi })}\mathrm{}`$ as $`k\mathrm{},`$ which contradicts condition (5).
This completes the proof of Theorem 2. โ
###### Corollary 2.
Suppose that $`\phi G`$, $`\psi G`$, and $`\psi (t)C_1\phi (t)`$ for $`t(0,1].`$ The following conditions are equivalent:
(1) (A) holds;
(2) the inclusion $`I:M(\stackrel{~}{\phi })M(\stackrel{~}{\psi })`$ is a DSS operator;
(3) there exist no sequence of nonzero disjoint functions $`x_n`$ and no constant $`C_2>0`$ such that for some $`C_2>0`$
$$x_n_{M(\stackrel{~}{\phi })}C_2x_n_{M(\stackrel{~}{\psi })}\text{for}n=1,2,\mathrm{}$$
The proof of Corollary 2 is similar to the proof of Corollary 1.
###### Corollary 3.
For an arbitrary SS $`EL_1`$ on $`[0,1],`$ the inclusion $`I:EL_1`$ is a DSS operator.
###### Proof.
First, $`L_1=M(1)`$ and an arbitrary SS $`E`$ is embedded in $`L_1`$ \[8, p.124\]. If $`f_E(t)=\phi (t),`$ then the function $`t/\phi (t)`$ increases, because $`\phi `$ is concave. Therefore, condition (A) is violated if and only if $`\phi (t)t`$ (i.e., if and only if $`C_1t\phi (t)C_2t`$ for some $`C_1>0`$ and $`C_2>0`$), and $`E=L_1.`$ It remains to apply Theorem 2. โ
###### Remark 1.
The assertion of Corollary 3 was proved in in a different way.
###### Remark 2.
Arguing as in the proof of Corollary 3 and applying Theorem 1, we can readily show that the inclusion $`I:L_{\mathrm{}}E`$ is a DSS operator for any SS $`EL_{\mathrm{}}.`$ Moreover, it is shown in that this operator is even strictly singular. This generalizes the Grothendieck theorem mentioned in the introduction.
In the next section we show that, generally, condition (A) is not sufficient for the inclusion of an SS with fundamental function $`\phi `$ into an SS with fundamental function $`\psi `$ to have the DSS property.
$`๐\mathrm{\hspace{0.17em}2}.`$ An example of symmetric spaces $`E`$ and $`F`$ such that $`EF`$
and their fundamental functions satisfy condition (A),
but $`I:EF`$ is not a DSS operator
###### Theorem 3.
There exist two symmetric spaces $`E`$ and $`F`$ on $`[0,1]`$ with fundamental functions $`\phi `$ and $`\psi ,`$ respectively, such that $`EF,`$ $`\phi `$ and $`\psi `$ satisfy condition (A), and the operator $`I:EF`$ has not the DSS property.
###### Proof.
Let the SS $`E`$ be the Marcinkiewicz space $`M(\stackrel{~}{\psi })`$ with $`\stackrel{~}{\psi }(t)=t/\psi (t),`$ where
$$\psi (t)=t^{1/2}\mathrm{log}_2^{1/2}\frac{4}{t},\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}<t1.$$
It is readily verified that $`\psi `$ is an increasing concave function on $`[0,1]`$ and $`\gamma _\psi =\delta _\psi =1/2.`$ Therefore, by \[8, p.156\],
$$x_{M(\stackrel{~}{\psi })}\underset{0<t1}{sup}\{x^{}(t)\psi (t)\}.$$
$`(6)`$
Let us define the space $`F`$. We put $`b_k=(k+2)^{1/2}2^{k/2}`$ and $`z_k(t)=b_k\chi _{(0,2^k]}(t)`$ and define a sequence of numbers $`n_0=1<n_1<n_2<\mathrm{}<n_m<\mathrm{}`$ by setting
$$n_{m+1}=\mathrm{max}\{n=1,2,..:\underset{k=n_m}{\overset{n1}{}}\frac{1}{k+2}1\}$$
$`(7)`$
and a sequence of functions $`w_m=w_m(t)`$ by setting
$$w_m(t)=\underset{n_mk<n_{m+1}}{\mathrm{max}}z_k(t),m=0,1,..$$
Since $`\{b_k\}`$ increases, (7) implies that the norms of $`w_m`$ in $`L_2`$ satisfy the inequalities:
$$w_m_2^2\underset{k=n_m}{\overset{n_{m+1}1}{}}b_k^22^{k1}=\mathrm{\hspace{0.17em}1}/2\underset{k=n_m}{\overset{n_{m+1}1}{}}\frac{1}{k+2}\frac{1}{4}$$
and
$$w_m_{2}^{}{}_{}{}^{2}\underset{k=n_m}{\overset{n_{m+1}1}{}}b_k^22^k=\underset{k=n_m}{\overset{n_{m+1}1}{}}\frac{1}{k+2}\mathrm{\hspace{0.17em}1}.$$
Therefore,
$$\frac{1}{2}w_m_21.$$
$`(8)`$
Consider $`\chi _b=b^{1/2}\chi _{(0,b)},`$ $`\overline{w}_m=w_m/w_m_2,`$ and $`V=\{\chi _b\}_{0<b1}\{\overline{w}_m\}_{m=0}^{\mathrm{}}.`$ Let $`F`$ be the set of all functions $`x=x(t)`$ measurable on $`[0,1]`$ and satisfying
$$x=sup\{_0^1x^{}(t)v(t)๐t:vV\}<\mathrm{}.$$
Then $`F`$ is an SS on $`[0,1]`$ as the intersection of Lorentz spaces determined by the functions $`_0^tv(s)๐s`$ with $`vV.`$ In addition, the definition of $`F`$ implies that $`x_{M(t^{1/2})}x_Fx_2.`$ Therefore the fundamental functions $`f_E(t)=\psi (t)`$ and $`f_F(t)=t^{1/2}`$ of the spaces $`E`$ and $`F`$ satisfy condition (A).
Let us prove that
$$M(\stackrel{~}{\psi })F.$$
$`(9)`$
By (6), it is suffices to show that $`1/\psi F.`$ Indeed,
$$_0^1\chi _b(t)\frac{dt}{\psi (t)}=\mathrm{\hspace{0.17em}2}b^{1/2}_0^b\frac{d(t^{1/2})}{\mathrm{log}_2^{1/2}4/t}2\text{for}\mathrm{\hspace{0.33em}\hspace{0.33em}0}<b1,$$
$$_0^1w_m(t)\frac{dt}{\psi (t)}\underset{k=n_m}{\overset{n_{m+1}1}{}}b_k_0^{2^k}\frac{dt}{\psi (t)}=\mathrm{\hspace{0.17em}2}\underset{k=n_m}{\overset{n_{m+1}1}{}}b_k_0^{2^k}\frac{d(t^{1/2})}{\mathrm{log}_2^{1/2}4/t}$$
$$\mathrm{\hspace{0.17em}2}\underset{k=n_m}{\overset{n_{m+1}1}{}}\frac{1}{k+2}\mathrm{\hspace{0.17em}2}.$$
Therefore (8) and the definition of $`F`$ imply that $`1/\psi _F4,`$ which proves (9).
Next, we put $`D_m=(2^{n_{m+1}},2^{n_m}]`$ and
$$v_m(t)=w_m(t)\chi _{D_m}(t)=\underset{k=n_m}{\overset{n_{m+1}1}{}}b_k\chi _{(2^{k1},2^k]}(t)\text{for}m=0,1,\mathrm{}$$
The functions $`v_m`$ are disjoint. Let us show that the norms of $`E`$ and $`F`$ are equivalent on their linear hull.
Suppose that
$$v(t)=\underset{m=0}{\overset{r}{}}a_mv_m(t).$$
Without loss of generality, we can assume that $`a_m0.`$ Consider $`w(t)=\mathrm{max}_{0mr}a_mw_m(t).`$ The function $`w(t)`$ monotonically decreases on $`(0,1],`$ and $`v(t)w(t).`$ Therefore, by (6),
$$v_Ew_EC\underset{0mr}{\mathrm{max}}\left\{a_m\underset{n_mk<n_{m+1}}{\mathrm{max}}b_k\psi (2^k)\right\}.$$
Since $`b_k\psi (2^k)=1`$ for $`k=0,1,2,..,`$ we obtain
$$v_EC\underset{0mr}{\mathrm{max}}a_m.$$
$`(10)`$
Now, let us estimate $`v_F`$ from below. By (7), for any $`m=0,1,..,r`$ we have
$$\begin{array}{c}_0^1v_m^{}(t)w_m(t)๐t_0^1(v_m^{}(t))^2๐t=_0^1v_m^2(t)๐t=\hfill \\ \hfill =\underset{k=n_m}{\overset{n_{m+1}1}{}}b_k^22^{k1}=\frac{1}{2}\underset{k=n_m}{\overset{n_{m+1}1}{}}\frac{1}{k+2}\frac{1}{4}.\end{array}$$
Hence (8) implies that $`v_m_F1/4.`$ Therefore, by (10),
$$v_F\underset{0mr}{\mathrm{max}}\{a_mv_m_F\}\frac{1}{4}\underset{0mr}{\mathrm{max}}a_m\frac{v_E}{4C}.$$
This together with (9) means that the norms of the spaces $`E`$ and $`F`$ are equivalent on the linear hull of the set of functions $`v_m`$ $`(m=0,1,..).`$ Hence there exists a $`B>0`$ such that, for an arbitrary $`a_m,`$
$$B^1\underset{m=0}{\overset{\mathrm{}}{}}a_mv_m_F\underset{m=0}{\overset{\mathrm{}}{}}a_mv_m_EB\underset{m=0}{\overset{\mathrm{}}{}}a_mv_m_F.$$
In other words, the identity inclusion operator $`I:EF`$ has not the DSS property.
This completes the proof of Theorem 3. โ
###### Remark 3.
Theorem 3 shows that relation (A) does not guarantee the presence of a โgapโ between spaces sufficient for the corresponding identity inclusion operator to have the DSS property. On the other hand, simple examples show that, even for symmetric spaces with the same fundamental function, this operator may have this property.
For instance, the inclusion of the Lorentz space $`\mathrm{\Lambda }(t^{1/p})`$ into the space $`L_p,`$ where $`1<p<\mathrm{},`$ has the DSS property. Indeed, it is easy to show that any sequence of normalized disjoint functions in the Lorentz space contains a subsequence equivalent to the standard basis in $`l_1`$. At the same time, any such sequence in $`L_p`$ is equivalent to the standard basis in $`l_p.`$
In the last section we show that, unlike (A), condition (B) is sufficient for the inclusion operator $`I:EF`$ of arbitrary symmetric spaces $`E`$ and $`F`$ with fundamental functions $`\phi `$ and $`\psi ,`$ respectively, to be a DSS operator.
$`๐\mathrm{\hspace{0.17em}3}.`$ Suffiiency of condition (B) for the operator
$`I:EF`$ to have the DSS property
###### Theorem 4.
Suppose that functions $`\phi G`$ and $`\psi G`$ satisfy condition (A), $`\delta _\phi <1,`$ and we have $`M(\stackrel{~}{\phi })\mathrm{\Lambda }(\psi ).`$ Then the operator $`I:M(\stackrel{~}{\phi })\mathrm{\Lambda }(\psi )`$ has the DSS property.
First, we prove the following auxiliary assertion.
###### Lemma 1.
Under the assumptions of Theorem 4, there exists a function $`\rho G`$ such that
$$(1)\underset{t0}{lim}\frac{\rho (t)}{\phi (t)}=\mathrm{\hspace{0.17em}0};$$
$$(2)M(\stackrel{~}{\rho })\mathrm{\Lambda }(\psi ).$$
###### Proof.
Since $`\delta _\phi <1,`$ \[8, p.156\] implies that
$$x_{M(\stackrel{~}{\phi })}\underset{0<t1}{sup}\{\phi (t)x^{}(t)\}.$$
Therefore the relation $`M(\stackrel{~}{\phi })\mathrm{\Lambda }(\psi )`$ is equivalent to
$$_0^1\frac{d\psi (s)}{\phi (s)}<\mathrm{}.$$
$`(11)`$
Since the function $`\phi `$ is concave, we have
$$_0^1\frac{d\psi (s)}{\phi (s)}\underset{k=0}{\overset{\mathrm{}}{}}\frac{\psi (2^k)\psi (2^{k1})}{\phi (2^k)};$$
hence (11) is equivalent to the condition
$$\underset{k=0}{\overset{\mathrm{}}{}}\frac{\psi (2^k)\psi (2^{k1})}{\phi (2^k)}<\mathrm{}.$$
$`(12)`$
Put
$$a_k=\psi (2^k)\psi (2^{k1})\text{and}S_n=\underset{k=n}{\overset{\mathrm{}}{}}\frac{a_k}{\phi (2^k)}.$$
Then $`S_n0,`$ and \[12, Chap. 3, Ex. 12\] and (12) imply that
$$\underset{k=0}{\overset{\mathrm{}}{}}\frac{a_k}{\sqrt{S_k}\phi (2^k)}<\mathrm{}$$
$`(13)`$
By the definition of upper dilation index, there exist $`u>0`$ and $`C>0`$ such that $`\delta _\phi +u<1`$ and
$$_\phi (t)Ct^{\delta _\phi +u/2}$$
$`(14)`$
for all $`t1.`$ Consider the sequence of numbers
$$g_0=S_0,g_k=\mathrm{max}(S_k,2^ug_{k1})\text{for}k=1,2,\mathrm{}$$
$`(15)`$
Suppose that $`g=g(t)`$ is a function linear on the intervals $`[2^{k1},2^k],`$ $`g(2^k)=g_k`$ for $`k=0,1,..,`$ and $`h(t)=\sqrt{g(t)}\phi (t).`$
Since $`\{S_k\}`$ decreases, $`\{g_k\}`$ also decreases; therefore, the functions $`g(t)`$ and $`h(t)`$ increase on $`(0,1].`$ It follows from (15) that, for $`j0`$ and $`2^{k1}<t2^k,`$
$$\frac{g(2^jt)}{g(t)}\frac{g(2^{jk})}{g(2^{k1})}=\frac{g_{kj}}{g_{k+1}}=\frac{g_{k+1(j+1)}}{g_{k+1}}2^{(j+1)u}.$$
Hence, by (14),
$$_h(2^j)=\underset{0<t2^j}{sup}\frac{\phi (2^jt)\sqrt{g(2^jt)}}{\phi (t)\sqrt{g(t)}}C2^{u/2}2^{(\delta _\phi +u)j}\text{for}j=0,1,\mathrm{}$$
This implies that $`\delta _h<1,`$ because $`\delta _\phi +u<1.`$ Therefore, according to \[8, p.78\], the function $`h(t)`$ is equivalent to its least concave majorant; we denote this majorant by $`\rho (t)`$.
The function $`\rho `$ belongs to $`G,`$ and $`\rho (2^k)h(2^k)=\sqrt{g_k}\phi (2^k)`$ for $`k=0,1,2,..;`$ in addition, since $`g_kS_k,`$ it follows from (13) that
$$\underset{k=0}{\overset{\mathrm{}}{}}\frac{a_k}{\rho (2^k)}C_1\underset{k=0}{\overset{\mathrm{}}{}}\frac{a_k}{\sqrt{S_k}\phi (2^k)}<\mathrm{}.$$
This equivalent to
$$M(\stackrel{~}{\rho })\mathrm{\Lambda }(\psi );$$
the equivalence is proved in the same way as for the function $`\phi `$ (see (11) and (12)).
Finally, relation (15) gives
$$\underset{t0}{lim}\frac{\rho (t)}{\phi (t)}C_1\underset{t0}{lim}\sqrt{g(t)}=\underset{k\mathrm{}}{lim}\sqrt{g_k}=\mathrm{\hspace{0.17em}0}.$$
This completes the proof of Lemma 1. โ
###### Proof of Theorem 4..
By Lemma 1, there exists a function $`\rho G`$ such that
$$M(\stackrel{~}{\phi })M(\stackrel{~}{\rho })\mathrm{\Lambda }(\psi )\text{and}\underset{t0}{lim}\frac{\rho (t)}{\phi (t)}=\mathrm{\hspace{0.17em}0}.$$
According to Corollary 2, the operator $`I:M(\stackrel{~}{\phi })M(\stackrel{~}{\rho })`$ has the DSS property; all the more, it has this property when regarded as an operator from $`M(\stackrel{~}{\phi })`$ into $`\mathrm{\Lambda }(\psi )`$.
This proves Theorem 4. โ
Theorem 4 makes it possible to prove the sufficiency of condition (B) for the identity inclusion operator from an SS $`E`$ into an SS $`F`$ with fundamental functions $`\phi `$ and $`\psi ,`$ respectively, to have the DSS property.
###### Theorem 5.
Suppose that functions $`\phi G`$ and $`\psi G`$ satisfy condition (B) and $`E`$ and $`F`$ are symmetric spaces with fundamental functions $`\phi `$ and $`\psi ,`$ respectively. Then $`EF`$ and $`I:EF`$ is a DSS operator.
###### Proof.
Let us verify that the assumptions of Theorem 4 hold, i.e., that $`\delta _\phi <1`$ and that
$$M(\stackrel{~}{\phi })\mathrm{\Lambda }(\psi ).$$
$`(16)`$
First, condition (B) implies the existence of $`u>0`$ and $`C>0`$ such that
$$\frac{\psi (ts)\phi (s)}{\psi (s)\phi (ts)}Ct^u$$
$`(17)`$
whenever $`0<t1`$ and $`0<s1.`$ This and the concavity of $`\psi `$ give
$$_\phi \left(\frac{1}{t}\right)C_\psi \left(\frac{1}{t}\right)t^u=Ct^{u1};$$
therefore, $`\delta _\phi 1u<1.`$
Next, since the function $`x^{}(t)`$ decreases, we have
$$x^{}(t)\frac{x_{M(\stackrel{~}{\phi })}}{\phi (t)}\text{if}xM(\stackrel{~}{\phi })\text{and}\mathrm{\hspace{0.33em}\hspace{0.33em}0}<t1.$$
Therefore, to prove (16), it suffices to verify that $`1/\phi \mathrm{\Lambda }(\psi ).`$
By (17), $`\psi (t)/\phi (t)C_1t^u.`$ Hence $`\psi (t)0`$ as $`t0,`$ and
$$1/\phi _{\mathrm{\Lambda }(\psi )}=_0^1\frac{\psi ^{}(t)}{\phi (t)}๐t_0^1\frac{\psi (t)}{\phi (t)}\frac{dt}{t}\frac{C_1}{u}<\mathrm{}.$$
This proves (16). The above-mentioned extremality of the Lorentz and Marcinkiewicz spaces in the class of symmetric spaces with the same fundamental function implies that $`EM(\stackrel{~}{\phi })\mathrm{\Lambda }(\psi )F.`$ Therefore, $`I:EF`$ has the DSS property, because it has this property as an operator from $`M(\stackrel{~}{\phi })`$ into $`\mathrm{\Lambda }(\psi )`$ by Theorem 4.
This completes the proof of Theorem 5. โ
###### Remark 4.
Theorem 5 was proved in for Orlicz spaces under a condition on functions $`\phi `$ and $`\psi `$ somewhat more restrictive than (B), namely, the inequality $`\delta _\phi <\gamma _\psi .`$
The author wishes to express his gratitude to S. Ya. Novikov for useful discussiones.
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# GUTPA/00/04/02WUP-TH 00-13HLRZ2000-11 A high precision study of the ๐โข๐ฬ potential from Wilson loops in the regime of string breaking
## I Introduction
Confinement of quarks is an issue of prime importance in the understanding of strong interaction physics. While the study of the static quark-antiquark potential from simulations of quantum chromodynamics has been pushed to rather high accuracy in quenched QCD and allows today for a precise determination of the string tension, the search for evidence of string breaking from Wilson loops in simulations of the full QCD vacuum has been futile so far. It seems that the linear rise of the static potential continues to prevail even in presence of vacuum polarization by quark loops .
The common explanation for this unexpected finding is that present studies cannot really resolve the asymptotic time behaviour of the Wilson loops and that a multichannel analysis including light fermion operators is required to achieve sufficient ground state overlap in the available time range . In fact, a fully fledged multichannel approach has been demonstrated to be a viable technique to realize breaking of the string between adjoint sources in pure gauge theories or fundamental colour sources in Higgs models . In the case of full QCD, however, it appears overly costly to achieve the required statistical accuracy of the generalized Wilson loops which incorporate light quark-antiquark pairs in the initial or final states . The reason is that one is prevented from exploiting the translational symmetry, since this would require computation of light propagators $`P(y,x)`$ on any source point location, $`x`$ (see e.g. Refs. ). In Refs. stochastic estimator methods with maximal variance reduction (so-called all-to-all methods) were applied to cope with the fluctuations on the multichannel correlator matrix, but failed so far to provide sufficient accuracy in the infrared regime.
On the other hand, if one scrutinizes existing QCD potential data from Wilson loops for colour screening one will notice that the errors become substantial in the region of interest, $`r1.2`$ fm<sup>*</sup><sup>*</sup>*We use capital (lower case) letters for quantities in lattice units (physical units).. Thus, there is room for suspicion that QCD colour screening has so far escaped detection simply for the lack of adequate precision of large Wilson loops.
The main purpose of the present note is to improve on this point by pushing for a high precision โclassicalโ Wilson loop calculation at large separations in $`N_f=2`$ QCD. This is achieved by squeezing maximal information out of each vacuum configuration through rotational invariance and comprehensive utilization of all possible $`R`$-values on the lattice. As a result of our โall-$`R`$ approachโ (ARA) we are able to present a long range static potential from Wilson loops in unprecedented accuracy.
In section II we shall describe how we go about in the build-up of nonplanar loops to any given $`๐`$. Section III contains our potential analysis from ARA which is used on top of existing signal enhancement techniques, such as conventional APE smearing and translational averaging. As a first step in the direction of a two channel investigation of string breaking we study in section IV the noise reduction effect from ARA on the transition correlator between static and static-light quark states, $`\overline{Q}Q`$ and $`\overline{Q}Q\overline{q}q`$, respectively.
## II Loop construction
Our aim is to increase statistics in the regime of colour screening, i.e. for large quark antiquark separations, $`R`$. Obviously, on such a length scale, a given QCD vacuum configuration contains plenty of information that can be exploited for self averaging and thus for error reduction: firstly, one can realize, on a hypercubic lattice, a fairly dense set of $`R`$ values; secondly, for a given large value of $`R=|๐|`$, there are many different three-vector realizations $`๐`$ on the lattice.
We wish to make use of this fact by a systematic construction of off-axis Wilson loops, $`W(R,T)`$, in the range $`R_{\mathrm{min}}RR_{\mathrm{max}}`$, with $`R_{\mathrm{min}}=101.7R_0`$ and $`R_{\mathrm{max}}=12\sqrt{3}3.5R_0`$, where $`R_0`$ is the Sommer radius (in lattice units) that amounts to $`r_00.5`$ fm . The construction proceeds by choosing all possible vectors, $`๐`$, with integer components $`C_{\mathrm{min}}`$, $`C_{\text{mid}}`$, and $`C_{\mathrm{max}}`$ (in any order of appearance) that obey the inequalities,
$$R_{\mathrm{min}}^2R^2=C_{\mathrm{min}}^2+C_{\text{mid}}^2+C_{\mathrm{max}}^2R_{\mathrm{max}}^2,$$
(1)
where $`|C_{\mathrm{min}}||C_{\text{mid}}||C_{\mathrm{max}}|`$.
Subsequently, the set of solutions to the constraint eq. (1) is sorted according to the correponding values of $`R`$. In Table I we display the large number of possible $`R`$ values and vectors, $`๐`$, obtained in this way, with the additional restriction, $`|C_{\mathrm{max}}|12`$. So far, only $`๐`$ vectors with $`(|C_{\mathrm{max}}|,|C_{\text{mid}}|,|C_{\mathrm{min}}|)`$ being multiples of $`(1,0,0),(1,1,0),(1,1,1),(2,1,0),(2,1,1)`$ or $`(2,2,1)`$ have been considered . Within the investigated regime, $`10R12\sqrt{3}`$, we achieve a gain factor of more than eight in terms of the spatial resolution in $`R`$ (from 21 to 175 different values). Moreover, the number of different $`๐`$ vectors that yield the same distance, $`R`$, is increased by an average gain factor of more than four! Accordingly, we find the ARA to reduce the statistical errors on potential values by factors around two.
We shall briefly discuss the construction of the gauge transporters connecting the quark and antiquark locations, $`๐_Q`$ and $`๐_{\overline{Q}}=๐_Q+๐`$, that appear within the ARA nonplanar Wilson loops. In order to achieve a large
overlap with the physical ground state we would like to construct lattice paths that follow as close as possible the straight line connecting $`๐_{\overline{Q}}`$ with $`๐_Q`$. This task can be accomplished by a procedure which is known as the Bresenham algorithm in computer graphics. There one wishes to map a straight continuous line between two points, say $`\mathrm{๐}`$ and $`๐=(C_{\mathrm{max}},C_{\mathrm{min}})`$, onto discrete pixels on a 2-d screen. Then one has to find the explicit sequence of pixel hoppings in max- and min-directions such that the resulting pixel set mimicks best the continuum geodesic between points $`\mathrm{๐}`$ and $`๐`$. The Bresenham prescription is simply to move in max-direction unless a step in min-direction brings you closer to the geodesic from $`\mathrm{๐}`$ to $`๐`$, where we assume $`|C_{\mathrm{max}}||C_{\mathrm{min}}|`$. It can easily be embodied into a fast algorithm based on local decision making only (see Fig. 1), by means of a characteristic lattice function $`\chi `$ that incorporates the aspect ratio $`C_{\mathrm{max}}/C_{\mathrm{min}}`$.
The algorithm in two dimensions looks like this:
cmax2 := 2\*cmax
cmin2 := 2\*cmin
chi := cmin2 -cmax
FOR i := 1 TO cmax DO
> step in max-direction
> IF chi $``$ 0 THEN
>
> > chi := chi - cmax2
> > step in min-direction
>
> ENDIF
> chi := chi + cmin2
ENDDO
The generalisation to three dimensions is achieved by combining two of these 2d-algorithms with different $`\chi `$โs for max-mid and max-min in just one loop over the max-direction.
In order to convey an idea about the statistics gain inherent in such a systematic approach we have listed the number of $`๐`$ vectors constructed in this manner in Table I. Note that for plane or space diagonal separations, we average over the two or six equivalent paths, respectively.
## III The static potential at large $`R`$
We base our analysis on 184 vacuum configurations, separated by one autocorrelation length, on $`L_\sigma ^3L_\tau =24^3\times 40`$ lattices of $`N_f=2`$ QCD at $`\beta =5.6`$ and $`\kappa _{sea}=.1575`$ (corresponding to $`m_\pi /m_\rho =.704(5)`$), produced by the T$`\chi `$L-collaboration on an APE 100 tower at INFN. These parameter values correspond to the biggest lattice volume at our disposal, $`L_\sigma a2`$ fm. In order to minimize finite-size effects and possible violations of rotational symmetry on the $`L_\sigma =24`$ torus, $`C_{\mathrm{max}}`$ has been restricted to $`|C_{\mathrm{max}}|12`$. The lattice constant $`a`$ was determined from the Sommer radius $`r_0=R_0a0.5`$ fm , as obtained in our previous investigation , $`R_0=5.89(3)`$.
Before entering the Wilson loop analysis the configurations are smoothened by spatial APE- or link smearing ,
$$\text{link}\alpha \times \text{link}+\text{staples},$$
(2)
followed up by a projection back into the gauge group , with 26 such iterative replacements and the parameter value, $`\alpha =2.3`$.
The potential values have been obtained by means of single and double exponential fits to smeared ARA Wilson loop data $`W(R,T)`$ within the range, $`T_{\mathrm{min}}T8`$. We shall quote statistical errors that are obtained by jackknifing. At the large $`R`$ values that are of interest in view of screening effects, the quality of the statistical signal did not allow to include $`T`$ values larger than $`8`$.
In Fig. 2 we display estimates on the potential in the range $`1.7r/r_03.5`$ obtained from a single exponential fit with $`T_{\mathrm{min}}=4`$ (they agree with results from a double exponential fit with $`T_{\mathrm{min}}=1`$). Note that for $`r>2.4r_0`$ the data are to be interpreted as strict upper limits on the potential. The slope is in agreement with the string tension, $`K=\sigma a^2=0.0372(8)=1.139(4)R_0^2`$, as quoted in Ref. from a fit to data obtained at smaller $`r<2.04r_0`$. Around the separation $`r_c2.3r_0`$ the potential energy crosses the expected threshold for string breaking, $`2m_{PS}a=1.256(13)`$, which is indicated by a horizontal error band. The errors quoted are statistical and are well below 1.5 % for $`r/r_03`$. We find the data to follow perfectly a straight line: a flattening of the Wilson loop potential is not visible within the accessible $`T`$ range and present statistical errors.
To complement this result, we computed the overlaps with the ground and first excited states by means of two-exponential fits, as displayed in Fig. 3. Again our two data sets exhibit linear dependencies on $`r`$, with nearly opposite slopes such that their sum turns out to be close to one. The remainder does only slightly depend on $`r`$ and is of order 10 % .
We conclude that Wilson loop operators are definitely not very well suited to uncover string breaking. To achieve the required overlap with four-quark states one has to introduce them explicitely into the calculation, in form of a coupled channel analysis. As a first step in such a program, we shall explore in the remaining part of this letter a number of signal enhancement techniques based on ARA, in application to the transition matrix element between two- and four-quark states.
## IV Noise on the transition operator
In a two-channel approach one extends the Wilson loop vacuum expectation value, $`C_{11}`$, to a $`2\times 2`$ correlation matrix $`C`$ as pictogrammed in Fig. 4. The task is then to solve the generalized eigenvalue problem ,
$$C(t+\tau )๐ฎ_i=\lambda _i(\tau )C(t)๐ฎ_i,$$
(3)
where the eigenvalues connect to the two energy levels $`E_i`$ at large enough $`t`$:
$$\lambda _i(\tau )=\mathrm{exp}(E_i\tau ).$$
(4)
Unless $`C`$ happens to be diagonal, the physical states $`๐ฎ_0`$ and $`๐ฎ_1`$ are mixtures of two- and four-quark states. It is obvious that the energy levels are insensitive to changes in the normalization of the wavefunctions, $`๐ฎ_i`$. Thus, in addition to spatial APE link smearing, source smearing techniques as known from spectrum calculations are applicable (see Fig. 5) and will be tested.
While the noise level on $`C_{11}`$ could be greatly reduced by applying ARA on top of standard volume self averaging (VSA), the quark propagators $`P(y,x)`$ (wiggly lines) entering the remaining components of the correlation matrix $`C`$ prevent the direct use of VSA that requires matrix inversions on all possible source points, $`x`$. Pennanen and Michael have tested noisy estimator techniques on $`P`$ for curing this problem but did not succeed to reach the accuracy necessary to solve eq. (3. This motivates us to explore noise controlling strategies based on ARA rather than VSA on the transition matrix element $`C_{21}(R,T)`$. One should keep in mind that ARA can be put to work at little extra cost since one inversion on a single source $`x`$ renders $`P(y,x)`$ for all sink positions $`y`$.
In addition to APE link smearing, we have explored source smearing as illustrated by the diagram, Fig. 5. This is applied to the light propagator source, $`\chi _x`$, and consists of iterative replacement (as in previous heavy-light spectrum analyses ),
$$\chi _x\chi _x+\alpha \underset{i=\pm 1}{\overset{\pm 3}{}}\chi _{x+\widehat{ฤฑ}}^{\text{parall.transported}},$$
(5)
which is repeated 50 times, with weight factor, $`\alpha =4.0`$. The source is put to zero subsequently outside the volume $`|๐ซ_s๐ฑ|R_qa`$, with $`R_q=5`$. Note that quark propagators are computed without link smearing throughout this work.
Fig. 6 shows the $`R`$-dependence of the relative errors on $`C_{21}(R,T=1)`$ throughout the stringbreaking region and illustrates the performance of various signal enhancement tools that are added on top of ARA: (a) the circles correspond to no source or link smearing; (b) crosses refer to local sources and smeared links; (c) stars refer to smeared quark sources and smeared gauge sinks.
Obviously link smearing helps a lot but leaves us still with errors of order 50 % in the region $`R18`$. It turns out that smearing is capable of cutting down noise amplitudes further. We find an additional reduction of more than a factor two, to a 20 % level at $`T=1`$.
Unfortunately, such accuracy does not suffice, as it cannot be sustained at larger $`T`$ values, where one wishes to analyse $`C`$ at the end of the day. Given the dense set of $`R`$-values available from ARA and in view of the fact that $`B(R)=\mathrm{ln}[C_{21}(R,T)]`$ is a smooth function of $`R`$, there is opportunity to further improve by filtering the sequence, $`\{B(R_i)\}\{B_f(R_i)\}`$. We have chosen a filter that weighs each individual jackknife sample with the fluctuation calculated on the entire data set,
$$B_f(R_i)=N_i^1\underset{|R_jR_i|R_f}{}\sigma _j^2B(R_j),$$
(6)
with normalization
$$N_i=\underset{|R_jR_i|R_f}{}\sigma _j^2.$$
(7)
We found best results with filter radius $`R_f=0.5`$. Let us now consider the transition matrix element at $`T=5`$. The effect of filtering is illustrated in Fig. 7 where we confront the signals from ARA plus link and source smearings, $`\{B(R_i)\}`$, (upper sequence with large error bars) with the filtered one, $`\{B_f(R_i)\}`$, (lower data sequence). In order to display the systematics of windowing, we have refrained from thinning the latter sequence as one should in actual applications. We find a striking noise reduction through filtering within the large distance regime, $`10R18`$: at $`T=3`$ and $`T=5`$ we encounter errors of less than 10 % and 20 %, respectively.
## V Summary and Conclusions
Exploiting the dense set of $`R`$ values available at large quark-antiquark separations on the lattice we succeeded in improving the precision of the $`Q\overline{Q}`$ potential from Wilson loops in the string breaking regime by a factor of two with respect to standard methods. This enables us to analyse Wilson loop data well beyond the point where string breaking is expected and corroborates previous conjectures that Wilson loops do not bear enough overlap with $`Q\overline{Q}q\overline{q}`$ states for uncovering string breaking within the $`T`$-range available at present.
The success of our all-$`R`$ approach to the Wilson loops in the string breaking regime has encouraged us to carry out a feasibility study on error control of the transition correlator, $`C_{21}`$. By additional use of source smearing and filtering techniques, we arrive at reasonable signal-to-noise ratios for $`\mathrm{ln}C_{21}`$.
For the final chord in a full two-channel analysis of $`Q\overline{Q}`$ and $`Q\overline{Q}q\overline{q}`$ states one would have to consider the correlator $`C_{22}`$ which contains both, a connected and a disconnected contribution: the former can be tackled with our present techniques and we expect sufficient accuracy with a factor of $`T_{max}`$ more effort than for $`C_{21}`$; the latter relates to the situation of $`B\overline{B}`$ pairs at large separations. From a previous $`B\overline{B}`$ study one would anticipate that the $`B\overline{B}`$ correlator is dominated in our $`R`$-range and for our purposes by $`B\overline{B}`$. This is currently being investigated.
###### Acknowledgements.
We appreciate useful discussions with S. Gรผsken und M. Peardon during early stages of this research. K.S. thanks F. Niedermayer for an interesting discussion. This work was supported by DFG Graduiertenkolleg โFeldtheoretische und Numerische Methoden in der Statistischen und Elementarteilchenphysikโ. G.B. acknowledges support from DFG grants Ba 1564/3-1, 1564/3-2 and 1564/3-3 as well as EU grant HPMF-CT-1999-00353. The HMC productions were run on an APE100 at INFN Roma. We are grateful to our colleagues F. Rapuano and G. Martinelli for the fruitful T$`\chi `$L-collaboration. Analysis was performed on the CRAY T3E system of ZAM at Research Center Jรผlich.
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# REFERENCES
HD-THEP-00-28
Optimisation of the exact renormalisation group
Daniel F. Litim <sup>*</sup><sup>*</sup>*D.Litim@thphys.uni-heidelberg.de
Institut fรผr Theoretische Physik
Philosophenweg 16, D-69120 Heidelberg, Germany.
Abstract
A simple criterion to optimise coarse-grainings for exact renormalisation group equations is given. It is aimed at improving the convergence of approximate solutions of flow equations. The optimisation criterion is generic, as it refers only to the coarse-grained propagator at vanishing field. In physical terms, it is understood as an optimisation condition for amplitude expansions. Alternatively, it can be interpreted as the requirement to move poles of threshold functions away from the physical region. The link to expansions in field amplitudes is discussed as well. Optimal parameters are given explicitly for a variety of different coarse-grainings. As a by-product it is found that the sharp cut-off regulator does not belong to the class of such optimal coarse-grainings, which explains the poor convergence of amplitude expansions based on it.
Introduction
Flow equations or exact renormalisation group equations are a promising tool for non-perturbative problems within quantum field theories. Based on the idea of integrating-out large momentum modes within a path integral formulation of quantum field theory, they permit to obtain an effective (or โcoarse-grainedโ) theory for the remaining light degrees of freedom. Flows interpolate between some initial (classical) action, and the full (quantum) effective action. The particular strength of the formalism is its flexibility, allowing for systematic approximations without being tied to the small coupling region .
On the conceptual side, important progress has been made within the recent past, in particular through the consistent extension to (non-Abelian) gauge theories , which may play an important role in future applications to QCD. In turn, comparatively little attention has been paid to the interplay between approximate solutions to flow equations on the one side and the coarse-graining procedure on the other. Any explicit application of the formalism requires some approximation, simply because it seems very difficult to simultaneously solve infinitely many coupled partial differential equations. Approximate solutions have a finite domain of validity , they may show a spurious dependence on the coarse-graining parameters at second order phase transitions or first order ones , and the corresponding expansions have a finite (often unknown) radius of convergence . It has also been argued that the spurious scheme dependence can be used to assess the reliability of a given approximation , resulting in error bars describing the scheme and truncation dependence , or to obtain a better convergence employing a minimum sensitivity condition .
The cross-dependence between the coarse-graining and a truncation of the flow has a very simple origin. Any coarse-graining couples to all operators in the effective action, with a strength depending on the details of the specific implementation. The sharp cutoff, for example, eliminates all momentum modes above some coarse-graining scale $`k`$, while smooth cutoff functions, like exponential, algebraic or mass-like ones, smoothly cut off the higher modes within some finite momentum interval about $`k`$. Hence, a change of the coarse-graining scheme alters the effective coupling to the operators in the effective action. A strong scheme dependence indicates whether some relevant operators have been neglected within a given approximation . Along similar lines it has also been argued that the smoothness of the coarse-graining can influence the contributions of higher order derivative operators .
More generally, this interplay suggests that a suitably chosen coarse-graining can significantly improve the convergence properties or the domain of validity of approximate solutions, which would be of great value for both analytical and numerical applications. In the present note we take profit of this freedom to propose a simple yet efficient optimisation criterion for Wilsonian flows, with the following characteristics: it optimises approximate solutions to flow equations; it is compatible with systematic approximation schemes like the derivative expansion; it only relies on a basic ingredient to the flow equation (the effective coarse-grained propagator at vanishing field); it is generically applicable, i.e. not based on a specific theory.
Flows, coarse-grainings and propagators
To begin with, we briefly review the basic ingredients needed for a Wilsonian flow equation. The modern way to implement an exact renormalisation group procedure amounts to add a regulator term $`d^dq\varphi (q)R_k(q^2)\varphi (q)`$ (for bosonic fields) to the action . The operator $`R_k(q^2)`$, sometimes referred to as the regulator scheme (RS), desribes the coarse-graining and introduces a fiducial scale parameter $`k`$, which is ultimately interpreted as a coarse-graining scale. This scale parameter induces a scale dependence, which, when written for the scale-dependent effective action $`\mathrm{\Gamma }_k`$, results in the flow equation
$$\frac{}{t}\mathrm{\Gamma }_k[\varphi ]=\frac{1}{2}\mathrm{Tr}\left\{\left(\mathrm{\Gamma }_k^{(2)}[\varphi ]+R_k\right)^1\frac{R_k}{t}\right\}.$$
(1)
Here, $`\varphi `$ denotes bosonic fields and $`t=\mathrm{ln}k`$ the logarithmic scale parameter. The right hand side of (1) contains the regulator function $`R_k`$ and the second functional derivative of the effective action with respect to the fields. The trace denotes a summation over all indices and integration over all momenta.
All specifications regarding the coarse-graining are given through the operator $`R_k(q^2)`$. This operator can be chosen at will though within some basic restrictions which are briefly reviewed. First of all, it is required that $`R_k(q^20)>0`$. This ensures that the effective propagator at vanishing field $`\mathrm{\Delta }_k(q^2)=1/[q^2+R_k(q^2)]`$ remains finite in the infrared limit $`q^20`$, and no infrared divergences are encountered in the presence of massless modes. The second requirement is the vanishing of $`R_k`$ in the infrared, $`R_k(q^2)0`$ for $`k0`$. This guarantees that the coarse-grained generating functional of 1PI Green functions $`\mathrm{\Gamma }_k`$ reduces to the usual generating functional of 1PI Green functions $`\mathrm{\Gamma }=lim_{k0}\mathrm{\Gamma }_k`$. The third condition to be met is that $`R_k(q^2)`$ diverges in the UV limit $`k\mathrm{\Lambda }`$. This way it is ensured that the microscopic action $`S=lim_{k\mathrm{\Lambda }}\mathrm{\Gamma }_k`$ is approached in the ultraviolet limit $`k\mathrm{\Lambda }`$. These conditions ensure that the flow (1) interpolates between the classical and the quantum effective action. For later convenience, we shall impose a further constraint on $`R_k`$, which, however, is nothing but a normalisation condition for the coarse-graining scale $`k`$, namely $`R_k(q^2=k^2)=k^2`$. If one would have chosen a regulator with $`R_k(q^2=k^2)=ck^2`$, then it is always possible to find a corresponding scale $`k_{\mathrm{eff}}=f(c)k`$ such that $`R(q^2=k_{\mathrm{eff}}^2)=k_{\mathrm{eff}}^2`$. Hence, the last condition guarantees that (trivial) rescalings of $`k`$ have been factored-out, and comparing different effective propagators is now sensible.
It is important to realise that the integrand of the flow equation (1), as a function of momenta $`q`$, is peaked about $`q^2k^2`$, and suppressed for large momenta. Consequently, at each infinitesimal integration step $`kk\mathrm{\Delta }k`$ only a narrow window of momentum modes contribute to the change of $`\mathrm{\Gamma }_k\mathrm{\Gamma }_{k\mathrm{\Delta }k}`$. Most importantly, modes with momenta $`qk`$ no longer contribute to the running at the scale $`k`$. It is this property which justifies the interpretation of $`\mathrm{\Gamma }_k`$ as a coarse-grained effective action with modes $`qk`$ already integrated out.
Let us now turn to the effective coarse-grained propagator $`\mathrm{\Delta }_k(q^2)`$, or, for convenience, to its dimensionless inverse given by
$$P^21/k^2\mathrm{\Delta }_k(q^2)=y[1+r(y)],$$
(2)
and $`y=q^2/k^2`$. We find it convenient to write the regulator function $`R_k(q^2)`$ in terms of a dimensionless function $`r(q^2/k^2)`$ as
$$R_k(q^2)=q^2r(q^2/k^2),$$
(3)
normalised as $`r(1)=1`$. As a consequence of the above conditions, we can establish the following properties of $`P^2`$ (see Fig. 1). In contrast to the โfreeโ inverse propagator $`q^2`$, the coarse-grained one is strictly positive for all momenta , $`P^2>0`$. This is due to the regulator term. For large momenta , $`P^2`$ grows linearly with $`y=q^2/k^2`$.
| Class | $`P^2(q^2=0)`$ | $`\mathrm{min}_{q^2}P^2`$ at |
| --- | --- | --- |
| Ia: | finite | $`q^2=0`$ |
| Ib: | finite | $`q^20`$ |
| II: | infinite | $`q^20`$ |
Tab. 1: Classification of different regulator schemes.
For small momenta, $`r(y)`$ diverges at least as $`1/y`$, hence $`P^2`$ either approaches a constant $`>0`$, or it diverges. We refer to the first class of regulators as mass-like, or Class I regulators, and to the second one as Class II regulators (see Tab. 1). Notice that their proper normalisation implies $`P^2(q^2=k^2)=2`$ for any regulator.
Optimisation
The effective coarse-grained inverse propagator $`P^2`$ depends on the particular regularisation scheme chosen. We can formulate a simple optimisation criterion, which is the requirement that the minimum of $`P^2`$ with respect to momenta be maximal with respect to the coarse-graining,
$$C_{\mathrm{opt}}=\underset{\mathrm{RS}}{\mathrm{max}}\left(\underset{y0}{\mathrm{min}}P^2(y)\right)\underset{\mathrm{RS}}{\mathrm{max}}\left(\underset{q^20}{\mathrm{min}}\left[k^2\mathrm{\Delta }_k(q^2)\right]^1\right).$$
(4)
Here, we denote those coarse-grainings as โoptimalโ for which the maximum is attained. It is interesting to note that this criterion is based only on the effective propagator at vanishing field, that is, on the basic ingredient to the flow equation (1). Stated differently, no reference is made to a specific model or theory considered. This is desirable insofar as it renders the optimisation condition universally applicable.
Some few observations concerning $`C_{\mathrm{opt}}`$ can be made at this stage. For an optimal scheme, we have $`C_{\mathrm{opt}}=2`$. This follows from combining $`r(1)=1`$ with (4), which implies $`\mathrm{min}_yP^2(y)P^2(y=1)=2`$ for any scheme. This also implies that the optimal $`P^2`$ reaches its (local) minimum at $`y=1`$ and $`r=1`$ for any optimised scheme. Such regulators are either of Class Ib or Class II. A Class Ia regulator can only be optimal in the event where $`P^2(y)`$ has two degenerate minima, one at $`y=0`$ (by definition), and the other one at $`y=1`$. Hence, an optimal Class Ia regulator is also of Class Ib. The optimal value $`C_{\mathrm{opt}}=2`$ cannot be reached for regulators, which are strictly Class Ia ($`P^2(y=0)`$ is the global minimum).
Threshold functions
The optimisation criterion can be interpreted in more physical terms when generic threshold functions are considered. To give a particular example, consider a $`N`$-component real scalar field theory in $`d`$ dimensions. To leading order in the derivative expansion, we approximate the effective action as $`\mathrm{\Gamma }_k=d^dx[\frac{1}{2}_\mu \varphi ^a^\mu \varphi _a+U_k(\varphi )+O(^4)]`$ and $`a=1,\mathrm{},N`$. It is useful to introduce rescaled (dimensionless) variables as $`u(\rho )=U_k/k^d`$ and $`\rho =\frac{1}{2}\varphi ^a\varphi _ak^{2d}`$. The flow equation for the first derivative of the effective potential $`u^{}`$ is a second-order partial differential equation, to wit
$$_tu^{}=2u^{}+(d2)\rho u^{\prime \prime }2v_d(N1)u^{\prime \prime }\mathrm{}_1^d(\omega _1)2v_d(3u^{\prime \prime }+2\rho u^{\prime \prime \prime })\mathrm{}_1^d(\omega _2)$$
(5)
with $`v_d^1=2^{d+1}\pi ^{d/2}\mathrm{\Gamma }[\frac{d}{2}]`$. Here, the relevant dimensionless amplitudes are $`\omega _1u^{}`$ and $`\omega _2u^{}+2\rho u^{\prime \prime }`$, and the threshold function $`\mathrm{}_1^d(\omega )`$ reads
$$_0^{\mathrm{}}๐yr^{}(y)y^{1+d/2}\frac{1}{(P^2+\omega )^2}.$$
(6)
Flow equations of the form (5) have been intensively studied in the literature .
In the general case, an Ansatz for $`\mathrm{\Gamma }_k`$ need not be restricted to the leading order in the derivative expansion. At higher order, one finds coupled sets of partial differential equations for the coefficient functions of the higher derivative operators, involving more complicated threshold functions. They are obtained from the flow equation (1), or functional derivatives thereof, within specific models and approximations. A generic threshold function has the following structure,
$$_0^{\mathrm{}}๐yK_{n_1,\mathrm{},n_i}(y)[P^2(y)+\omega _1]^{n_1}\mathrm{}[P^2(y)+\omega _i]^{n_i}$$
(7)
and all $`n_i>0`$. Typically, the amplitudes $`w_i`$ describe (field-dependent) mass terms $`\omega m_k^2(\varphi )/k^2`$, which grow large once the coarse-graining scale falls below the mass threshold, $`km_k(\varphi )`$. This leads to the decoupling of the corresponding modes (hence the name โthresholdโ functions). In dimensionful units the physical amplitudes are $`\mathrm{\Omega }k^2\omega m_k^2(\varphi )`$. The explicit form of the kernel $`K_{n_1,\mathrm{},n_i}`$ depends on the specific quantity under investigation. In the example above, the threshold function (6) is of the form (7) for the specific kernel $`K_2(y)=r^{}(y)y^{1+d/2}`$.
As an immediate consequence of the properties of the basic flow equation (1) we can conclude that the kernel of (7) is well-defined, finite and peaked as a function of momenta. Given the properties of $`P^2`$, we deduce that such threshold functions always have a pole on the negative $`\omega `$-axis, located at $`\omega _{\mathrm{pole}}=\mathrm{min}_{y0}P^2(y)<0`$. The location of the pole depends on the particular coarse-graining chosen. Notice that the pole cannot be shifted arbitrarily far away towards the negative axis. Requiring that the pole of any threshold function be as negative as possible corresponds precisely to the condition (4). Hence, the maximally attainable point is given by $`C_{\mathrm{opt}}\omega _{\mathrm{pole}}=2`$. The condition (4) applies to any threshold function of the form (7), and therefore holds for any amplitude $`\omega >\omega _{\mathrm{pole}}`$. In dimensionful units, the pole is located at $`\mathrm{\Omega }_{\mathrm{pole}}=k^2\omega _{\mathrm{pole}}<0`$. For $`k`$ close to some ultraviolet scale $`\mathrm{\Lambda }`$ the pole is far away on the negative axis and outside the physical domain. It is moving inward for decreasing $`k`$, ultimately reaching $`\mathrm{\Omega }_{\mathrm{pole}}=0`$ in the infrared limit, where the potential becomes convex.
The pole structure of threshold functions is closely linked to the flattening of the non-convex parts of effective potentials in the limit $`k0`$, like $`U_k`$ in the example above . The convexity is a consequence of the 1PI generating functional being obtained by a Legendre transformation. For finite $`k>0`$, though, convexity is not required. In the non-convex region, the amplitudes $`\omega _1k^2\varphi ^1U_k/\varphi `$ or $`\omega _2k^2^2U_k/\varphi ^2`$ turn negative. In particular, for $`k0`$, they approach the poles of threshold functions. Hence, choosing a coarse-graining for which this pole is the most negative corresponds to choosing a coarse-graining for which the transition towards a convex effective action is the smoothest. Such a choice is particularly useful for the stability and an improved convergence of numerical implementations of flow equations.
Amplitude expansions
An alternative physical explanation of the optimisation condition is given in terms of an expansion in (small) amplitudes. With amplitude expansion we mean any expansion in the amplitudes $`\omega `$ in units of $`P^2`$. Such expansions always exist, simply because $`P^2`$ is strictly non-zero. Then, all information about the regulator scheme is encoded in the related expansion coefficients, introduced in . Here, they can be written as
$$a_n=_0^{\mathrm{}}๐yK[r]P^n(y).$$
(8)
The explicit form of the kernel $`K`$ depends on the specific quantity studied. For the example given above, the threshold function is Taylor-expanded in powers of the amplitudes $`\mathrm{}_1^d(\omega )=_{m=1}^{\mathrm{}}\frac{2m}{d}a_{2md}(\omega )^{m1}`$ where the coefficients $`a_n`$ are those defined in (8) with the particular kernel
$$K[r]=\frac{d}{2}\frac{r^{}(y)}{[1+r(y)]^{1+d/2}},$$
(9)
normalised such that $`a_0=1`$ .
The main characteristics of amplitude expansions are deduced from the knowledge that the kernel $`K`$ for general expansion coefficients (8) is finite, peaked, and suppressed for sufficiently large momenta $`q^2/k^2`$. Conceptually speaking, it would be most desirable if only a few terms in such a series have to be taken into account, and higher order corrections be small. Hence, we look for coarse-grainings for which the main physical information is already encoded within a few leading order terms. For fixed amplitudes, this amounts to the requirement that the higher order coefficients $`a_n`$ be as small as possible, or, that the radius of convergence for amplitude expansions be maximal. The (dimensionless) radius of convergence $`C`$ is defined as the ratio $`a_n/a_{n+2}`$ of two successive expansion coefficients in the limit $`n\mathrm{}`$,
$$C\underset{n\mathrm{}}{lim}\frac{a_n}{a_{n+2}}.$$
(10)
This limit is computed as follows. First consider the definition (8) and notice that the term $`P^n`$ will strongly suppress $`a_n`$ in the limit $`n\mathrm{}`$ because $`P^2`$ is strictly positive and diverging for large momenta (see Fig. 1). The sole contribution to the integrand will then come from the minimum of $`P^2`$ where the integrand is the least suppressed. Second, taking into account that the kernel $`K`$ is well-behaved, we can conclude that the radius of convergence is given by the โsaddle point approximationโ to (10), namely
$$C=\underset{y0}{\mathrm{min}}P^2(y).$$
(11)
The radius of convergence does not depend on the kernel $`K`$, but still on the parameters of the coarse-graining. The optimisation criterion for amplitude expansions $`C_{\mathrm{opt}}=\mathrm{max}_{\mathrm{RS}}C`$ coincides with the optimisation condition introduced earlier. The finiteness of $`C`$ is due to the pole structure of threshold functions. This implies that the radius of convergence cannot be increased by changing to another expansion point. The optimal expansion point is $`\omega =0`$.
The optimal coarse-grainings have another interesting interpretation. Consider again the definition (8) for the expansion coefficients, and let us split them as $`a_n=a_n^<+a_n^>`$ into its low momentum contributions $`a_n^<=_0^1๐yK[r]P^n`$ and its high momentum contributions $`a_n^>=_1^{\mathrm{}}๐yK[r]P^n`$. We assume that the coarse-graining scheme is parametrised by some parameter $`b`$. Let us define implicitly the parameters $`b_n`$ from the requirement that the hard modes and soft modes give equal contributions to the coefficients $`a_n`$, hence $`a_n^>=a_n^<`$. This is certainly possible for $`n`$ sufficiently large. The optimal choice for the coarse-graining $`b_{\mathrm{opt}}`$, as defined by the optimisation condition, can be shown to correspond precisely to the limit $`b_{\mathrm{opt}}=lim_n\mathrm{}b_n`$.
Field amplitude expansions
It is worth pointing out a crucial difference between expansions based on the amplitudes $`\omega _i`$, as considered here, and expansions based directly on field amplitudes $`\varphi ^a`$ or $`\rho `$ (sometimes also denoted as a local polynomial approximation), which have often been used in the literature .
We consider the $`N`$-component scalar field theory with the flow equation (5) as example. Expansions in field amplitudes $`\rho `$ of the effective potential read $`u^{}=_{n=1}^M\frac{1}{n!}\lambda _n(\rho \lambda _0)^n`$, where $`\lambda _00`$ corresponds to the chosen expansion point.For expansions about $`\rho =0`$ we write $`u^{}=_{n=0}^M\frac{1}{n!}\lambda _n\rho ^n`$ instead. The partial differential equation (5) reduces to a set of coupled ordinary differential equations $`\beta _nd\lambda _n/dt`$ for the $`(M+1)`$ couplings $`\lambda _n`$. These $`\beta `$-functions depend on the coarse-graining only through the threshold function $`\mathrm{}_1^d(\omega )`$ and derivatives thereof, evaluated at the expansion point $`\omega (\lambda _0)`$. The relevant amplitudes are $`\omega _1=u^{}`$ and $`\omega _2=u^{}+2\rho u^{\prime \prime }`$, and the optimal expansion point is $`\omega _1=\omega _2=0`$. When expressed in terms of field amplitudes $`\rho `$, these expansion points are not equivalent.The condition $`\omega _1=\omega _2=0`$ over-constrains the choice for an expansion point $`\lambda _0`$ as it implies either $`u^{}=u^{\prime \prime }=0`$, or $`u^{}(0)=0`$ and $`u^{\prime \prime }(0)=`$ finite. An expansion in $`\rho `$ therefore has a smaller radius of convergence than expansions in the amplitudes $`\omega _i`$.
Interestingly, there are two exceptions in the present example. The first one is the limit $`N\mathrm{}`$, in which case only the amplitude $`\omega _1=u^{}`$ remains. The expansion about the local minimum of the scaling potential $`\rho _1`$ with $`\omega _1(\rho _1)=0`$ is the optimal choice and equivalent to the expansion in the amplitude $`\omega _1`$. A second exceptional case concerns $`N=1`$, where only the amplitude $`\omega _2=u^{}+2\rho u^{\prime \prime }`$ appears, suggesting that an expansion about the inflection point $`\rho _2`$ with $`\omega _2(\rho _2)=0`$ has best convergence properties, as long as $`\omega `$ is Taylor-expandable about $`\rho _2`$.
It was already argued that expansions about the local potential minimum $`\rho _1`$ show better convergence than expansions about vanishing field $`\lambda _0=0`$ . The present discussion clarifies from a more general perspective why an appropriately chosen field expansion point does improve the convergence of the entire series. While the optimal expansion point for amplitude expansions is always $`\omega _i=0`$, the optimal expansion point for field amplitudes $`\rho _{\mathrm{opt}}`$ depends on the particular theory studied (in the example discussed here on the number of real scalar fields $`N`$). Clearly, all these expansions, either in $`\rho `$ or in the amplitudes $`\omega _i`$, are optimised through an appropriate choice of the coarse-graining, as defined in (4).
Applications
We now turn our attention to a number of specific coarse-graining functions. Our list of regulators is by no means exhaustive, though it covers the main regulators used and discussed in the literature. It includes exponential, algebraic and mass-like ones, and the sharp cut-off for comparison (see Tab. 2; the parameter $`c`$ ensures the consistent normalisation of the scale parameter $`k`$). They all differ in the way how high momentum modes are cut off, with the parameter $`b`$ measuring the relative โsmoothnessโ โ the smaller $`b`$ the โsmootherโ the coarse-graining.
| Regulator | Parameter range | | Class | $`b_{\mathrm{opt}}`$ | $`C_{\mathrm{opt}}`$ |
| --- | --- | --- | --- | --- | --- |
| $`r_{\mathrm{exp}}`$ = $`\left[\mathrm{exp}cy^b1\right]^1`$ | $`c=\mathrm{ln}2`$ | $`1b\mathrm{}`$ | II | $`1.44`$ | 2 |
| $`r_{\mathrm{mexp}}`$ = $`b\left[\mathrm{exp}cy1\right]^1`$ | $`c=\mathrm{ln}(1+b)`$ | $`0b\mathrm{}`$ | I | $`3.92`$ | 2 |
| $`r_{\mathrm{mod}}`$ = $`\left[\mathrm{exp}\frac{c}{b}(y+(b1)y^b)1\right]^1`$ | $`c=\mathrm{ln}2`$ | $`1b\mathrm{}`$ | I | $`1.92`$ | 2 |
| $`r_{\mathrm{mix}}`$ = $`\mathrm{exp}\left[\frac{b}{2a}(y^ay^a)\right]`$ | $`a0`$ | $`1b\mathrm{}`$ | II | 2 | 2 |
| $`r_{\mathrm{step}}`$ = $`\frac{2b2}{b}y^{b2}\left[\mathrm{exp}cy^{b1}1\right]^1`$ | $`c=\mathrm{ln}\frac{3b2}{b}`$ | $`1b\mathrm{}`$ | I | $`4.40`$ | $`1.65`$ |
| $`r_{\mathrm{power}}`$ = $`y^b`$ | | $`1b\mathrm{}`$ | II | 2 | 2 |
| $`r_{\mathrm{sharp}}`$ = $`\frac{1}{\theta (y1)}1`$ | | | II | | 1 |
Tab. 2: The definition of various classes of regulator functions and their optimal parameter choices, including the sharp cutoff for comparison. Class II denotes Class II regulators except for $`b=1`$, where they correspond to Class I.
More specifically, the regulators $`r_{\mathrm{exp}},r_{\mathrm{mod}},r_{\mathrm{step}}`$, $`r_{\mathrm{mexp}}`$ for $`b0`$ and $`r_{\mathrm{mix}}`$ for $`a1`$ are exponential ones, that is, exponentially suppressed for large momenta. In contrast, $`r_{\mathrm{power}}`$, $`r_{\mathrm{mexp}}`$ for $`b=0`$, and $`r_{\mathrm{mix}}`$ for $`a=0`$ are algebraic regulators. In the classifaction of Tab. 1, the mass-like exponential regulator $`r_{\mathrm{mexp}}`$, the modified exponential regulator $`r_{\mathrm{mod}}`$ and $`r_{\mathrm{step}}`$ are Class I (mass-like) regulators, as are the exponential one $`r_{\mathrm{exp}}`$ and the algebraic one $`r_{\mathrm{power}}`$ for $`b=1`$. On the other hand, the mixed exponential regulator $`r_{\mathrm{mix}}`$, the sharp cut-off $`r_{\mathrm{sharp}}`$, and $`r_{\mathrm{exp}}`$ and $`r_{\mathrm{power}}`$ for $`b>1`$, are Class II regulators.
Some cross-dependences amongst these functions are worth noticing. For $`b\mathrm{}`$, $`r_{\mathrm{exp}},r_{\mathrm{mexp}},r_{\mathrm{mod}}`$ and $`r_{\mathrm{power}}`$ approach the sharp cut-off limit $`r_{\mathrm{sharp}}`$, while $`r_{\mathrm{step}}`$ approaches $`y^1\theta (1y)`$, which corresponds to a mass term which is cut-off for momenta above $`k`$, $`R_{\mathrm{step},k}k^2\theta (k^2q^2)`$. The two-parameter family of mixed exponential regulators $`r_{\mathrm{mix}}`$, which obey $`r_{\mathrm{mix}}(1/y)=1/r_{\mathrm{mix}}(y)`$, contain the algebraic regulator $`r_{\mathrm{power}}`$ as the limiting case $`a0`$. $`r_{\mathrm{exp}}`$ and $`r_{\mathrm{mod}}`$ for $`b=1`$ coincide with $`r_{\mathrm{step}}`$ for $`b=2`$. Also, $`r_{\mathrm{mexp}}`$ for $`b=0`$, and $`r_{\mathrm{step}}`$ for $`b=1`$ correspond to $`r_{\mathrm{power}}`$ for $`b=1`$.
We now distinguish the different coarse-grainings by their associated radii of convergence. In Fig. 2, we have displayed $`C`$, as defined in (10), for different coarse-graining schemes and as functions of the coarse-graining parameter $`b`$. In the sharp cutoff limit $`b\mathrm{}`$ the value $`C_{\mathrm{sharp}}=1`$ is attained. For the classes of regulators considered here, we notice that the maximal radius is reached for specific smooth coarse-grainings. The radii vary between $`1C2`$ (though it is possible to find regulators such that $`C`$ is arbitrarily small). Optimal parameters are found for both Class I and Class II regulators (see Fig. 3). The most efficient way for determining optimal parameters for all Class Ib and Class II regulators is the following. We have established that $`P^2(y)`$ takes its maximum at $`y=1`$ and $`r=1`$ for an optimal coarse-graining. The vanishing of $`dP^2/dy=1+r+yr^{}`$ at the extremum implies $`r_{\mathrm{opt}}^{}(1)=2`$. This allows to fix the optimal coarse-graining parameters.
For $`r_{\mathrm{exp}}`$ one finds $`b_{\mathrm{exp}}=1/\mathrm{ln}2`$, for $`r_{\mathrm{power}}`$ the radius of convergence reads $`C=b(1b)^{1/b1}`$ which attains its maximum at $`b_{\mathrm{power}}=2`$. The optimal parameter for $`r_{\mathrm{mix}}`$ is $`b_{\mathrm{mix}}=2`$ for all $`a`$. The value for $`b_{\mathrm{mexp}}`$ is obtained numerically by solving $`2b=(b+1)\mathrm{ln}(b+1)`$, and the one for $`b_{\mathrm{mod}}`$ reads $`b_{\mathrm{mod}}=[1+\frac{1}{\mathrm{ln}2}][\frac{1}{2}+\sqrt{\frac{1}{4}1/(1+\frac{1}{\mathrm{ln}2})^2}]`$ (see Tab. 2). The regulator $`r_{\mathrm{step}}`$ does not attain the maximal radius $`C=2`$. The reason behind it is that $`P_{\mathrm{step}}^2(y)`$ has two local minima for $`bb_{\mathrm{opt}}`$, one at $`y1`$, and another one at $`y=0`$. The minima are degenerate for $`b=b_{\mathrm{opt}}`$, and their competition only allows for a sub-optimal value $`C_{\mathrm{step}}=1.65`$.
The fact that $`C_{\mathrm{sharp}}=1`$ is half as big as the possible optimal value explains why amplitude expansions of sharp cut-off flows have a notoriously bad convergence behaviour . This property is not a particularity of the sharp cut-off; the smooth cut-offs $`r_{\mathrm{power}}`$, $`r_{\mathrm{exp}}`$ and $`r_{\mathrm{mexp}}`$ (all for $`b=1`$) also have a radius of convergence $`C=1`$. A significant improvement is achieved only for the optimal coarse-grainings, which all belong to smooth cut-offs.
The present discussion can be extended in several directions. Fermionic fields can be included using a regulator term like $`d^dq\overline{\psi }(q)R_{\mathrm{F},k}(q)\psi (q)`$, with $`R_{\mathrm{F},k}^2(q)=q^2r_\mathrm{F}^2(q^2/k^2)`$. Effectively, it appears in threshold functions in the combination $`P_\mathrm{F}^2=y[1+r_\mathrm{F}(y)]^2`$, which is the fermionic analogue of (2). Hence, the optimisation condition for fermionic degrees of freedom is the same as (4), replacing $`P^2`$ by $`P_\mathrm{F}^2`$, yielding $`C_{\mathrm{F},\mathrm{opt}}=4`$. For gauge fields, the flow equation is amended by a (modified) Ward or a (modified) BRST Identity, which ensures that physical Green functions obey the usual Ward or BRST Identities . The criterion (2) is compatible with such an additional constraint. Also, wave function renormalisations can be taken into account in the usual manner, setting $`R_k=Z_kq^2r(q^2/k^2)`$. This modification will not change the optimisation as discussed in the present note, and allows for a systematic inclusion of higher derivative operators. The same holds for field theories at finite temperature within the imaginary or the real-time formalism .
In conclusion, we have presented a simple and generic optimisation criterion for exact renormalisation group equations. It is interesting to mention its close link to a minimisation of the spurious scheme dependence of flow equations, a more detailed discussion of which shall be presented elsewhere .
The author thanks Ulrich Ellwanger, Filipe Freire, Jan M. Pawlowski, Janos Polonyi and Christof Wetterich for useful discussions.
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# Generalised Inflation with a Gravitational Wave Background
## 1 Introduction
The observational reconstruction of the cosmological density perturbation (CDP) spectrum is a key problem of modern cosmology. It provided a dramatic challenge after detecting the primordial cosmic microwave background (CMB) anisotropy by DMR COBE (Smoot et al. 1992, Bennett et al. 1992) as the signal found at $`10^0`$, $`\mathrm{\Delta }T/T=1.06\times 10^5`$ (Bennett et al. 1996), appeared to be few times higher than the expected value of $`\mathrm{\Delta }T/T`$ in the most simple and best developed cosmological model with standard cold matter (sCDM)<sup>1</sup><sup>1</sup>1The matter density $`\mathrm{\Omega }_m=\mathrm{\Omega }_b+\mathrm{\Omega }_{cm}=1`$, there is no cosmological term , $`h=0.5`$, the slope of the CDP spectrum $`n_S=1`$, i.e. of Harrison-Zeldovich type, no contribution from cosmological gravitational waves.. Indeed, the CDP spectrum of sCDM normalized by the biasing parameter $`b^1\sigma _8=0.6`$ (White, Efstathiou, Frenk 1993, Eke, Coles, Frenk 1996, Viana, Liddle 1996) can reproduce only 30% of the COBE measured CMB anisotropy. If $`\sigma _8`$ is less than $`0.6`$, this contradiction gets even stronger.
During recent years there were many proposals to improve sCDM<sup>2</sup><sup>2</sup>2In the simplest terms, to remove the discrepancy between the CDP spectrum amplitudes at $`8h^1\mathrm{Mpc}`$ as determined by galaxy clusters, and at large scales, $`1000h^1\mathrm{Mpc}`$, by the $`\mathrm{\Delta }T/T`$ inhomogeneity. by adding hot dark matter, a $`\mathrm{\Lambda }`$-term, or considering non-flat primordial CDP spectra. Below, we present another, presumably more natural way to solve the sCDM problem based on taking into account a possible contribution of cosmic gravitational waves (CGWs) into the large-scale CMB anisotropy. Further, we will also try to preserve the original near-scale-invariant CDP spectrum in the CDM universe. Thus, the problem is reduced to the construction of a simple inflation producing a near Harrison-Zelโdovich (HZ) spectrum of CDPs ($`n_S1`$) and a large relative contribution of CGWs into the $`\mathrm{\Delta }T/T`$ at COBE scale (i.e. the ratio of the tensor to scalar mode contributions T/S $``$ 1).
A basic physical reason for the production of tensor and scalar perturbations in the expanding Universe is the parametric amplification effect (Grishchuk 1974, Lukash 1980): the spontaneous creation of quantum physical fields in a non-stationary gravitational background. From the theoretical point of view a cosmology with negligible contribution of CGWs may be considered as a phenomenological model only, because any inflationary model predicts a non-zero CGW amplitude. Generally, the tensor mode is not discriminated in inflation (e.g. Starobinsky 1979, Rubakov, Sazhin, Veryaskin 1982). Different models with nearly scale-invariant spectra of CDPs predict different abundance of CGWs (Lucchin, Matarrese 1985, Linde 1994, Garsia-Bellido, Linde, Wands 1996, Lukash, Mikheeva 1996). However, in chaotic and power-law inflations it is usually small, T/S $`{}_{}{}^{}{}_{}{}^{<}`$ 1, which corresponds to the spectrum slope $`n0.8`$. An always โredโ CDP spectrum ($`n_S<1`$) originating in the power-law inflation model helps to satisfy the โdouble-normalizationโ (i.e. to reconcile the $`\sigma _8`$ vs. $`\mathrm{\Delta }T/T`$ problem).
So, the first candidate of a successful model fitting both COBE and galaxy cluster normalizations could be a power-law spectrum of CDPs with a slope $`n_S0.85`$ predicted by power-law inflation (e.g. Lucchin, Matarrese 1985, there an exponential potential of the inflation field is used). One has a simple estimate for the fraction of CGWs,
$$\frac{\mathrm{T}}{\mathrm{S}}6n_T,$$
(1)
where $`n_T`$ is the slope of the CGW spectrum related trivially to $`n_S`$ in the case of power-law inflation: $`n_T=n_S1`$.
Notice that large T/S and hence the desired double-normalization can be reached in this model only at the expence of a rejection from the HZ spectrum: the red power-law CDP spectrum helps on large scale, however it is getting undesirable on Mpc scale producing a too late galaxy formation epoch (Gardner et al. 1997, and references cited therein).
Therefore, it is interesting to consider other models with high abundance of CGWs but at the same time with $`n_S1`$ unchanged (with a preference of HZ or a slightly blue CDP spectrum on short scales to provide early formation of high-redshift quasars and early galaxy formation).
A simple model of such kind is $`\mathrm{\Lambda }`$-inflation, an inflationary model with an effective metastable $`\mathrm{\Lambda }`$-term (Lukash, Mikheeva 1996, 1997, 2000, Mikheeva 1997). This model produces both curvature and primordial gravitational wave perturbations which have a non-power-law spectrum, with a shallow minimum in the CDP spectrum, located at a scale $`k_{cr}`$ (where the $`\mathrm{\Lambda }`$-term and the scalar field have equal energy densities). Around this scale, the CDP spectrum is exactly of the scale invariant HZ form, and the ratio T/S is close to its maximum; it is of the order unity depending on the model parameters.
The estimate (1) remains true for $`\mathrm{\Lambda }`$-inflation and, probably, keeps its universality for any type of inflationary dynamics (however, the relationship between $`n_T`$ and $`n_S`$ can vary). The cost to be paid for the possibility of having a high ratio T/S on scales where $`n_S1`$ is the non-power-law CDP global spectrum: it is โredโ at $`k<k_{cr}`$ and โblueโ at $`k>k_{cr}`$. This smooth transition in the spectrum slope from red to blue makes T/S obviously dependent on how $`k_{COBE}`$ is related to $`k_{cr}`$. One may benefit from the HZ local slope or a blue spectrum enhancement on Mpc scales (to initiate the early structure formation) by simply adjusting $`k_{cr}`$ to a galaxy cluster scale.
## 2 $`\mathrm{\Lambda }`$-inflation with self-interaction
We summarize briefly the main properties of $`\mathrm{\Lambda }`$-inflation and determine the basic model for our investigation. Let us consider a general potential of $`\mathrm{\Lambda }`$-inflation:
$$V(\phi )=V_0+\underset{\kappa =2}{\overset{\kappa _{max}}{}}\frac{\lambda _\kappa }{\kappa }\phi ^\kappa ,$$
(2)
where $`\phi `$ is the inflaton scalar field, $`V_0>0`$ and $`\lambda _\kappa `$ are constants, and $`\kappa =2,3,4,..`$.
In the case of a massive inflaton ($`\kappa =\kappa _{max}=2`$, $`\lambda _2m^2>0`$, this model is called $`\mathrm{\Lambda }m`$-inflation) T/S can be larger than unity only when the CDP spectrum slope in the โblueโ asymptote is very steep, $`n_S^{blue}>1.8`$. To avoid such a strong spectral bend on short scales ($`k>k_{cr}`$), we choose here another simple version of $`\mathrm{\Lambda }`$-inflation โ the case with self-interaction: $`\kappa =\kappa _{max}=4`$, $`\lambda _4\lambda >0`$; this model is called $`\mathrm{\Lambda }\lambda `$-inflation.
The scalar field $`\phi `$ drives an inflationary evolution if $`\gamma \dot{H}/H^2<1`$, where<sup>3</sup><sup>3</sup>3We assume Planckian units $`8\pi G=c=\mathrm{}=1`$, a dot denotes the time derivative ($`d/ad\eta `$), and $`a`$ and $`H\dot{a}/a`$ are scale and Hubble factors, respectively. $`H=\sqrt{V/(3\gamma )}\sqrt{V/3}`$. This condition holds true for all values of $`\phi `$ (at worst, for $`\phi \phi _{cr}`$, where $`V_0=\lambda \phi _{cr}^4/4`$ and $`\gamma `$ reaches its maximum) if
$$c\frac{1}{4}\phi _{cr}^2=\frac{1}{2}\sqrt{\frac{V_0}{\lambda }}>1,$$
(3)
which we imply hereafter.
The last inequality simultaneously ensures the validity of the slow-roll approximation ($`|\ddot{\phi }/\dot{\phi }|3H`$) as
$$\gamma \frac{2}{c}\left(\frac{y^3}{1+y^4}\right)^2,\frac{\ddot{\phi }}{3H\dot{\phi }}\frac{y^2(3+y^4)}{3c(1+y^4)^2},$$
(4)
where $`y\phi /\phi _{cr}=\phi /\left(2\sqrt{c}\right)`$. Then we get the value of the inflaton field at horizon-crossing time ($`k=aH1/\eta `$),
$$cy^2=\sqrt{c^2+x^2}x,$$
and the gravitational perturbation spectra $`q_k`$ and $`h_k`$ generated in S and T modes, respectively (see the Appendix):
$$q_k=\frac{H}{2\pi \sqrt{2\gamma }}=\frac{\sqrt{2\lambda /3}}{\pi }\left(c^2+x^2\right)^{3/4},$$
(5)
$$h_k=\frac{H}{\pi \sqrt{2}}=\frac{2c\sqrt{\lambda /3}}{\pi }\left(1+\frac{x}{\sqrt{c^2+x^2}}\right)^{1/2},$$
(6)
where
$`x`$ $`=`$ $`\mathrm{ln}\left[{\displaystyle \frac{k}{k_{cr}y}}\left({\displaystyle \frac{2}{1+y^4}}\right)^{1/6}\right]`$
$`=`$ $`\mathrm{ln}\left[{\displaystyle \frac{k}{k_{cr}}}\left(1+\left({\displaystyle \frac{x}{c}}\right)^2\right)^{1/4}\left(1+{\displaystyle \frac{x}{\sqrt{c^2+x^2}}}\right)^{2/3}\right]`$
$``$ $`\mathrm{ln}(k/k_{cr}).`$
Fig.1 shows the power spectra for $`c=5,9,11`$.
The dimensionless power spectrum of CDPs is directly related to the fundamental $`q_k`$ spectrum
$$\delta ^2=\underset{0}{\overset{\mathrm{}}{}}\mathrm{\Delta }_k^2\frac{dk}{k},\mathrm{\Delta }_k=3.6\times 10^6\left(\frac{k}{h}\right)^2q_kT(k),$$
(7)
where the wave number $`k`$ is measured in $`h/\mathrm{Mpc}`$, and $`T(k)`$ is the transfer function.
We can find also the local slopes of the fundamental power spectra (5), (6):
$$n_S12\frac{d\mathrm{ln}q_k}{d\mathrm{ln}k}=\frac{3x}{c^2+x^2}[\frac{3}{2c},\frac{3}{2c}],$$
(8)
$`n_T2{\displaystyle \frac{d\mathrm{ln}h_k}{d\mathrm{ln}k}}=2\gamma ={\displaystyle \frac{1}{\sqrt{c^2+x^2}}}\left(1{\displaystyle \frac{x}{\sqrt{c^2+x^2}}}\right)`$ (9)
$`[{\displaystyle \frac{3\sqrt{3}}{4c}},0].`$
Obviously, the maximal deviations from the HZ spectrum take place in S-mode at $`x=\pm c`$, and in T-mode at $`x=c/\sqrt{3}`$. At the latter point the ratio of spectra reaches its maximum,
$$\left(\frac{h_k}{q_k}\right)^2=4\gamma =2|n_T|\frac{3\sqrt{3}}{2c}.$$
(10)
## 3 CDM cosmology from $`\mathrm{\Lambda }`$-inflation
Let us consider the CDP spectrum $`\mathrm{\Delta }_k`$ (7) with CDM transfer function, normalized both at the large-scale $`\mathrm{\Delta }T/T|_{10^0}`$ (including the contribution from $`h_k`$ (6)) and the galaxy cluster abundance at $`z=0`$ ($`\sigma _8`$), to find the family of the most realistic S-spectra $`q_k`$ produced in $`\mathrm{\Lambda }\lambda `$-inflation.
In total, we have three parameters entering the $`q_k`$ spectrum: $`\lambda `$ (the overall amplitude), $`c`$ (the measure of T/S) and $`k_{cr}`$ (the scale where the CDP spectrum is locally HZ, $`n_S=1`$<sup>4</sup><sup>4</sup>4$`n_S=1`$ is also in the asymptotics $`|x|c`$.). Constraining them by two observational tests, we are actually left with only one free parameter (say, $`k_{cr}`$) which may be restricted elsewhere by other observations (e.g. cluster power spectra, acoustic peaks, bulk velocities, etc.).
To demonstrate explicitly how the three parameters are mutually related, we first employ simple analytical estimates for the $`\sigma _8`$ and $`\mathrm{\Delta }T/T`$ tests to derive the key equation relating $`c`$ and $`k_{cr}`$, and then solve this equation numerically to obtain the range of interesting physical parameters.
Instead of taking the $`\sigma `$-integral numerically ($`R=8h^1\mathrm{Mpc}`$),
$$\sigma _R^2=\underset{0}{\overset{\mathrm{}}{}}\mathrm{\Delta }_k^2W^2(kR)\frac{dk}{k},W(z)=\frac{3}{z^3}(\mathrm{sin}zz\mathrm{cos}z),$$
(11)
we may estimate the wavelength $`k_1`$ at which $`\mathrm{\Delta }_{k_1}=1.6\sigma _8`$ (to be equal unity for $`\sigma _8=0.6`$)<sup>5</sup><sup>5</sup>5The integral is roughly estimated as $`\sigma _8^2\mathrm{\Delta }_{k_1}^2/\alpha `$ assuming that the integrand grows sharp with $`k`$ and $`\mathrm{\Delta }_k^2k^\alpha `$ near $`kk_1`$. For sCDM, $`\alpha 2.5`$ at $`k_10.3h/\mathrm{Mpc}`$.. This will fix the spectrum amplitude on the cluster scale (see eq.(7)):
$$q_{k_1}4.5\times 10^7\frac{h^2\sigma _8}{k_1^2T(k_1)}.$$
(12)
On the other hand, the spectrum amplitude on large scale ($`k_2=k_{COBE}10^3h/\mathrm{Mpc}`$) can be taken from $`\mathrm{\Delta }T/T`$ due to the Sachs-Wolfe-effect (Sachs, Wolfe 1967):
$$\left(\frac{\mathrm{\Delta }T}{T}\right)^2_{10^0}=\mathrm{S}+\mathrm{T}1.1\times 10^{10},\mathrm{S}=0.04\mathrm{q}^2_{10^0}.$$
(13)
The relationship between the power spectrum at COBE scale and the variance of the $`q`$ potential averaged in $`10^0`$-angular-scale at the last scattering surface, involves an effective interval of the spectral wavelengths contributing to the latter:
$$q^2_{10^0}^{1/2}=fq_{k_2},f^2\mathrm{ln}\left(\frac{k_2}{k_{hor}}\right).$$
(14)
To estimate $`\mathrm{T}/\mathrm{S}`$, we will use the approximation formula (1) for $`x_2=x_{COBE}`$ (cf. eqs.(9), (10)):
$`{\displaystyle \frac{\mathrm{T}}{\mathrm{S}}}6n_T=12\gamma `$ $`=`$ $`3\left({\displaystyle \frac{h_{k_2}}{q_{k_2}}}\right)^2`$
$`=`$ $`{\displaystyle \frac{6}{\sqrt{c^2+x_2^2}}}\left(1{\displaystyle \frac{x_2}{\sqrt{c^2+x_2^2}}}\right).`$
Evidently, both normalizations determine essentially the corresponding $`q_k`$ amplitudes at the locations of cluster ($`k_1`$) and COBE ($`k_2`$) scales. Accordingly, the parameters $`k_1`$ and $`f`$ can slightly vary while changing the local spectrum slopes at the respective wavelengths. However, when the deviation of the $`q_k`$ slopes from HZ is small (e.g. $`c>5`$, see eq.(8)) we can just identify the parameters $`k_1`$ and $`f`$ for CDM models with their sCDM values:
$$k_10.3h/\mathrm{Mpc},\mathrm{f}1.26.$$
Finally, comparing the two normalization conditions with the theoretical spectrum $`q_k`$ in eq.(5) we get the key equation for the relationship between $`c`$ and $`k_{cr}`$,
$$\left(\frac{q_{k_1}}{q_{k_2}}\right)^2D\left(1+\frac{\mathrm{T}}{\mathrm{S}}\right),$$
(16)
which looks like an algebraic equation for finding $`x_2`$ by given $`c`$:
$$\left(1+\frac{d(d+2x_2)}{c^2+x_2^2}\right)^{3/2}D\left(1+\frac{6}{\sqrt{c^2+x_2^2}}\frac{6x_2}{c^2+x_2^2}\right).$$
(17)
Here $`D\sigma _8^2`$, and
$$d=x_1x_2\mathrm{ln}\left[300\left(\frac{c^2+x_2^2}{c^2+x_1^2}\right)^{1/12}\left(\frac{\sqrt{c^2+x_2^2}x_2}{\sqrt{c^2+x_1^2}x_1}\right)^{2/3}\right].$$
The $`\lambda `$-parameter is then obtained as
$$\sqrt{\lambda }\frac{10^4\sigma _8}{(c^2+x_1^2)^{3/4}}.$$
Eq.(16) has a clear physical meaning: the ratio of the S-spectral power at cluster and COBE scales is proportional to $`\sigma _8^2`$ and inversely proportional to the fraction of the scalar mode contributing to the large-scale temperature anisotropy variance, S/(S+T). This simple argument makes eq.(16) independent of a particular way how it was obtained, just the proportionality coefficient should be taken properly.
Eq.(16) provides quite a general constraint on the fundamental inflation spectra (both S and T) in a wide set of dark matter models using only two basic measurement (the cluster abundance and large scale $`\mathrm{\Delta }T/T`$). Actually, the DM information (a transfer function) is contained in the $`D`$-coefficient which can be calculated using the same equation (16) for a simple inflationary spectrum (e.q. power-law) preserving the given DM model. E.g. for CDM with $`h=0.5`$ we have:
$$D\frac{0.6\sigma _8^2}{13.1\mathrm{\Omega }_b},\mathrm{\Omega }_b<0.2.$$
(18)
The solution of eq.(17) is shown in the plane $`x_2c`$ for D=0.2, 0.3, 0.4 (see Fig. 2). For the whole range $`0.1<D<0.5`$, it can be analytically approximated with a precision better than 10% as follows:
$$\mathrm{ln}^2\left(\frac{k_0}{k_{cr}}\right)E\left(c_0c\right)\left(c+c_1\right),\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}2}<cc_0.$$
(19)
Notice there exists not any solution of eq.(17) for high enough $`c`$ ($`c>c_0`$). We have found the following best fit coefficients $`E,k_0[h/\mathrm{Mpc}]`$ and $`c_{0,1}`$:
$$E1,\mathrm{ln}k_049D^2+1.3,$$
$$c_061D^2+6.2,c_144D^2+4.0.$$
The tensor-mode-contribution is approximated similarly:
$$\frac{\mathrm{T}}{\mathrm{S}}\frac{2.534.3D}{(\mathrm{ln}k_{cr}+4.65)^{2/3}}+\frac{1}{3}.$$
(20)
Recall, $`k_{cr}`$ is measured in the units $`h/\mathrm{Mpc}`$.
## 4 Discussion
We have presented a new inflationary model predicting a near scale-invariant spectrum of density perturbations and large amount of CGWs. Our model is consistent with COBE $`\mathrm{\Delta }T/T`$ and cluster abundance data. The perturbation spectra depend on one free scale-parameter, $`k_{cr}`$, which can be found in further analysis by fitting other observational data. At the location of $`k_{cr}`$, the CDP spectrum transients smoothly from the red ($`k<k_{cr}`$) to the blue ($`k>k_{cr}`$) parts (see eq.(5)).
By adjusting $`k_{cr}`$ with the galaxy cluster scale, we easily gain the boosts in power (in comparison with sCDM) on both scales, large (voids and superclusters) and small (quasars and $`Ly_\alpha `$ clouds). Say, for $`k_{cr}=k_1,D=0.4`$ and $`\mathrm{\Omega }_b=0.1`$, we get $`c11`$, $`\sqrt{\lambda }1.610^6`$ and
$$\mathrm{T}/\mathrm{S}0.7,\mathrm{n}_\mathrm{S}0.9$$
(21)
at large scale ($`1000h^1\mathrm{Mpc}`$). The boost on Mpc-scale is about $`8\%`$,
$$\left(\frac{q_{10k_1}}{q_{k_1}}\right)^21.08,$$
which is really a lot when compared with the red spectrum (21) extrapolated from large scales.
We conclude that the wing-like S-power-spectra similar to those in eq.(5), can provide a simple solution to the cosmological problem in a matter-dominated universe (cp. also a similar spectrum of CDPs in Semig, Mรผller 1996). Indeed, the spectrum power at the current dynamical scale is strongly suppressed by the present galaxy cluster abundance ($`\sigma _80.6`$), whereas other โclassicalโ observations persistently require a large CDP power both on scales $`100h^1\mathrm{Mpc}`$ (the existence of large scale structures) and $`1h^1\mathrm{Mpc}`$ (QSOs, $`Ly_\alpha `$ forest, early galaxy formation).
More of that, today we seriously discuss a nearly flat shape of the dimensionless (linear) CDP-spectrum within the scale range encompassing clusters and superclusters,
$$\mathrm{\Delta }_k^2k^{(0.9\pm 0.2)},k(0.05,0.2)h\mathrm{Mpc}^1,$$
(22)
(with a break towards the HZ slope on higher scales) which stays in obvious disagreement with the sCDM predictions. The first arguments supporting eq.(22) came from the analysis of large-scale galaxy distribution (Guzzo et al. 1991) and the discovery of large quasar groups (Komberg, Lukash 1994, Komberg, Kravtsov, Lukash 1996). Recent measurements of the galaxy cluster power spectrum (Tadros, Efstathiou, Dalton 1998, Retzlaff et al. 1998) brought a higher statistical support for eq.(22) (see Fig.3).
A possible explanation of eq.(22) from a theoretical point of view can be a fundamental red power spectrum established on large scales; then the transition to the spectrum (22) at $`100h^1\mathrm{Mpc}`$ would be much easier understood with help of a traditional modification of the transfer function $`T(k)`$ (e.g. for mixed hot+cold dark matter, see Mikheeva et al. 2000). A boost of power of the high-multipole CMB anisotropy (at $`l220`$) looks self-consistent with the above argument on a slightly red S-spectrum at large scales. The rediness may be not too high, remaining in the range (0.9, 1). Notice any introduction of a blue power spectrum at large scale appears extremely unfavorable in this connection as it suppresses severely all the LSS effects at $`100h^1\mathrm{Mpc}`$ (provided the spectrum meets the $`\sigma _8`$ constraints).
The main problem for the LSS in matter-dominated models remains a low number of $`\sigma _8`$: if $`\sigma _80.6`$, then the first acoustic peak in $`\mathrm{\Delta }T`$ cannot be as high as $`70\mu K`$ in any model with $`\mathrm{\Omega }_m=1`$ (see Fig.4, cp. a similar conclusion for MDM models in Mikheeva et al. 2000).
While the red power spectrum could help in solving the LSS problem, it is undesirable when continued to Mpc scale: here we need an enhanced power of CDPs to trigger the observed early structure formation. This is the way to see how a wing-like power spectrum can be of help.
Summarizing, our arguments are very simple:
(i) The wing-like S-power-spectra is a common feature of $`\mathrm{\Lambda }`$-inflation;
(ii) An early galaxy formation period together with the developed structure at $`(10100)h^1\mathrm{Mpc}`$ can be reconciled (for models with standard cold or mixed dark matter) in the considered wing-like primordial spectrum of CDPs.
### Acknowledgments
The work was partly supported by German Scientific Foundation (DFG-436 RUS 113/357/0) and INTAS grant (97-1192).
## Appendix
Here, we recall the basic properties of the parametric amplification effect and derive the cosmological perturbation spectra originating in $`\mathrm{\Lambda }\lambda `$-inflation (Lukash, Mikheeva 2000).
The linear perturbations in the Friedmann geometry are irreducibly presented in terms of the uncoupled scalar and tensor parts:
$$ds^2=(1+h_{00})dt^2+2ah_{0\alpha }dtdx^\alpha a^2(\delta _{\alpha \beta }+h_{\alpha \beta })dx^\alpha dx^\beta ,$$
$`(A.1)`$
$$\frac{1}{2}h_{\alpha \beta }=A\delta _{\alpha \beta }+B_{,\alpha \beta }+G_{\alpha \beta },h_{0\alpha }=C_{,\alpha },$$
where $`G_{\alpha \beta }`$ is a trace-free divergence-less tensor field ($`G_\alpha ^\alpha =G_{\alpha ,\beta }^\beta =0`$), and potentials $`h_{00}`$, $`A`$, $`B`$, $`C`$ are coupled to the perturbation of the scalar field $`\phi `$ (which is also the source of the Friedmann geometry). The vector mode is not included here for the absence of proper source.
The gauge-invariant canonical 4-scalar determining the physical scalar perturbations, $`q=q(t,\stackrel{}{x})`$, is uniquely fixed by the appearance of the S-part of the perturbative Lagrangian similar to a massless field (Lukash 1980, 1996)<sup>6</sup><sup>6</sup>6For simplicity, we assume the total field Lagrangian density in the form $`\frac{1}{2}\phi _{,i}\phi ^{,i}V(\phi )`$ and the Hilbert minimal action for gravity.:
$$\delta LL(q,G_{\alpha \beta })=\gamma q_{,i}q^{,i}+\frac{1}{2}G_{\alpha \beta ,\gamma }G^{\alpha \beta ,\gamma }.$$
$`(A.2)`$
The relation of $`q`$ to the original potentials have the following form:
$$\delta \phi =\alpha \left(q+A\right),a^2\dot{B}+C=\frac{\mathrm{\Phi }+A}{H},$$
$$\frac{1}{2}h_{00}=\gamma q+\left(\frac{A}{H}\right)^.,\mathrm{\Phi }=\frac{H}{a}a\gamma q๐t,$$
$`(A.3)`$
$$\frac{\delta \rho }{\rho +p}=\frac{\dot{q}}{H}3(q+A),\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}4}\pi G\delta \rho _c\gamma H\dot{q}=\frac{\mathrm{}\mathrm{\Phi }}{a^2},$$
where $`a`$, $`H`$, $`\gamma `$, $`\rho =\dot{\phi ^2}/2+V`$, $`\alpha =\dot{\phi }/H=\pm \sqrt{2\gamma }`$ are solutions of the Friedmann equations of the background cosmological model, i.e. they are pure time dependent, and $`\mathrm{}^2/\stackrel{}{x}^2`$ is the spatial Laplacian. Any two potentials taken from the triple $`A`$, $`B`$, $`C`$ are arbitrary functions of all coordinates specifying the gauge choice.
The equations of motions of the scalar and tensor fields propagating in the Friedmann Universe are two harmonic oscillators,
$$\ddot{q}+\left(3H+\frac{\dot{\gamma }}{\gamma }\right)\dot{q}\frac{\mathrm{}}{a^2}q=0,$$
$`(A.4)`$
$$\ddot{G_{\alpha \beta }}+3H\dot{G_{\alpha \beta }}\frac{\mathrm{}}{a^2}G_{\alpha \beta }=0,$$
$`(A.5)`$
which reduce the problem of the generation of cosmological perturbations to the well established parametric amplification effect (particle creation in intensive gravitational fields).
The quantum-generated perturbations inflating outside the horizon become frozen in time and can be treated as classical random Gaussian fields with the variances
$$q^2=_0^k_{}q_k^2\frac{dk}{k},G_{\alpha \beta }G^{\alpha \beta }=_0^k_{}h_k^2\frac{dk}{k},$$
$`(A.6)`$
where $`k_{}=|\eta |^1`$ is the horizon-crossing time ($`\eta ๐t/a`$). The power spectra are given as usually in the limit $`k|\eta |1`$:
$$q_k=\frac{k^{\frac{3}{2}}|\nu _k|}{2\pi a\sqrt{\gamma }},h_k=\frac{k^{\frac{3}{2}}|\nu _k^\lambda |}{\pi a},$$
$`(A.7)`$
by solving the respective Klein-Gordon equations for the functions $`\nu _k^\lambda (\eta )`$ (see (A.4), (A.5)),
$$\frac{d^2\nu _k^{(\lambda )}}{d\eta ^2}+\left(k^2U^{(\lambda )}\right)\nu _k^{(\lambda )}=0,$$
$`(A.8)`$
with $`U=d^2\left(a\sqrt{\gamma }\right)/a\sqrt{\gamma }d\eta ^2`$, $`U^\lambda =d^2a/d\eta ^2`$, and the vacuum initial conditions,
$$k|\eta |1:\nu _k^{(\lambda )}=\frac{\mathrm{exp}(ik\eta )}{\sqrt{2k}}.$$
$`(A.9)`$
In particular, for slow-roll inflation (e.g. the $`\mathrm{\Lambda }\lambda `$-inflation) we have
$$\nu _k^{(\lambda )}\frac{\mathrm{exp}(ik\eta )}{\sqrt{2k}}\left(1\frac{i}{k\eta }\right),$$
$`(A.10)`$
and
$$q_k\frac{H}{2\pi \sqrt{2\gamma }},h_k\frac{H}{\pi \sqrt{2}},k=aH1/\eta .$$
$`(A.11)`$
The $`h_k`$ spectrum takes into account both polarizations of CGWs.
In conclusion, we comment on eq.(13) recalling that $`\mathrm{\Phi }=0.6q`$ at the matter-dominated stage (see (A.3)), and
$$\frac{\mathrm{\Delta }T}{T}=\frac{\mathrm{\Phi }}{3}=\frac{q}{5}$$
$`(A.12)`$
at large angular scale ($`>1^0`$); both Newtonian and $`q`$ potentials in (A.12) are taken at the last-scattering surface.
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# รtude de la rรฉgularitรฉ analytique de lโapplication de symรฉtrie CR formelle
## ยง1. Introduction
La recherche de formes normales (cf. travaux de Chern-Moser , Moser-Webster , Webster , Huang-Krantz , Gong et Ebenfelt ) pour les sous-variรฉtรฉs analytiques rรฉelles de $`^n`$ soulรจve la question de la convergence des normalisations formelles. Moser et Webster ont donnรฉ des exemples de surfaces $`๐^\omega `$ dans $`^2`$ ร tangente complexe isolรฉe et hyperboliques au sens de Bishop, formellement mais non holomorphiquement normalisables (ร cause dโun phรฉnomรจne de petits diviseurs, voir ), mรชme lorsque la forme normale est elle-mรชme analytique ou algรฉbrique. En revanche, il apparaรฎt quโun tel phรฉnomรจne ne se produit pas pour les objets CR, dโaprรจs les rรฉsultats rรฉcents de Baouendi-Rothschild obtenus en collaboration avec Ebenfelt ou avec Zaitsev, et รฉnoncรฉs avec des hypothรจses gรฉnรฉriques de non-dรฉgรฉnรฉrescence (voir ). Ces auteurs รฉtablissent notamment quโune application formelle inversible entre deux sous-variรฉtรฉs CR-gรฉnรฉriques $`๐^\omega `$ finiment non-dรฉgรฉnรฉrรฉes (ou essentiellement finies) et minimales de $`^n`$ est convergente. On dรฉmontre ici un thรฉorรจme de convergence plus gรฉnรฉral, valable sans hypothรจse de non-dรฉgรฉnรฉrescence, et qui confirme la rigiditรฉ du cas CR. Ce rรฉsultat sโinterprรจte alors comme un principe de symรฉtrie de Schwarz formel pour les applications CR et il borne aussi le degrรฉ de transcendance de lโapplication formelle par rapport au corps des fractions rationnelles (cf. Coupet-Pinchuk-Sukhov ).
## ยง2. Application de symรฉtrie et รฉquivalences formelles de variรฉtรฉs CR
### 2.1. Application de symรฉtrie
Soit $`h:(M,p)_{}(M^{},p^{})`$ une รฉquivalence formelle entre deux sous-variรฉtรฉs analytiques rรฉelles ($`๐^\omega `$) CR-gรฉnรฉriques de $`^n`$. On suppose $`p=p^{}=0`$ dans des coordonnรฉes centrรฉes $`t^n`$, $`t^{}^n`$ et on note $`m:=\text{dim}_{CR}M=m^{}`$, $`d:=\text{codim}_{}M=d^{}`$, $`m+d=n`$. Dans un travail antรฉrieur, lโauteur a remarquรฉ lโexistence dโun invariant plus gรฉnรฉral que $`h`$, quโil convient dโappeler application de symรฉtrie et suggรฉrรฉ lโintรฉrรชt dโรฉtablir sa rรฉgularitรฉ, lorsque $`(M^{},p^{})`$ est holomorphiquement dรฉgรฉnรฉrรฉe. Cette application synthรฉtise en une seule sรฉrie formelle le โjet dโordre infini en $`\overline{\nu }^{}`$ de la variรฉtรฉ de Segre complexifiรฉe conjuguรฉe (voir ) $`\underset{ยฏ}{๐ฎ}_{h(t)}^{}`$โ: $`(\overline{\nu }^{},t)j_{\overline{\nu }^{}}^{\mathrm{}}\underset{ยฏ}{๐ฎ}_{h(t)}^{}`$, $`\overline{\nu }^{}^n`$, $`\overline{\nu }^{}\underset{ยฏ}{๐ฎ}_{h(t)}^{}`$, de la maniรจre suivante. En coordonnรฉes locales $`t^{}=(w^{},z^{})^m\times ^d`$, $`\tau ^{}=(\zeta ^{},\xi ^{})`$, telles que les $`d`$ รฉquations de la complexifiรฉe extrinsรจque $`(^{},0):=((M^{})^c,0^c)`$ sโรฉcrivent $`\xi ^{}=\mathrm{\Theta }^{}(\zeta ^{},t^{})=_{\gamma ^m}\zeta _{}^{}{}_{}{}^{\gamma }\mathrm{\Theta }_\gamma ^{}(t^{})`$, on a $`\underset{ยฏ}{๐ฎ}_{h(t)}^{}=\{(\overline{\lambda }^{},\overline{\mu }^{})\overline{\mu }^{}=\mathrm{\Theta }^{}(\overline{\lambda }^{},h(t))\}`$, oรน $`\overline{\nu }^{}=(\overline{\lambda }^{},\overline{\mu }^{})`$, et lโapplication de symรฉtrie sera par dรฉfinition la sรฉrie formelle $`d`$-vectorielle $`_h^{}(\overline{\nu }^{},t):=\overline{\mu }^{}_{\gamma ^m}\overline{\lambda ^{}}^\gamma \mathrm{\Theta }_\gamma ^{}(h(t))\{\overline{\nu }^{}\}[[t]]^d`$. Dans ces conditions, le lien de $`_h^{}(\overline{\nu }^{},t)`$ avec $`j_{\overline{\nu }^{}}^{\mathrm{}}\underset{ยฏ}{๐ฎ}_{h(t)}^{}`$ sโexprime par une collection infinie de sรฉries :
$$j_{\overline{\nu }^{}}^{\mathrm{}}S_{h(t)}^{}=(\overline{\nu }^{},\{_{\overline{\lambda }^{}}^\beta [_h^{}(\overline{\nu }^{},t)]\}_{\beta ^m})=(\overline{\nu }^{},\{_{\overline{\lambda }^{}}^\beta \left[\overline{\mu }^{}\mathrm{\Theta }^{}(\overline{\lambda }^{},h(t))\right]\}_{\beta ^m}),$$
$`2.2`$
en prenant les jets avec dรฉpendance du point base, e.g. $`j_t^k\mathrm{\Psi }(t):=(t,\{_t^\alpha \mathrm{\Psi }(t)\}_{|\alpha |k})`$. Par dรฉfinition, cette application formelle $`_h^{}`$ dรฉpend du systรจme de coordonnรฉes $`t^{}`$, mais sa convergence en est indรฉpendante (fait qui dรฉcoule de lโinvariance biholomorphe des variรฉtรฉs de Segre). Voici notre rรฉsultat principal :
###### Th\\'eor\\\`eme 2.3
Si $`(M,p)`$ est minimale, alors lโapplication de symรฉtrie $`_h^{}`$ est convergente, i.e. $`_h^{}(\overline{\nu }^{},t)\{\overline{\nu }^{},t\}^d`$.
###### Remark Remarques
(a) De maniรจre รฉquivalente : toutes les applications$`\mathrm{\Theta }_\gamma ^{}(h(t))=:\phi _\gamma ^{}(t)`$ $`(`$une infinitรฉ$`)`$ appartiennent ร $`\{t\}^d`$, et $`\epsilon ,C>0`$, $`t\epsilon \phi _\gamma ^{}(t)C^{|\gamma |+1}`$.
(b) Lโรฉtude de la rรฉgularitรฉ de $`_h^{}`$ gรฉnรฉralise adรฉquatement le principe de symรฉtrie de Schwarz en plusieurs variables. En effet, aucune condition de non-dรฉgรฉnรฉrescence nโest supposรฉe sur la variรฉtรฉ image $`(M^{},p^{})`$, comme dans le cas $`n=1`$.
(c) On supposera dans la suite $`(M,p)`$ minimale au sens de Tumanov (voir ).
Applications. Voici maintenant deux applications importantes de ce thรฉorรจme.
$``$ Premiรจrement, rappelons que $`(M^{},p^{})`$ est holomorphiquement non-dรฉgรฉnรฉrรฉe si et seulement si le rang gรฉnรฉrique de lโapplication $`t^{}(\mathrm{\Theta }_\gamma ^{}(t^{}))_{\beta ^m}`$ est รฉgal ร $`n`$ (critรจre de Stanton ). Dans ce cas, les hypothรจses du lemme suivant (qui dรฉcoule directement de ) seront satisfaites :
###### Lemme 2.4
Soient $`R(t,t^{})\{t,t^{}\}^n`$, $`t,t^{}^n`$, et $`h(t)[[t]]^n`$, $`h(0)=0`$, vรฉrifiant : $`R(t,h(t))0`$ et $`\text{dรฉt}(\frac{R_k}{y_l}(t,h(t)))_{1k,ln}0`$. Alors $`h(t)\{t\}^n`$.
En effet, dโaprรจs , il existe $`\gamma ^1,\mathrm{},\gamma ^n^m`$ tels que $`\text{dรฉt}(\frac{\mathrm{\Theta }_{\gamma ^k}^{}}{t_l^{}}(t^{}))_{1k,ln}0`$, et comme $`\text{dรฉt}(\frac{h_k(t)}{t_l})_{1k,ln}(0)0`$, lโhypothรจse est satisfaite avec $`R_k(t,t^{}):=\mathrm{\Theta }_{\gamma ^k}^{}(t^{})\phi _{\gamma ^k}^{}(t)`$, $`1kn`$, oรน $`\mathrm{\Theta }_\gamma ^{}(h(t))=:\phi _\gamma ^{}(t)\{t\}^d`$ dโaprรจs le Thรฉorรจme 2.3.
Par consรฉquent, dโaprรจs le critรจre de Stanton et le Lemme 2.4 :
###### Corollaire 2.5
Si $`(M^{},p^{})`$ est holomorphiquement non-dรฉgรฉnรฉrรฉe, $`h(t)\{t\}^n`$.
Rรฉciproquement, il est connu que si $`(M^{},p^{})`$ est holomorphiquement dรฉgรฉnรฉrรฉe, il existe $`h^{\mathrm{}}:(M^{},p^{})_{}(M^{},p^{})`$ inversible et non convergente (cf. et ยง9 ci-dessous). Le Corollaire 2.5 donne ainsi une condition nรฉcessaire et suffisante pour la convergence de lโapplication formelle $`h(t)`$.
Plus gรฉnรฉralement, notre rรฉsultat peut aussi sโinterprรฉter comme un rรฉsultat dโanalyticitรฉ partielle de $`h`$, dans lโesprit de Coupet-Pinchuk-Sukhov (voir ) :
###### Corollaire 2.6
Le degrรฉ de transcendance de lโextension de corps $`\text{Frac}(\{t\})\text{Frac}((\{t\}))(h_1(t),\mathrm{},h_n(t)))`$ $`\text{Frac}([[t]])`$ est infรฉrieur ou รฉgal au rang gรฉnรฉrique $`e^{}`$ de lโapplication holomorphe $`t^{}(\mathrm{\Theta }_\gamma ^{}(t^{}))_{\gamma ^m}`$.
En effet, le graphe formel de $`h`$ est alors contenu dans une composante irrรฉductible de dimension $`n+e^{}`$ de lโensemble analytique complexe $`\{(t,t^{})\mathrm{\Theta }_\gamma ^{}(t^{})=\phi _\gamma ^{}(t),\gamma \}`$ et le Corollaire 2.6 dรฉcoule donc de la caractรฉrisation du degrรฉ de transcendance par la dimension minimale dโun ensemble analytique contenant ce graphe (voir ).
$``$ Deuxiรจmement, dโaprรจs le thรฉorรจme dโapproximation dโArtin appliquรฉ aux รฉquations analytiques $`\mathrm{\Theta }_\gamma ^{}(h(t))=\phi _\gamma ^{}(t)\{t\}^d`$ ($`\gamma ^m`$), il existe $`H(t)\{t\}^n`$ vรฉrifiant $`\mathrm{\Theta }_\gamma ^{}(H(t))=\phi _\gamma ^{}(t)\{t\}^d`$, $`\gamma `$. Lโapplication $`tH(t)`$ รฉtablissant alors une รฉquivalence convergente (voir ยง8), on en dรฉduit :
###### Corollaire 2.7
Les variรฉtรฉs CR-gรฉnรฉriques $`๐^\omega `$ minimales $`(M,p)`$ et $`(M^{},p^{})`$ sont biholomorphes si et seulemement si elles sont formellement รฉquivalentes.
###### Remark Remarque
Ce rรฉsultat a รฉtรฉ obtenu rรฉcemment par Baouendi, Rothschild et Zaitsev dans en supposant lโapplication $`t^{}(\mathrm{\Theta }_\gamma ^{}(t^{}))_{\gamma ^m}`$ de rang constant au voisinage de $`p^{}`$, et en travaillant avec un raffinement du thรฉorรจme dโArtin dรป ร Wavrick.
### 2.8. Conditions de non-dรฉgรฉnรฉrescence
Ainsi, tout repose sur le Thรฉorรจme 2.3, mais habituellement, les thรฉorรจmes de rรฉgularitรฉ dโapplications CR sont รฉtablis avec des conditions supplรฉmentaires de non-dรฉgรฉnรฉrescence. Toutes ces conditions de non-dรฉgรฉnรฉrescence peuvent sโexprimer ร travers le morphisme des $`k`$-jets de variรฉtรฉs de Segre, introduit par Diederich et Webster (), et qui a servi ร dรฉfinir lโapplication de symรฉtrie ci-dessus. Soit $`j_t^{}^k๐ฎ_\tau ^{}^{}`$ le $`k`$-jet au point $`t^{}`$ de la variรฉtรฉ de Segre complexifiรฉe $`๐ฎ_\tau ^{}^{}=\{(w^{},z^{})z^{}=\overline{\mathrm{\Theta }}^{}(w^{},\tau ^{})\}`$, qui dรฉfinit une application holomorphe $`\phi _k^{}:^{}(t^{},\tau ^{})j_t^{}^k๐ฎ_\tau ^{}^{}=(t^{},\{_w^{}^\beta [z^{}\overline{\mathrm{\Theta }}^{}(w^{},\tau ^{})]\}_{|\beta |k})^{n+d\frac{(m+k)!}{m!k!}}`$. Soit $`p_{}^{}{}_{}{}^{c}:=(p^{},\overline{p}^{})^{}`$. Classons toutes ces conditions par ordre croissant de gรฉnรฉralitรฉ (pour lโimplication nd4 $``$ nd5, voir ). On dit que $`(M^{},p^{})`$ est :
nd1. Levi-non-dรฉgรฉnerรฉe en $`p^{}`$ si $`\phi _1^{}`$ est une immersion en $`p_{}^{}{}_{}{}^{c}`$.
nd2. Finiment non-dรฉgรฉnรฉrรฉe en $`p^{}`$ si $`\phi _k^{}`$ est immersive en $`p_{}^{}{}_{}{}^{c}`$, $`kk_0`$.
nd3. Essentiellement finie en $`p^{}`$ si $`\phi _k^{}`$ est finie en $`p_{}^{}{}_{}{}^{c}`$, $`kk_0`$.
nd4. Segre-non-dรฉgรฉnรฉrรฉe en $`p^{}`$ si $`\text{rg-gรฉn}(\phi _k^{}|_{\underset{ยฏ}{๐ฎ}_p^{}^{}})=m^{}`$, $`kk_0`$.
nd5. Holomorphiquement non-dรฉgรฉnรฉrรฉe si $`\text{rg-gรฉn}(\phi _k^{})=\text{dim}_{}^{}`$, $`kk_0`$.
En particulier, lโรฉquivalence formelle $`h`$ converge sous chacune des conditions nd1 ร nd5 ci-dessus, dโaprรจs le Corollaire 2.6, qui gรฉnรฉralise des rรฉsultats dรฉjร connus.
###### Remark Comparaison avec les rรฉsultats prรฉcรฉdents
La dรฉmonstration du cas nd1 suit implicitement des formes normales de Chern-Moser ; le cas nd2 est traitรฉ dans ; le cas nd3 dans ; le cas nd4 par lโauteur dans ; et enfin, le cas nd5 par lโauteur, dans , en codimension un (cas hypersurface). Ensuite Mir a traitรฉ le Thรฉorรจme 2.3 en codimension un (voir ). Mir avait รฉcrit auparavant un preprint traitant le Thรฉorรจme 2.3 lorsque $`(M^{},p^{})`$ est algรฉbrique (cf. aussi ), cas trรจs spรฉcial, puisque toutes les sรฉries $`\{\mathrm{\Theta }_\gamma ^{}(t^{})\}_{\gamma ^m}`$ sont algรฉbriquement dรฉpendantes sur un nombre fini dโentre elles. Cependant, mentionnons que lโutilisation dโidentitรฉs de rรฉflexion conjuguรฉes, qui constitue lโidรฉe principale dans la dรฉmonstration du prรฉsent article (voir ยง5 ci-dessous) et qui permet de traiter la codimension quelconque sans hypothรจse dโalgรฉbricitรฉ, est une idรฉe qui apparaรฎt dรฉjร clairement dans sans y avoir รฉtรฉ toutefois exploitรฉe de maniรจre appropriรฉe.
###### Remark Remarque sur les conditions de non-dรฉgรฉnรฉrescence
ร chaque รฉtape de la dรฉmonstration par rรฉcurrence du Corollaire 2.6, chacune des conditions nd1, nd2, nd3 et nd4 permet de substituer ร la considรฉration de lโinfinitรฉ de sรฉries formelles $`\{\mathrm{\Theta }_\gamma ^{}(h(t))\}_{\gamma ^m}`$ celle des $`n`$ composantes de $`h`$. Dans ce cas, il nโest pas nรฉcessaire de travailler avec les identitรฉs de rรฉflexion conjuguรฉes (voir ). Mais ici, dans le cas gรฉnรฉral, on travaillera avec cette collection infinie $`\{\mathrm{\Theta }_\gamma ^{}(h(t))\}_{\gamma ^m}`$ en utilisant de maniรจre cruciale la symรฉtrie du problรจme par conjugaison complexe.
###### Remark Remerciements
Lโauteur tient ร reconnaรฎtre ici sa dette envers les travaux de Baouendi-Rothschild et co-auteurs dont il sโest inspirรฉ.
## ยง3 Notations et prรฉliminaires
### 3.1. รquations dรฉfinissantes
Dans des coordonnรฉes $`t=(w,z)^m\times ^d`$ et $`t^{}=(w^{},z^{})^m\times ^d`$ sโannulant en $`p`$ et en $`p^{}`$ et telles que $`T_0^cM_z^d=\{0\}`$ et $`T_0^cM^{}_z^{}^d=\{0\}`$, les variรฉtรฉs $`(M,0)(_t^n,0)`$ et $`(M^{},0)(_t^{}^n,0)`$ sont donnรฉes par deux systรจmes de $`d`$ รฉquations scalaires analytiques $`z_j=\overline{\mathrm{\Theta }}_j(w,\overline{t})`$ et $`z_j^{}=\overline{\mathrm{\Theta }}_j^{}(w^{},\overline{t}^{})`$, $`1jd`$. Soient $`\tau :=(\zeta ,\xi ):=(\overline{t})^c`$ et $`\tau ^{}:=(\zeta ^{},\xi ^{}):=(\overline{t}^{})^c`$ les variables complexifiรฉes formelles auxquelles correspondent les complexifications extrinsรจques $`(,0)_t^n\times _\tau ^n`$ de $`(M,0)`$ et $`(^{},0)_t^{}^n\times _\tau ^{}^n`$ de $`(M^{},0)`$, dont voici les รฉquations donnรฉes sous les formes conjuguรฉes รฉquivalentes :
$$\{\begin{array}{cc}& z=\underset{\gamma ^m}{}w^\gamma \overline{\mathrm{\Theta }}_\gamma (\tau )\text{et}z^{}=\underset{\gamma ^m}{}w_{}^{}{}_{}{}^{\gamma }\overline{\mathrm{\Theta }}_\gamma ^{}(\tau ^{}),\hfill \\ & \xi =\underset{\gamma ^m}{}\zeta ^\gamma \mathrm{\Theta }_\gamma (t)\text{et}\xi ^{}=\underset{\gamma ^m}{}\zeta _{}^{}{}_{}{}^{\gamma }\mathrm{\Theta }_\gamma ^{}(t^{}).\hfill \end{array}$$
$`3.2`$
Quitte ร choisir des coordonnรฉes normales, on supposera $`\mathrm{\Theta }_\gamma (0)=0`$ et $`\mathrm{\Theta }_\gamma ^{}(0)=0`$, $`\gamma `$. Ici, les fonctions analytiques $`d`$-vectorielles $`\mathrm{\Theta }_\gamma (t):=_{\alpha ^n}\theta _{\gamma ,\alpha }t^\alpha \{t\}^d`$, satisfont une inรฉgalitรฉ de Cauchy : $`\mathrm{\Theta }_\gamma (t)C^{|\gamma |+1}`$ pour $`t\epsilon `$, oรน $`\epsilon ,C>0`$, et de mรชme pour $`\mathrm{\Theta }_\gamma ^{}`$. On pose $`r(t,\tau ):=z\overline{\mathrm{\Theta }}(w,\tau )`$ et $`r^{}(t^{},\tau ^{}):=z^{}\overline{\mathrm{\Theta }}^{}(w^{},\tau ^{})`$. Ainsi, $`\overline{r}(\tau ,t)=\xi \mathrm{\Theta }(\zeta ,t)`$ et $`\overline{r}^{}(\tau ^{},t^{})=\xi ^{}\mathrm{\Theta }^{}(\zeta ^{},t^{})`$. Par hypothรจse, il existe deux matrices inversibles $`a(t,\tau )\{t,\tau \}^{d\times d}`$ et $`a^{}(t^{},\tau ^{})\{t^{},\tau ^{}\}^{d\times d}`$ telles que $`r(t,\tau )a(t,\tau )\overline{r}(\tau ,t)`$ et $`r^{}(t^{},\tau ^{})a^{}(t^{},\tau ^{})\overline{r}^{}(\tau ^{},t^{})`$, avec $`a(0)=I_{d\times d}=a^{}(0)`$.
### 3.3. Application formelle
Lโapplication formelle donnรฉe $`h(t)`$ est par dรฉfinition une sรฉrie formelle vectorielle $`h(t)=(h_1(t),\mathrm{},h_n(t))[[t]]^n`$, $`h(0)=0`$ vรฉrifiant $`\text{dรฉt}(\frac{h_k}{t_l})_{1k,ln}(0)0`$ et (aprรจs complexification) $`r^{}(h(t),\overline{h}(\tau ))c(t,\tau )r(t,\tau )`$, pour une matrice (nรฉcessairement inversible) $`c(t,\tau )[[t,\tau ]]^{d\times d}`$. Afin de la distinguer dโune vรฉritable application ensembliste, on รฉcrira $`h^c=(h,\overline{h})(,0)_{}(^{},0)`$ avec lโindice $``$ pour โformelโ ou bien on รฉcrira โ$`r^{}(h(t),\overline{h}(\tau ))=0`$ lorsque $`r(t,\tau )=0`$โ. Une telle expression a un sens, puisque sur la variรฉtรฉ complexe $`(2m+d)`$-dimensionnelle $`(,0)=\{(t,\tau )r(t,\tau )=0\}`$, on peut choisir (indiffรฉremment) les coordonnรฉes $`(w,\tau )`$ ou $`(\zeta ,t)`$ et alors on entendra lโidentitรฉ $`r^{}(h(t),\overline{h}(\zeta ,\mathrm{\Theta }(\zeta ,t)))0`$ dans $`[[\zeta ,t]]^d`$. Bien entendu, on pourrait utiliser le langage adรฉquat de la thรฉorie des morphismes dโalgรจbres locales, mais les calculs prรฉsentant dรฉjร une certaine complexitรฉ, nous prรฉfรฉrerons les exposer sous une forme directe. Ainsi, lโhypothรจse $`h(,0)_{}(^{},0)`$ sera exprimรฉe par les deux รฉquations รฉquivalentes :
$$f(t)=\overline{\mathrm{\Theta }}^{}(g(t),\overline{h}(\tau ))()\overline{f}(\tau )=\mathrm{\Theta }^{}(\overline{g}(\tau ),h(t)),\text{โlorsque}r(t,\tau )=0\text{}.$$
$`3.4`$
Enfin, nous aurons besoin dโune propriรฉtรฉ de symรฉtrie par conjugaison complexe des objets et des sรฉries formelles. Soit $`h^c=(h,\overline{h})(,0)_{}(^{},0)`$ comme ci-dessus, $`h^c(t,\tau )=(h(t),\overline{h}(\tau ))`$ et posons $`\sigma (t,\tau ):=(\overline{\tau },\overline{t})`$, $`\sigma ^{}(t^{},\tau ^{}):=(\overline{\tau }^{},\overline{t}^{})`$. Alors :
$$\sigma ^{}(h^c(t,\tau ))=\sigma ^{}(h(t),\overline{h}(\tau ))=(h(\overline{\tau }),\overline{h}(\overline{t}))=h^c(\overline{\tau },\overline{t})=h^c(\sigma (t,\tau )).$$
$`3.5`$
### 3.6. Application de symรฉtrie
Lโapplication de symรฉtrie $`_h^{}(\overline{\nu }^{},t)`$, $`t^n`$, $`\overline{\nu }^{}=(\overline{\lambda }^{},\overline{\mu }^{})^m\times ^d`$ sera par dรฉfinition la sรฉrie formelle $`d`$-vectorielle :
$$_h^{}(\overline{\nu }^{},t)=\overline{\mu }^{}\underset{\gamma ^m}{}\overline{\lambda ^{}}^\gamma \mathrm{\Theta }_\gamma ^{}(h(t))\{\overline{\nu }^{}\}[[t]]^d$$
$`3.7`$
###### Remark Remarque
Soient des variables $`x_1,x_2`$. Lโanneau $`[[x_1]]\{x_2\}`$ nโa pas de sens.
###### Lemme 3.8
La propriรฉtรฉ $`_h^{}(\overline{\nu }^{},t)\{\overline{\nu }^{},t\}^d`$ est indรฉpendante du choix des coordonnรฉes $`(w^{},z^{})`$ telles que $`T_0^cM^{}_z^{}^d=\{0\}`$.
###### Demonstration Preuve
Consรฉquence aisรฉe de lโinvariance biholomorphe des variรฉtรฉs de Segre. โ
## ยง4. Convergence de $`_h^{}`$ et de ses jets sur les chaรฎnes de Segre
On note $`:=(^1,\mathrm{},^m)`$ et $`\underset{ยฏ}{}:=(\underset{ยฏ}{}^1,\mathrm{},\underset{ยฏ}{}^m)`$ des bases de $`T^{1,0}`$ et $`T^{0,1}`$ ร coefficients holomorphes qui commutent, donnรฉes par $`^j:=\frac{}{w_j}+\frac{\overline{\mathrm{\Theta }}(w,\tau )}{w_j}\frac{}{z}`$ et $`\underset{ยฏ}{}^j:=\frac{}{\zeta _j}+\frac{\mathrm{\Theta }(\zeta ,t)}{\zeta _j}\frac{}{\xi }`$, et $`_w(0)=_{w^1}^1(\mathrm{}_{w^m}^m(0))`$ le $`m`$-flot de $``$, $`w^m`$ (de mรชme pour $`\underset{ยฏ}{}_\zeta (0)`$, $`\zeta ^m`$). Clairement, $`_w(0)=(w,\overline{\mathrm{\Theta }}(w,0),0,0)^{2n}`$ et $`\underset{ยฏ}{}_\zeta (0)=(0,0,\zeta ,\mathrm{\Theta }(\zeta ,0))^{2n}`$. Les concatรฉnations alternรฉes de tels flots sont appelรฉes $`k`$-chaรฎnes de Segre, par exemple, pour $`k=2j`$, $`(w_1,\mathrm{},w_{2j})\underset{ยฏ}{}_{w_{2j}}(_{w_{2j1}}(\mathrm{}\underset{ยฏ}{}_{w_2}(_{w_1}(0))))`$, i.e. $`\underset{ยฏ}{}_{w_2}(_{w_1}(0))=(w_1,\overline{\mathrm{\Theta }}(w_1,0),w_2,\mathrm{\Theta }(w_2,`$ $`w_1,\overline{\mathrm{\Theta }}(w_1,0)))`$, etc. (voir ). On a aussi $`\sigma (_w(0))=\underset{ยฏ}{}_{\overline{w}}(0)`$, etc.. Si on convient de noter $`w_{(k)}:=(w_1,\mathrm{},w_k)`$, oรน $`w_1,\mathrm{},w_k^m`$ sont proches de $`0`$, ces $`k`$-chaรฎnes seront abrรฉgรฉes dans la suite par $`\mathrm{\Gamma }_k(w_{(k)})`$.
###### Remark Remarque
Un point de vue ensembliste รฉquivalent sur les chaรฎnes de Segre a รฉtรฉ dรฉveloppรฉ par Baouendi, Ebenfelt et Rothschild (voir ). Dans , lโauteur a ensuite gรฉomรฉtrisรฉ ce point de vue en introduisant les flots de champs de vecteurs.
On note $`_t^\kappa \chi (t):=(_t^\beta \chi (t))_{|\beta |\kappa }[[t]]^{K\frac{(n+\kappa )!}{n!\kappa !}}`$ le $`\kappa `$-jet dโune sรฉrie formelle vectorielle $`\chi (t)[[t]]^K`$, oรน $`K_{}`$, $`\kappa `$. Par exemple, $`_t^\kappa _h^{}(\overline{\nu }^{},t)=_{\gamma ^m}\overline{\lambda ^{}}^\gamma _t^\kappa [\mathrm{\Theta }_\gamma ^{}(h(t))]`$.
###### Lemme 4.1
Soit $`q(x)[[x]]^{2n}`$, avec $`q(0)=0`$. On a $`_t^\kappa h(\underset{ยฏ}{}_\zeta (q(x)))_t^\kappa h(q(x))`$ dans $`[[x,\zeta ]]^{n\frac{(n+\kappa )!}{n!\kappa !}}`$ et $`_\tau ^\kappa \overline{h}(_w(q(x)))_\tau ^\kappa \overline{h}(q(x))`$ dans $`[[x,w]]^{n\frac{(n+\kappa )!}{n!\kappa !}}`$.
###### Demonstration Preuve
On note $`q(x)=(q_1(x),q_2(x))^n\times ^n`$ et $`q_2(x):=(q_2^1(x),q_2^2(x))^m\times ^d`$. Il est facile de voir que $`\underset{ยฏ}{}_\zeta (q(x))=(q_1(x),\zeta +q_2^1(x),\mathrm{\Theta }(\zeta +q_2^1(x),q_1(x)))`$, et comme $`_t^\kappa h(t,\tau )=_t^\kappa h(t)`$, on a bien $`_t^\kappa h(\underset{ยฏ}{}_\zeta (q(x)))_t^\kappa h(q_1(x))_t^\kappa h(q(x))`$. La deuxiรจme propriรฉtรฉ dรฉcoule de la premiรจre grรขce ร la symรฉtrie par conjugaison complexe 3.5 et grรขce ร la relation $`\sigma (_w(q(x)))=\underset{ยฏ}{}_{\overline{w}}(\sigma (q(x)))`$. โ
Dโaprรจs le critรจre de minimalitรฉ de reformulรฉ dans :
###### Lemme 4.2
La variรฉtรฉ CR-gรฉnรฉrique $`๐^\omega `$ $`(M,p)`$ est minimale si et seulement si $`\mathrm{\Gamma }_k^{mk}(,0)`$ induit une submersion en $`0`$ pour $`k2d+1`$.
Ainsi, il nous suffira de dรฉmontrer que $`_h^{}(\overline{\nu }^{},\mathrm{\Gamma }_k(w_{(k)}))\{\overline{\nu }^{},w_{(k)}\}^d`$, $`k`$ pour en dรฉduire que $`_h^{}(\overline{\nu }^{},t)\{\overline{\nu }^{},t\}^d`$, comme souhaitรฉ (cf. ).
### 4.3. Estimรฉes de Cauchy
Dans cet objectif, il est nรฉcessaire dโestimer la croissance des normes $`\mathrm{\Theta }_\gamma ^{}(h(t))`$ quand $`|\gamma |\mathrm{}`$. Commenรงons par une formule standard et universelle de dรฉrivation composรฉe.
###### Lemme 4.4
Soit $`\beta ^m`$. Il existe un polynรดme universel $`d`$-vectoriel $`Q_\beta `$ tel que
$$\{\begin{array}{cc}& _t^\beta [\mathrm{\Theta }_\gamma ^{}(h(t))]=Q_\beta (\{_t^\delta h(t)\}_{\delta \beta },\{[_t^{}^\delta \mathrm{\Theta }_\gamma ^{}](h(t))\}_{\delta \beta })\hfill \\ & Q_\beta (\{_t^\delta h(t)\}_{\delta \beta },0)0.\hfill \end{array}$$
$`4.5`$
###### Demonstration Preuve
On note $`Q_\beta (\{h_j^\delta \}_{\delta \beta }^{1jn},\{\theta _{k}^{}{}_{}{}^{\delta }\}_{\delta \beta }^{1kd})`$ ce polynรดme. Il vรฉrifiera 4.5 si et seulement si $`Q_0=\theta _{}^{}{}_{}{}^{0}`$ et, par rรฉcurrence,
$$Q_{\beta +\beta ^1}:=\underset{\delta \beta }{}\underset{j=1}{\overset{n}{}}\frac{Q_\beta }{h_j^\delta }h_j^{\delta +\beta ^1}+\underset{\delta \beta }{}\underset{k=1}{\overset{d}{}}\underset{j=1}{\overset{n}{}}\frac{Q_\beta }{\theta _{k}^{}{}_{}{}^{\delta }}\theta _{k}^{}{}_{}{}^{\delta +\text{1}_j}h_j^{\beta ^1}$$
$`4.6`$
$`\beta ^1^m`$ avec $`|\beta ^1|=1`$ et oรน $`\text{1}_j=(0,\mathrm{},1,\mathrm{},0)`$ avec $`1`$ ร la $`j`$-iรจme place. โ
Si $`\mathrm{\Psi }(t^{})[[t^{}]]^n`$, on note $`[_t^\kappa \mathrm{\Psi }(h)](\mathrm{\Gamma }_k(w_{(k)})):=\{[_t^\beta \mathrm{\Psi }(h(t))]|_{t:=\mathrm{\Gamma }_k(w_{(k)})}\}_{|\beta |\kappa }`$. Maintenant, en utilisant les polynรดmes $`Q_\beta `$, lโestimรฉe de Cauchy $`_t^{}^\kappa \mathrm{\Theta }_\gamma ^{}(t^{})C_{}^{}{}_{}{}^{|\gamma |+1}`$ et le thรฉorรจme dโArtin (), on peut estimer $`\mathrm{\Theta }_\gamma ^{}(h(t))`$ quand $`|\gamma |\mathrm{}`$ :
###### Lemme 4.7
Soient $`k`$ et $`\kappa `$. Les propriรฉtรฉs suivantes sont รฉquivalentes :
###### Demonstration Preuve
On a (1) $``$ (2) et (2) $``$ (3). รtablissons lโimplication (3) $``$ (2). Soit $`\beta ^m`$, $`|\beta |\kappa `$. Soient $`C>0`$ et $`\epsilon >0`$ deux constantes positives dont la valeur pourra varier suivant le contexte. Par hypothรจse, $`[_t^\beta \mathrm{\Theta }_\gamma ^{}(h)](\mathrm{\Gamma }_k(w_{(k)}))=:\phi _\gamma ^\beta (w_{(k)})\{w_{(k)}\}^d`$, $`\gamma ^m`$. Dโaprรจs le Lemme 4.4, cette relation peut sโรฉcrire :
$$Q_\beta (\{_t^\delta h(\mathrm{\Gamma }_k(w_{(k)}))\}_{\delta \beta },\{[_t^{}^\delta \mathrm{\Theta }_\gamma ^{}](h(\mathrm{\Gamma }_k(w_{(k)})))\}_{\delta \beta })=\phi _\gamma ^\beta (w_{(k)}).$$
$`4.8`$
On lโinterprรจte comme une relation analytique satisfaite par les sรฉries formelles $`\{_t^\delta h(\mathrm{\Gamma }_k(w_{(k)}))\}_{\delta \beta }`$. Dโaprรจs le thรฉorรจme dโArtin, il existe une solution convergente de lโรฉq. 4.8, soit $`\{H^\delta (w_{(k)})\}_{\delta \beta }`$, oรน $`H^\delta (w_{(k)})\{w_{(k)}\}^n`$. Or, grรขce ร la relation $`Q_\beta (\{_t^\delta h(t)\}_{\delta \beta },0)0`$ et aux estimรฉes de Cauchy satisfaites par les diffรฉrentielles $`_t^{}^\delta \mathrm{\Theta }_\gamma ^{}(t^{})`$, $`\delta \beta `$, on en dรฉduit aisรฉment une estimรฉe de Cauchy $`Q_\beta (\{h_j^\delta \}_{\delta \beta }^{1jn},\{\theta _{k}^{}{}_{}{}^{\delta }\}_{\delta \beta }^{1kd})C^{|\gamma |+1}`$ pour $`(\{h_j^\delta \}_{\delta \beta }^{1jn},\{\theta _{k}^{}{}_{}{}^{\delta }\}_{\delta \beta }^{1kd})\epsilon `$. Par composition avec la solution convergente $`\{H^\delta (w_{(k)})\}_{\delta \beta }`$, $`H^\delta (0)=0`$, on en dรฉduit lโestimรฉe de Cauchy souhaitรฉe pour $`w_{(k)}\epsilon `$ :
$$||Q_\beta (\{H^\delta (w_{(k)})\}_{\delta \beta },[_t^{}^\delta \mathrm{\Theta }_\gamma ^{}](H^0(w_{(k)}))\}_{\delta \beta })||=||\phi _\gamma ^\beta (w_{(k)})||C^{|\gamma |+1}.\mathit{}$$
$`4.9`$
Ainsi, il suffit dโรฉtablir (3) par rรฉcurrence sur $`k`$ en admettant ร chaque รฉtape lโรฉquivalence avec (2). Cette rรฉcurrence procรจdera en deux moments (ยง6 et ยง7) :
###### \\'Etape 1
Pour tout $`k`$, on a :
$$\{\begin{array}{cc}& [_t^\kappa \mathrm{\Theta }_\gamma ^{}(h)](\mathrm{\Gamma }_k(w_{(k)}))\text{et}[_\tau ^\kappa \overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})](\mathrm{\Gamma }_k(w_{(k)}))\{w_{(k)}\}^{d\frac{(n+\kappa )!}{n!\kappa !}},\kappa ,\gamma \hfill \\ & [\mathrm{\Theta }_\gamma ^{}(h)](\mathrm{\Gamma }_{k+1}(w_{(k+1)}))\text{et}[\overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})](\mathrm{\Gamma }_{k+1}(w_{(k+1)}))\{w_{(k+1)}\}^{d\frac{(n+\kappa )!}{n!\kappa !}},\gamma .\hfill \end{array}$$
###### \\'Etape 2
En supposant lโรtape 1 vraie pour $`k1`$, pour tout $`\kappa `$, on a :
$$\{\begin{array}{cc}& [_t^\kappa \mathrm{\Theta }_\gamma ^{}(h)](\mathrm{\Gamma }_k(w_{(k)}))\text{et}[_\tau ^\kappa \overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})](\mathrm{\Gamma }_k(w_{(k)}))\{w_{(k)}\}^{d\frac{(n+\kappa )!}{n!\kappa !}},\gamma \hfill \\ & [_t^{\kappa +1}\mathrm{\Theta }_\gamma ^{}(h)](\mathrm{\Gamma }_k(w_{(k)}))\text{et}[_\tau ^{\kappa +1}\overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})](\mathrm{\Gamma }_k(w_{(k)}))\{w_{(k)}\}^{d\frac{(n+\kappa +1)!}{n!(\kappa +1)!}},\gamma .\hfill \end{array}$$
###### Remark Remarque
Il est clair que pour $`k=0`$, les hypothรจses de lโรtape 1 et celles de lโรtape 2 sont satisfaites.
## ยง5. Symรฉtrie par conjugaison complexe et dรฉrivations CR
Mais tout dโabord, si $`\beta ^m`$, on note $`|\beta |:=|\beta _1|+\mathrm{}+|\beta _m|`$ et $`\underset{ยฏ}{}^\beta :=\underset{ยฏ}{}_{}^{1}{}_{}{}^{\beta _1}\mathrm{}\underset{ยฏ}{}_{}^{m}{}_{}{}^{\beta _m}`$. Appliquant ces dรฉrivations aux deux identitรฉs $`r^{}(h(t),\overline{h}(\tau ))=0`$ et$`\overline{r}^{}(\overline{h}(\tau ),h(t))=0`$, $`t^n`$, $`\tau ^n`$, $`r(t,\tau )=0`$, on obtient deux familles infinies dโรฉquations, deux identitรฉs de rรฉflexion conjuguรฉes, qui sont satisfaites sur $``$ :
$$\{\begin{array}{cc}\hfill ():f\overline{\mathrm{\Theta }}^{}(g,\overline{h}),0& \underset{\gamma ^m}{}g^\gamma \underset{ยฏ}{}^\beta (\overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})),\beta _{}^m.\hfill \\ \hfill (\overline{}):\overline{f}=\mathrm{\Theta }^{}(\overline{g},h),\underset{ยฏ}{}^\beta \overline{f}& \underset{\gamma ^m}{}\underset{ยฏ}{}^\beta (\overline{g}^\gamma )\mathrm{\Theta }_\gamma ^{}(h),\beta _{}^m.\hfill \end{array}$$
$`5.1`$
Or, il existe une matrice $`a^{}(t^{},\tau ^{})\{t^{},\tau ^{}\}^{d\times d}`$, telle que $`a^{}(0,0)=I_{d\times d}`$ et $`r^{}(t^{},\tau ^{})a^{}(t^{},\tau ^{})\overline{r}^{}(\tau ^{},t^{})`$. On en dรฉduit dans $`\{t^{}\}[[\tau ]]^d`$ :
$$\underset{ยฏ}{}^\beta [r^{}(t^{},\overline{h}(\tau ))]=0,\beta ^m\underset{ยฏ}{}^\beta [\overline{r}^{}(\overline{h}(\tau ),t^{})]=0,\beta ^m,$$
$`5.2`$
et en particulier, lโรฉquivalence des systรจmes $`()`$ et $`(\overline{})`$, que lโon va exploiter. Mais habituellement, seul le systรจme $`(\overline{})`$ est considรฉrรฉ (cf. ).
Enfin, on utilisera plusieurs fois deux propriรฉtรฉs formelles qui dรฉcoulent de transformations linรฉaires รฉlรฉmentaires sur des systรจmes trigonaux infinis.
$``$ Premiรจrement, dโaprรจs un calcul qui utilise lโhypothรจse $`\text{dรฉt}(\frac{h_k}{t_l})_{1k,ln}(0)0`$ et qui est classique dans les travaux de Baouendi-Rothschild (cf. ), on a :
$$\{\begin{array}{cc}& [\underset{ยฏ}{}^\beta \overline{f}](t,\tau )=\underset{\gamma ^m}{}[\underset{ยฏ}{}^\beta \overline{g}^\gamma ](t,\tau )\mathrm{\Theta }_\gamma ^{}(t^{}),\beta ^m\hfill \\ & \underset{ยฏ}{\mathrm{\Omega }}_\beta (t,\tau ,^{|\beta |}\overline{h}(\tau ))=\mathrm{\Theta }_\beta ^{}(t^{})+\underset{\gamma _{}^m}{}\frac{(\beta +\gamma )!}{\beta !\gamma !}\overline{g}(\tau )^\gamma \mathrm{\Theta }_{\beta +\gamma }^{}(t^{}),\beta ^m,\hfill \end{array}$$
$`5.3`$
oรน les termes $`\underset{ยฏ}{\mathrm{\Omega }}_\beta `$ sont analytiques prรจs de $`0\times 0\times ^{|\beta |}\overline{h}(0)`$ et $`r(t,\tau )=0`$.
$``$ Deuxiรจmement, on a la rรฉsolution formelle directe du systรจme trigonal infini :
$$\{\begin{array}{cc}& \psi _\beta +\underset{\gamma _{}^m}{}\frac{(\beta +\gamma )!}{\beta !\gamma !}\overline{g}^\gamma \psi _{\beta +\gamma }=\underset{ยฏ}{\omega }_\beta ,\beta ^m\hfill \\ & \psi _\beta =\underset{ยฏ}{\omega }_\beta +\underset{\gamma _{}^m}{}(1)^\gamma \frac{(\beta +\gamma )!}{\beta !\gamma !}\overline{g}^\gamma \underset{ยฏ}{\omega }_{\beta +\gamma },\beta ^m.\hfill \end{array}$$
$`5.4`$
Posons $`\underset{ยฏ}{\mathrm{\Gamma }}_k(w_{(k)}):=\sigma (\mathrm{\Gamma }_k(\overline{w}_{(k)}))`$, i.e. $`\underset{ยฏ}{\mathrm{\Gamma }}_1(w_1)=\underset{ยฏ}{}_{w_1}(0)`$, $`\underset{ยฏ}{\mathrm{\Gamma }}_2(w_{(2)})=_{w_2}(\underset{ยฏ}{}_{w_1}(0))`$, etc. Alors on a dans $`\{w_{(k)}\}^{d\frac{(n+\kappa )!}{n!\kappa !}}`$ :
$$\overline{[_t^\kappa \mathrm{\Theta }_\gamma ^{}(h)](\mathrm{\Gamma }_k(w_{(k)}))}=\{\begin{array}{cc}& [_\tau ^\kappa \overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})](\underset{ยฏ}{\mathrm{\Gamma }}_k(\overline{w_{(k)}})),\text{si}k\text{est impair},\hfill \\ & [_\tau ^\kappa \overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})](\underset{ยฏ}{\mathrm{\Gamma }}_{k1}(\overline{w_{(k1)}})),\text{si}k\text{est pair}.\hfill \end{array}$$
$`5.5`$
## ยง6. Saut ร la chaรฎne de Segre supรฉrieure
###### Demonstration Preuve de lโรtape 1
On traite seulement le cas $`k`$ pair (le cas $`k`$ impair est similaire et sโy ramรจne formellement grรขce ร la symรฉtrie par conjugaison complexe 3-5.5). Tout dโabord, comme $`k`$ est pair, on a dรฉjร : $`[\overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})](\mathrm{\Gamma }_{k+1}(w_{(k+1)}))=[\overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})](\mathrm{\Gamma }_k(w_{(k)}))]\{w_{(k)}\}^{d\frac{(n+\kappa )!}{n!\kappa !}}`$, $`\gamma `$, cโest-ร -dire la deuxiรจme moitiรฉ de la conclusion de lโรtape 1. Plus gรฉnรฉralement, en appliquant encore le Lemme 4.1, on a : $`[_\tau ^\kappa \overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})](\mathrm{\Gamma }_{k+1}(w_{(k+1)}))=[_\tau ^\kappa \overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})](\mathrm{\Gamma }_k(w_{(k)}))]\{w_{(k)}\}^{d\frac{(n+\kappa )!}{n!\kappa !}}`$, convergent par hypothรจse. Ensuite, les coefficients de $`\underset{ยฏ}{}`$ รฉtant analytiques, on a facilement :
###### Lemme 6.1
Il existe $`P_\beta `$ analytique tq. $`\underset{ยฏ}{}^\beta (\overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h}(\tau )))=P_\beta (t,\tau ,[_\tau ^{|\beta |}\overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})])(\tau )`$.
Par consรฉquent, tous les termes suivants sont convergents ($`\gamma ^m`$) :
$`[\underset{ยฏ}{}^\beta \overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})](\mathrm{\Gamma }_{k+1}(w_{(k+1)}))=P_\beta (\mathrm{\Gamma }_{k+1}(w_{(k+1)}),[_\tau ^{|\beta |}\overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})](\mathrm{\Gamma }_k(w_{(k)})))\{w_{(k+1)}\}^d.`$ $`6.2`$
Le thรฉorรจme dโArtin sโapplique donc aux รฉquations $`()`$ รฉcrites au point $`(t,\tau ):=\mathrm{\Gamma }_{k+1}(w_{(k+1)})`$, dont les coefficients sont analytiques grรขce ร 6.2, รฉquations qui sont satisfaites par la sรฉrie formelle $`h(\mathrm{\Gamma }_{(k+1)}(w_{(k+1)}))`$. Ainsi, il existe une solution convergente de ces รฉqs. $`()`$ que lโon note $`H(w_{(k+1)})\{w_{(k+1)}\}^n`$. Dโaprรจs lโรฉquivalence 5.2 des systรจmes $`()`$ et $`(\overline{})`$, cette solution satisfait aussi le systรจme :
$$\{\begin{array}{cc}& \overline{f}(\mathrm{\Gamma }_k(w_{(k)}))\underset{\gamma ^m}{}\overline{g}^\gamma (\mathrm{\Gamma }_k(w_{(k)}))\mathrm{\Theta }_\gamma ^{}(H(w_{(k+1)})),\hfill \\ & [\underset{ยฏ}{}^\beta \overline{f}](\mathrm{\Gamma }_{k+1}(w_{(k+1)}))\underset{\gamma ^m}{}[\underset{ยฏ}{}^\beta \overline{g}^\gamma ](\mathrm{\Gamma }_{k+1}(w_{(k+1)}))\mathrm{\Theta }_\gamma ^{}(H(w_{(k+1)})),\beta _{}^m.\hfill \end{array}$$
$`6.3`$
Aprรจs la transformation linรฉaire 5.3 sur ce sytรจme et sur $`(\overline{})`$, on a
$$\{\begin{array}{cc}& \mathrm{\Theta }_\beta ^{}(H(w_{(k+1)}))+\underset{\gamma _{}^m}{}\frac{(\beta +\gamma )!}{\beta !\gamma !}\overline{g}^\gamma (\overline{\mathrm{\Gamma }}_k(w_{(k)}))\mathrm{\Theta }_{\beta +\gamma }^{}(H(w_{(k+1)}))=\hfill \\ & =\underset{ยฏ}{\mathrm{\Omega }}_\beta (\mathrm{\Gamma }_{k+1}(w_{(k+1)}),_\tau ^{|\beta |}\overline{h}(\overline{\mathrm{\Gamma }}_k(w_{(k)})))=(\beta ^m)\hfill \\ & =\mathrm{\Theta }_\beta ^{}(h(\mathrm{\Gamma }_{k+1}(w_{(k+1)})))+\underset{\gamma _{}^m}{}\frac{(\beta +\gamma )!}{\beta !\gamma !}\overline{g}^\gamma (\overline{\mathrm{\Gamma }}_k(w_{(k)}))\mathrm{\Theta }_{\beta +\gamma }^{}(h(\mathrm{\Gamma }_{k+1}(w_{(k+1)}))).\hfill \end{array}$$
$`6.4`$
Pour terminer, on applique aux รฉqs. 6.4 la rรฉsolution 5.4. On en dรฉduit que $`\mathrm{\Theta }_\beta ^{}(h(\mathrm{\Gamma }_{(k+1)}(w_{(k+1)})))\mathrm{\Theta }_\beta ^{}(H(w_{(k+1)}))\{w_{(k+1)}\}^d`$ converge, $`\beta ^m`$. Cโest la premiรจre moitiรฉ de la conclusion de lโรtape 1. โ
## ยง7. Rรฉcurrence des jets sur une chaรฎne de Segre
###### Demonstration Preuve de lโรtape 2
On traite seulement le cas $`k`$ impair (le cas $`k`$ pair est similaire et sโy ramรจne par symรฉtrie en utilisant 3-5.5). Tout dโabord, comme $`k`$ est impair, on a dรฉjร $`[_\tau ^{\kappa +1}\overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})](\mathrm{\Gamma }_k(w_{(k)}))=[_\tau ^{\kappa +1}\overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})](\mathrm{\Gamma }_{k1}(w_{(k1)}))]\{w_{(k1)}\}^{d\frac{(n+\kappa +1)!}{n!(\kappa +1)!}}`$ convergent $`\gamma `$, puisque lโon suppose vraie lโรtape 1 pour $`k1`$. Cโest la deuxiรจme moitiรฉ de la conclusion de lโรtape 2. En vรฉritรฉ, on va effectuer un calcul direct qui gรฉnรฉralise lโรtape 1 du ยง6 pour montrer $`[_t^\kappa \mathrm{\Theta }_\gamma ^{}(h)](\mathrm{\Gamma }_k(w_{(k)}))\{w_{(k)}\}^{d\frac{(n+\kappa )!}{n!\kappa !}}`$, $`\kappa `$, en faisant au passage une rรฉcurrence sur les jets (i.e. sur $`\kappa `$) du type de lโรtape 2 (voir la preuve du Lemme 7.20).
Soient donc $`\xi ^d`$ et $`w_{(k)}^{mk}`$. Soit $`\mathrm{{\rm Y}}`$ le $`d`$-champ de vecteurs tangent ร $``$ dรฉfini par $`\mathrm{{\rm Y}}:=\frac{}{z}+\frac{\mathrm{\Theta }(\zeta ,t)}{z}\frac{}{\xi }`$ et soit $`\underset{ยฏ}{\mathrm{{\rm Y}}}:=\frac{}{\xi }+\frac{\overline{\mathrm{\Theta }}(w,\tau )}{\xi }\frac{}{z}`$. On a $`[,\underset{ยฏ}{\mathrm{{\rm Y}}}]=0`$.
###### Lemme 7.1
Les propriรฉtรฉs suivantes sont รฉquivalentes :
###### Demonstration Preuve
Comme $`=\frac{}{w}+\frac{\overline{\mathrm{\Theta }}(w,\tau )}{w}\frac{}{z}`$, $`\underset{ยฏ}{\mathrm{{\rm Y}}}=\frac{}{\xi }+\frac{\overline{\mathrm{\Theta }}(w,\tau )}{\xi }\frac{}{z}`$, $`_t=(_w,_z)`$ et $`\frac{\overline{\mathrm{\Theta }}(0,0)}{\xi }=I_{d\times d}`$, on a (a) $``$ (b) par transformations linรฉaires. De plus, on a (c) $``$ (b), car $`_{w_k}^\delta [\mathrm{\Psi }(\mathrm{\Gamma }_k(w_{(k)}))]=_{w_k}^\delta [\mathrm{\Psi }(_{w_k}(\mathrm{\Gamma }_{k1}(w_{(k1)})))]=[^\delta \mathrm{\Psi }](\mathrm{\Gamma }_k(w_{(k)}))`$. โ
Il suffit donc de prouver (c). Soit $`(\xi ,p)\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (p)`$ le $`d`$-flot de $`\underset{ยฏ}{\mathrm{{\rm Y}}}`$. Bien sรปr, on a $`_\xi ^\alpha [\mathrm{\Psi }(\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (q(x)))]=[\underset{ยฏ}{\mathrm{{\rm Y}}}^\alpha \mathrm{\Psi }](\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (q(x)))`$. Considรฉrons $`\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)}))`$. On pose :
$$\{\begin{array}{cc}& E_\beta (w_{(k)},\xi ,t^{}):=\underset{\gamma ^m}{}w_{}^{}{}_{}{}^{\gamma }[\underset{ยฏ}{}^\beta \overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})](\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)}))),\hfill \\ & F_\beta (w_{(k)},\xi ,t^{}):=[\underset{ยฏ}{}^\beta \overline{f}](\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)})))\underset{\gamma ^m}{}[\underset{ยฏ}{}^\beta \overline{g}^\gamma ](\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)})))\mathrm{\Theta }_\gamma ^{}(t^{}),\hfill \end{array}$$
$`7.2`$
pour tout $`\beta ^m`$ (cf. $`()`$-$`(\overline{})`$). Dโaprรจs la dรฉfinition de $`a^{}(t^{},\tau ^{})`$ (voir ยง3.1), on a :
$$\{\begin{array}{cc}& E_0(w_{(k)},\xi ,t^{})=a^{}(t^{},\overline{h}(\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)}))))F_0(w_{(k)},\xi ,t^{}),\hfill \\ & F_0(w_{(k)},\xi ,t^{})=\overline{a}^{}(\overline{h}(\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)}))),t^{})E_0(w_{(k)},\xi ,t^{}).\hfill \end{array}$$
$`7.3`$
Appliquant toutes les dรฉrivations $`\underset{ยฏ}{}^\beta `$ aux รฉqs. 7.3 :
$$\{\begin{array}{cc}& E_\beta (w_{(k)},\xi ,t^{})a^{}(t^{},w_{(k)},\xi )F_\beta (w_{(k)},\xi ,t^{})+\underset{\delta <\beta }{}a_{\delta }^{}{}_{}{}^{\beta }(t^{},w_{(k)},\xi )F_\delta (w_{(k)},\xi ,t^{}),\hfill \\ & F_\beta (w_{(k)},\xi ,t^{})b^{}(t^{},w_{(k)},\xi )E_\beta (w_{(k)},\xi ,t^{})+\underset{\delta <\beta }{}b_{\delta }^{}{}_{}{}^{\beta }(t^{},w_{(k)},\xi )E_\delta (w_{(k)},\xi ,t^{}).\hfill \end{array}$$
$`7.4`$
Ici, $`a_{\delta }^{}{}_{}{}^{\beta }(t^{},w_{(k)},\xi )`$, $`b_{\delta }^{}{}_{}{}^{\beta }(t^{},w_{(k)},\xi )\{t^{}\}[[w_{(k)},\xi ]]^{d\times d}`$. Dโaprรจs 7.4, on a (cf. 5.2) $`\text{Idรฉal}\{E_\beta (w_{(k)},\xi ,t^{})\}_{\beta ^m}=\text{Idรฉal}\{F_\beta (w_{(k)},\xi ,t^{})\}_{\beta ^m}`$. Plus gรฉnรฉralement, par rรฉcurrence sur $`\alpha ^d`$, on dรฉfinit une collection $`\{E_\beta ^{(\alpha )}\}_{\alpha ^d,\beta ^m}`$ de fonctions $`d`$-vectorielles comme suit. Soit $`\alpha ^1^d`$ avec $`|\alpha ^1|=1`$. On pose $`T_0^{}=t^{}`$ et :
$$\{\begin{array}{cc}& E_\beta ^{(0)}(w_{(k)},\xi ,t^{}):=E_\beta (w_{(k)},\xi ,t^{});\text{et :}E_\beta ^{(\alpha +\alpha ^1)}(w_{(k)},\xi ,\{T_\alpha ^{}^{}\}_{\alpha ^{}\alpha +\alpha ^1}):=\hfill \\ & :=\frac{E_\beta ^{(\alpha )}}{\xi ^{\alpha ^1}}(w_{(k)},\xi ,\{T_\alpha ^{}^{}\}_{\alpha ^{}\alpha })+\underset{\alpha ^{}\alpha }{}\frac{E_\beta ^{(\alpha )}}{T_\alpha ^{}^{}}(w_{(k)},\xi ,\{T_\alpha ^{}^{}\}_{\alpha ^{}\alpha })T_{\alpha ^{}+\alpha ^1}^{}.\hfill \end{array}$$
$`7.5`$
On dรฉfinit aussi la collection similaire $`\{F_\beta ^{(\alpha )}\}_{\alpha ^d,\beta ^m}`$. Par construction :
$$\{\begin{array}{cc}& \left[E_\beta ^{(\alpha )}(w_{(k)},\xi ,\{T_\alpha ^{}^{}\}_{\alpha ^{}\alpha })\right]|{}_{T_\alpha ^{}^{}:=[\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi ^\alpha ^{}h](\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)}))),\alpha ^{}\alpha }{}^{}=\hfill \\ & =_\xi ^\alpha [E_\beta (w_{(k)},\xi ,h(\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)}))))].\hfill \end{array}$$
$`7.6`$
Voici maintenant une propriรฉtรฉ gรฉnรฉralisant 7.4 qui se vรฉrifie par un calcul formel direct en utilisant les relations 7.3 et les dรฉfinitions 7.5 de $`E_\beta ^{(\alpha )}`$ et de $`F_\beta ^{(\alpha )}`$ : โ
###### Lemme 7.7
Dans lโanneau $`\{\{T_\alpha ^{}^{}\}_{\alpha ^{}\alpha }\}[[w_{(k)},\xi ]]^d`$, on a pour tout $`\alpha ^d`$ :
$$\text{Idรฉal}\{E_\beta ^{(\alpha )}(w_{(k)},\xi ,\{T_\alpha ^{}^{}\}_{\alpha ^{}\alpha })\}_{\beta ^m}=\text{Idรฉal}\{F_\beta ^{(\alpha )}(w_{(k)},\xi ,\{T_\alpha ^{}^{}\}_{\alpha ^{}\alpha })\}_{\beta ^m}.$$
$`7.8`$
Suite de la dรฉmonstration. Dโaprรจs $`()`$, la sรฉrie formelle $`h(\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)})))`$ est une solution des รฉquations $`E_\beta (w_{(k)},\xi ,`$ $`h(\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)}))))0`$, $`\beta ^m`$. Par consรฉquent :
$$\{\begin{array}{cc}& 0_\xi ^\alpha |_{\xi =0}[E_\beta (w_{(k)},\xi ,h(\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)}))))]=\hfill \\ & =E_\beta ^{(\alpha )}(w_{(k)},0,\{[\underset{ยฏ}{\mathrm{{\rm Y}}}^\alpha ^{}h](\mathrm{\Gamma }_k(w_{(k)}))\}_{\alpha ^{}\alpha }),\alpha ^d,\beta ^m.\hfill \end{array}$$
$`7.9`$
Maintenant, on fixe $`\alpha ^d`$ et on considรจre le sous-systรจme fini :
$$E_\beta ^{(\alpha ^{})}(w_{(k)},0,\{[\underset{ยฏ}{\mathrm{{\rm Y}}}^{\alpha ^{\prime \prime }}h](\mathrm{\Gamma }_k(w_{(k)}))\}_{\alpha ^{\prime \prime }\alpha ^{}})=0,\alpha ^{}\alpha .$$
$`7.10`$
###### Lemme 7.11
Les รฉquations 7.10 sont analytiques, i.e. :
$$E_\beta ^{(\alpha ^{})}(w_{(k)},0,\{T_{\alpha ^{\prime \prime }}\}_{\alpha ^{\prime \prime }\alpha ^{}})\{w_{(k)},\{T_{\alpha ^{\prime \prime }}\}_{\alpha ^{\prime \prime }\alpha ^{}}\}^d,\alpha ^{}\alpha ,$$
$`7.12`$
###### Demonstration Preuve
En effet, comme $`k`$ est impair, on a $`\mathrm{\Gamma }_k(w_{(k)})=_{w_k}(\mathrm{\Gamma }_{k1}(w_{(k1)}))`$, dโoรน $`[_\tau ^\kappa \overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})](\mathrm{\Gamma }_k(w_{(k)}))[_\tau ^\kappa \overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})](\mathrm{\Gamma }_{k1}(w_{(k1)}))\{w_{(k1)}\}^{d\frac{(n+\kappa )!}{n!\kappa !}}`$, puisquโon suppose vraie lโรtape 1 pour $`k1`$. Par consรฉquent, en appliquant le Lemme 6.1, on voit que les dรฉrivรฉes $`_\xi ^{\alpha ^{\prime \prime }}|_{\xi =0}[[\underset{ยฏ}{}^\beta \overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h})](\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)})))]\{w_{(k)}\}^d`$ des coefficients de $`E_\beta ^{(\alpha ^{})}`$ (cf. รฉq. 7.2$`^{1^{i\stackrel{`}{e}re}}`$) par rapport ร $`\{T_{\alpha ^{\prime \prime }}\}_{\alpha ^{\prime \prime }\alpha }`$ convergent toutes. โ
Ainsi, il existe des solutions $`H_\alpha ^{}(w_{(k)})\{w_{(k)}\}^n`$, $`\alpha ^{}\alpha `$, satisfaisant :
$$E_\beta ^{(\alpha ^{})}(w_{(k)},0,\{H_{\alpha ^{\prime \prime }}(w_{(k)})\}_{\alpha ^{\prime \prime }\alpha ^{}})0,\alpha ^{}\alpha .$$
$`7.13`$
Grรขce ร la propriรฉtรฉ 7.8, on dรฉduit de 7.13 :
$$F_\beta ^{(\alpha ^{})}(w_{(k)},0,\{H_{\alpha ^{\prime \prime }}(w_{(k)})\}_{\alpha ^{\prime \prime }\alpha ^{}})0,\alpha ^{}\alpha .$$
$`7.14`$
Mais on a aussi dโun autre cรดtรฉ en dรฉrivant 7.2$`^{2^{i\stackrel{`}{e}me}}`$ par rapport ร $`\xi `$ :
$$F_\beta ^{(\alpha ^{})}(w_{(k)},0,\{[\underset{ยฏ}{\mathrm{{\rm Y}}}^{\alpha ^{\prime \prime }}h](\mathrm{\Gamma }_k(w_{(k)}))\}_{\alpha ^{\prime \prime }\alpha ^{}})0,\alpha ^{}\alpha ,$$
$`7.15`$
une identitรฉ que lโon peut rรฉรฉcrire plus explicitement comme suit :
$$\{\begin{array}{cc}& _\xi ^\alpha ^{}|_{\xi =0}[[\underset{ยฏ}{}^\beta \overline{f}](\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)})))]\underset{\gamma ^m}{}\underset{\alpha ^{\prime \prime }\alpha ^{}}{}\frac{\alpha ^{}!}{\alpha ^{\prime \prime }!(\alpha ^{}\alpha ^{\prime \prime })!}\hfill \\ & _\xi ^{\alpha ^{}\alpha ^{\prime \prime }}[[\underset{ยฏ}{}^\beta \overline{g}^\gamma ](\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)})))]_\xi ^{\alpha ^{\prime \prime }}[\mathrm{\Theta }_\gamma ^{}(h(\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)}))))]|_{\xi =0},\alpha ^{}\alpha .\hfill \end{array}$$
$`7.16`$
Or, il existe clairement des fonctions $`\mathrm{\Theta }_{\gamma }^{}{}_{}{}^{\alpha ^{\prime \prime }}`$ analytiques telles que :
$$\mathrm{\Theta }_{\gamma }^{}{}_{}{}^{\alpha ^{\prime \prime }}(\{[\underset{ยฏ}{\mathrm{{\rm Y}}}^{\alpha ^{\prime \prime \prime }}h](\mathrm{\Gamma }_k(w_{(k)}))\}_{\alpha ^{\prime \prime \prime }\alpha ^{\prime \prime }}):=_\xi ^{\alpha ^{\prime \prime }}|_{\xi =0}[\mathrm{\Theta }_\gamma ^{}(h(\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)}))))].$$
$`7.17`$
Par consรฉquent, on peut rรฉรฉcrire les รฉquations 7.14 et 7.16 comme suit :
$$\{\begin{array}{cc}& _\xi ^\alpha ^{}|_{\xi =0}[[\underset{ยฏ}{}^\beta \overline{f}](\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)})))]\underset{\gamma ^m}{}\underset{\alpha ^{\prime \prime }\alpha ^{}}{}\frac{\alpha ^{}!}{\alpha ^{\prime \prime }!(\alpha ^{}\alpha ^{\prime \prime })!}\hfill \\ & _\xi ^{\alpha ^{}\alpha ^{\prime \prime }}|_{\xi =0}[[\underset{ยฏ}{}^\beta \overline{g}^\gamma ](\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)})))]\mathrm{\Theta }_{\gamma }^{}{}_{}{}^{\alpha ^{\prime \prime }}(\{H_{\alpha ^{\prime \prime \prime }}(w_{(k)})\}_{\alpha ^{\prime \prime \prime }\alpha ^{\prime \prime }}),\alpha ^{}\alpha .\hfill \end{array}$$
$`7.18`$
$$\{\begin{array}{cc}& _\xi ^\alpha ^{}|_{\xi =0}[[\underset{ยฏ}{}^\beta \overline{f}](\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)})))]\underset{\gamma ^m}{}\underset{\alpha ^{\prime \prime }\alpha ^{}}{}\frac{\alpha ^{}!}{\alpha ^{\prime \prime }!(\alpha ^{}\alpha ^{\prime \prime })!}\hfill \\ & _\xi ^{\alpha ^{}\alpha ^{\prime \prime }}|_{\xi =0}[[\underset{ยฏ}{}^\beta \overline{g}^\gamma ](\underset{ยฏ}{\mathrm{{\rm Y}}}_\xi (\mathrm{\Gamma }_k(w_{(k)})))]\mathrm{\Theta }_{\gamma }^{}{}_{}{}^{\alpha ^{\prime \prime }}(\{[\underset{ยฏ}{\mathrm{{\rm Y}}}^{\alpha ^{\prime \prime \prime }}h](\mathrm{\Gamma }_k(w_{(k)}))\}_{\alpha ^{\prime \prime \prime }\alpha ^{\prime \prime }})0.\hfill \end{array}$$
$`7.19`$
On va maintenant utiliser 5.3-4 afin de dรฉduire des รฉquations 7.18-19 :
###### Lemme 7.20
Pour tout $`\gamma ^m`$ et tout $`\alpha ^{}\alpha `$, on a :
$$\mathrm{\Theta }_{\gamma }^{}{}_{}{}^{\alpha ^{}}(\{[\underset{ยฏ}{\mathrm{{\rm Y}}}^{\alpha ^{\prime \prime }}h](\mathrm{\Gamma }_k(w_{(k)}))\}_{\alpha ^{\prime \prime }\alpha ^{}})\mathrm{\Theta }_{\gamma }^{}{}_{}{}^{\alpha ^{}}(\{H_{\alpha ^{\prime \prime }}(w_{(k)})\}_{\alpha ^{\prime \prime }\alpha ^{}})\{w_{(k)}\}^d.$$
$`7.21`$
###### Demonstration Preuve
En appliquant directement 5.3-4, ceci est vrai pour $`|\alpha ^{}|=0`$ en considรฉrant les รฉqs. 7.18-19 seulement au rang $`\alpha ^{}=0`$, comme ร la fin du ยง6. Supposons par rรฉcurrence que 7.21 est vrai pour $`\alpha ^{}<\alpha `$, $`|\alpha ^{}|=\kappa _{}`$. Soit $`\alpha _0^{}\alpha `$, $`|\alpha _0^{}|=\kappa +1`$. On รฉcrit les รฉqs. 7.18-19 au rang $`\alpha ^{}:=\alpha _0^{}`$ et on les soustrait deux ร deux. Grรขce ร cette hypothรจse de rรฉcurrence, on obtient :
$$\{\begin{array}{cc}& \underset{\gamma ^m}{}[\underset{ยฏ}{}^\beta \overline{g}^\gamma ](\mathrm{\Gamma }_{(k)}(w_{(k)}))[\mathrm{\Theta }_{\gamma }^{}{}_{}{}^{\alpha _0^{}}(\{H_{\alpha ^{\prime \prime }}(w_{(k)})\}_{\alpha ^{\prime \prime }\alpha _0^{}})\hfill \\ & \mathrm{\Theta }_{\gamma }^{}{}_{}{}^{\alpha _0^{}}(\{[\underset{ยฏ}{\mathrm{{\rm Y}}}^{\alpha ^{\prime \prime }}h](\mathrm{\Gamma }_k(w_{(k)}))\}_{\alpha ^{\prime \prime }\alpha _0^{}})]0,\hfill \end{array}$$
$`7.22`$
pour tout $`\beta ^m`$. Dโaprรจs 5.3-4, on a alors 7.21 pour $`\alpha ^{}=\alpha _0^{}`$. Ainsi, lโรฉq. 7.21 conclut que $`[\underset{ยฏ}{\mathrm{{\rm Y}}}^\alpha \mathrm{\Theta }_\gamma ^{}(h)](\mathrm{\Gamma }_k(w_{(k)}))\{w_{(k)}\}^d`$, $`\alpha ^d`$. โ
###### Conclusion 7.23
Grรขce au Lemme 7.1, on a ainsi achevรฉ dโรฉtablir la premiรจre (et la seconde) moitiรฉ de lโรtape 2 (cas $`k`$ impair). En appliquant le Lemme 4.7 pas ร pas dans le procรฉdรฉ de rรฉcurrence en deux moments dรฉfini par lโรtape 1 et par lโรtape 2, on en conclut que $`_h^{}(\overline{\nu }^{},\mathrm{\Gamma }_k(w_{(k)}))\{\overline{\nu }^{},w_{(k)}\}^d`$, $`k`$, et donc $`_h^{}(\overline{\nu }^{},t)\{\overline{\nu }^{},t\}^d`$ par minimalitรฉ de $`(,0)`$.
La dรฉmonstration du Thรฉorรจme 2.3 est terminรฉe. โ
## ยง8. รquivalences formelles et รฉquivalences holomorphes
Lโรฉnoncรฉ suivant prรฉcise le contenu du Corollaire 2.7 :
###### Th\\'eor\\\`eme 8.1
Soit $`h(M,0)_{}(M^{},0)`$ une รฉquivalence formelle entre sous-variรฉtรฉs CR-gรฉnรฉriques minimales $`๐^\omega `$ de $`^n`$. Pour tout $`N_{}`$, il existe une รฉquivalence holomorphe $`H_N(M,0)_{}(M^{},0)`$ avec $`H_N(t)h(t)(\text{mod}t^N)`$.
###### Demonstration Preuve
Dโaprรจs le Thรฉorรจme 2.3, il existe $`\phi _\gamma ^{}(t)\{t\}^d`$ tel que $`\mathrm{\Theta }_\gamma ^{}(h(t))\phi _\gamma ^{}(t)`$, $`\gamma ^m`$. Soit $`N_{}`$. Appliquant le thรฉorรจme dโArtin, on obtient une application holomorphe $`H_N`$ satisfaisant $`\mathrm{\Theta }_\gamma ^{}(H_N(t))\phi _\gamma ^{}(t)`$, $`\gamma _m`$ et $`H_N(t)h(t)(\text{mod}t^N)`$. Bien sรปr, on a lโestimรฉe de Cauchy $`\mathrm{\Theta }_\gamma ^{}(H_N(t))C^{|\gamma |+1}`$ et $`\overline{f}(\zeta ,\mathrm{\Theta }(\zeta ,t))_{\gamma ^m}\overline{g}^\gamma (\zeta ,\mathrm{\Theta }(\zeta ,t))\mathrm{\Theta }_\gamma ^{}(H_N(t))`$. Rappelons $`r^{}(H_N(t),\overline{h}(\tau ))a^{}(H_N(t),\overline{h}(\tau ))\overline{r}^{}(\overline{h}(\tau ),H_N(t))`$, dโoรน $`F_N(t)_{\gamma ^m}G_N^\gamma (t)\overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h}(\zeta ,\mathrm{\Theta }(\zeta ,t)))`$, qui รฉquivaut encore ร $`F_N(w,\overline{\mathrm{\Theta }}(w,\tau )))_{\gamma ^m}G_N^\gamma (w,\overline{\mathrm{\Theta }}(w,\tau ))\overline{\mathrm{\Theta }}_\gamma ^{}(\overline{h}(\tau ))`$, et finalement $`F_N(w,\overline{\mathrm{\Theta }}(w,\tau ))_{\gamma ^m}G_N^\gamma (w,\overline{\mathrm{\Theta }}(w,\tau ))\overline{\mathrm{\Theta }}_\gamma ^{}(\overline{H}_N(\tau ))`$. En conclusion, lโapplication $`(H_N,\overline{H}_N)(,0)(^{},0)`$ รฉtablit une รฉquivalence convergente. โ
## ยง9 Nรฉcessitรฉ de la non-dรฉgรฉnรฉrescence holomorphe
On rappelle quโune hypersurface $`๐^\omega `$ $`(M^{},p^{})`$ est holomorphiquement dรฉgรฉnรฉrรฉe $`p^{}`$ si et seulement si il existe un germe de champ de vecteurs $`L^{}=_{j=1}^na_j^{}(t^{})\frac{}{t_j^{}}`$ ร coefficients holomorphes non tous nuls, tangent ร $`(M^{},p^{})`$ (voir ). La nรฉcessitรฉ de la non-dรฉgรฉnรฉrescence holomorphe pour la rรฉgularitรฉ de $`h`$ a รฉtรฉ รฉtablie dans pour les diffรฉomorphismes CR $`๐^{\mathrm{}}`$, mais lโauteur ne connaรฎt pas de rรฉfรฉrence publiรฉe pour la nรฉcessitรฉ dans le cas formel. En voici une dรฉmonstration brรจve utilisant .
###### Proposition 9.1
Il existe $`\varpi ^{}(t^{})[[t^{}]]\backslash \{t^{}\}`$, $`\varpi ^{}(0)=0`$, tel que le flot $`^nt^{}\mathrm{exp}(\varpi ^{}(t^{})L^{})(t^{})^n`$ induit une auto-application formelle inversible non-convergente $`h^{\mathrm{}}(M^{},0)_{}(M^{},0)`$.
###### Demonstration Preuve
Soit $`\phi ^{}(t^{},u^{})\mathrm{exp}(u^{}L^{})(t^{})=\phi ^{}(t^{},u^{})`$, le flot local de $`L^{}`$, qui est holomorphe en $`t^{}^n`$ et $`u^{}`$, pour $`t^{}`$, $`|u^{}|\epsilon `$, $`\epsilon >0`$. Ce flot satisfait $`\phi ^{}(t^{},0)t^{}`$ et $`_u^{}\phi _k^{}(t^{},u^{})a_k^{}(\phi ^{}(t^{},u^{}))`$. Comme $`L^{}0`$, on a $`_u^{}\phi ^{}(t^{},u^{})0`$. On peut supposer $`_u^{}\phi _1^{}(t^{},u^{})0`$. Soit $`\varpi ^{}(t^{})[[t^{}]]\backslash \{t^{}\}`$, $`\varpi ^{}(0)=0`$, une sรฉrie formelle non convergente, satisfaisant de plus $`_u^{}\phi _1^{}(t^{},\varpi ^{}(t^{}))0`$ dans $`[[t^{}]]`$ (il en existe beaucoup). Si la sรฉrie formelle $`h^{\mathrm{}}t^{}_{}\phi ^{}(t^{},\varpi ^{}(t^{}))`$ รฉtait convergente, alors $`t^{}_{}\varpi ^{}(t^{})`$ le serait aussi (par le Lemme 2.4), contrairement au choix de $`\varpi ^{}`$. Enfin, $`L^{}`$ รฉtant tangent ร $`(M^{},0)`$, il est clair que $`h^{\mathrm{}}(M^{},0)_{}(M^{},0)`$. โ
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# References
Extreme black hole entropy obtained in an operational approach
Bin Wang<sup>a,b,</sup><sup>1</sup><sup>1</sup>1e-mail:binwang@fma.if.usp.br, Ru-Keng Su<sup>c,</sup><sup>2</sup><sup>2</sup>2e-mail:rksu@fudan.ac.cn and Elcio Abdalla<sup>a,</sup><sup>3</sup><sup>3</sup>3e-mail:eabdalla@fma.if.usp.br
<sup>a</sup> Instituto De Fisica, Universidade De Sao Paulo, C.P.66.318, CEP 05315-970, Sao Paulo, Brazil
<sup>b</sup> Department of Physics, Shanghai Teachersโ University, P. R. China
<sup>c</sup> Department of Physics, Fudan University, Shanghai 200433, P. R. China
## Abstract
The entropy of anti-de Sitter Reissner-Nordstr$`\ddot{o}`$m black hole is found to be stored in the material which gathers to form it and equals to $`A/4`$ regardless of material states. Extending the study to two kinds of extreme black holes, we find different entropy results for the first kind of extreme black hole due to different material states. However for the second kind of extreme black hole the results of entropy are uniform independently of the material states. Relations between these results and the stability of two kinds of extreme black holes have been addressed.
PACS number(s): 04.70.Dy, 97.60.Lf
Traditionally it is widely believed that black holes have a gravitational entropy given by $`S_{BH}=A/4`$, where $`A`$ is its area and units are such that $`c=G=\mathrm{}=k=1`$. However until now the full understanding of the origin of this entropy is still lacking, though some possible explanations have been raised \[1-3\]. Recently, Pretorius, Vollick and Israel have made significant progress on this problem . By examining the reversible contraction of a thin spherical shell down to the Reissner-Nordstr$`\ddot{o}`$m (RN) black hole event horizon, they suggested that $`S_{BH}`$ is the equilibrium thermodynamic entropy that would be stored in the material which gathers to form the black hole, if one imagines all of this material compressed into a thin layer near its gravitational radius. It is of interest to extend their study to other black hole models and investigate whether their operational definition for black hole entropy is valid for other black holes. This is the first motivation of the present paper. In view of recent interest in anti-de Sitter geometries, we will extend the study of ref. to an asymptotically anti-de Sitter version of RN black hole .
The second motivation of this paper is to extend the operational approach to extreme black hole (EBH) entropy. Recently there have been heated discussions on EBH entropies and different results obtained by using different treatments \[6-10\]. Starting with the original RN EBH, Hawking et al claimed that a RN EBH has zero entropy, infinite proper distance $`l`$ between the horizon and any fixed point . However, in the grand canonical ensemble, Zaslavskii argued that EBH can be obtained as the limit of nonextreme counterpart by first adopting the boundary condition $`r_+=r_B`$, where $`r_+`$ is the event horizon and $`r_B`$ is boundary of the cavity, and then the extreme condition. The final extreme hole is in the topological sector of nonextreme configuration and its entropy still obeys the Bekenstein-Hawking formula \[8-10\]. Recently by using these two treatments, the geometry and intrinsic thermodynamics have been investigated in detail for a wide class of EBHs including 4D and two-dimensional (2D) cases \[11-13\]. It was found that these different treatments lead to two different topological objects represented by different Euler characteristics and show drastically different intrinsic thermodynamical properties both classically and quantum-mechanically. Based upon these results it was suggested that there maybe two kinds of EBHs in nature: the first kind suggested by Hawking et al with the extreme topology and zero entropy, which can only be formed by pair creation in the early universe; on the other hand, the second kind, suggested by Zaslavskii, has the topology of the nonextreme sector and the entropy is still described by the Bekenstein-Hawking formula, which can be developed from its nonextreme counterpart through second order phase transition \[11-13\]. This speculation has been further confirmed recently in a Hamiltonian framework and using the grand canonical ensemble as well as canonical ensemble formulation for RN anti-de Sitter black hole by finding that the Bekenstein-Hawking entropy and zero entropy emerge for extreme cases respectively. It is worth extending the operational approach to investigate the EBH entropy, especially for these two kinds of EBHs, and compare the operational definitions of EBH entropy to the results available. In some attempts have been made to study the EBH entropy and some ambiguous results have been obtained, however their study is only limited in the first kind of RN EBH. By extending the study to two kinds of RN EBH as well as anti-de Sitter RN EBH, we will show that ambiguous results for EBH entropy do appear for the first kind of EBH, disappearing for the second kind of EBH. Some physical understanding on this problem will be given.
The RN black hole solution of Einstein equations in free space with a negative cosmological constant $`\mathrm{\Lambda }={\displaystyle \frac{3}{l^2}}`$ is given by
$$\mathrm{d}s^2=h\mathrm{d}t^2+h^1\mathrm{d}r^2+r^2\mathrm{d}\mathrm{\Omega }^2$$
(1)
where
$$h=1\frac{r_+}{r}\frac{r_+^3}{l^2r}\frac{Q^2}{r_+r}+\frac{Q^2}{r^2}+\frac{r^2}{l^2}$$
(2)
The asymptotic form of this spacetime is anti-de Sitter. There is an outer horizon located at $`r=r_+`$. The mass of the black hole is given by
$$m=\frac{1}{2}(r_++\frac{r_+^3}{l^2}+\frac{Q^2}{r_+})$$
(3)
In the extreme case $`r_+,Q`$ satisfy the relation
$$1\frac{Q^2}{r_+^2}+\frac{3r_+^2}{l^2}=0.$$
(4)
Now we consider compressing a spherical shell reversibly from an infinite radius down to the black hole event horizon. We will concentrate our attention on the NEBH case at the beginning with the left-hand-side of Eq.(4) bigger than zero. To maintain reversibility at each stage the shell must be in equilibrium with the acceleration radiation that would be measured by an observer on the shell. To ensure this equilibrium, an โadiabaticโ diaphragm must be interposed between the faces. As done in , we picture the shell as a pair of concentric spherical plates, with inner and outer masses $`M_1`$ and $`M_2`$, separated by a massless and thermally inert interstitial layer of negligible thickness. These two plates separate three concentric spherical regions: an inner region where $`h(r)=h_1(r)`$, a very thin intermediate flat region with $`h(r)=1`$ and an outer region where $`h(r)=h_2(r)`$. The local temperature $`T_i`$ of the plates is given by the expression
$$T_i=\frac{h_i^{}}{4\pi \sqrt{h_i}}(i=1,2)$$
(5)
Introducing the Gaussian normal coordinates near every point on the plate $`\mathrm{\Sigma }`$, the coordinates on $`\mathrm{\Sigma }`$ are $`(\tau ,\theta ,\varphi )`$, where $`\tau `$ is the proper time for an observer on $`\mathrm{\Sigma }`$. Defining $`\stackrel{}{N}`$ as the unit spacelike vector orthogonal to $`\mathrm{\Sigma }`$ and $`\stackrel{}{U}`$ the velocity of a mass element of this surface, the orthogonal condition becomes $`\stackrel{}{N}\stackrel{}{U}=0`$. The velocity is $`\stackrel{}{U}=\dot{t}_t+\dot{r}_r`$ where the overdot denotes differentiation with respect to $`\tau `$. We obtain $`\stackrel{}{N}=(|g_{tt}|)^1\dot{r}_t+|g_{tt}|\dot{t}_r`$. The normalization conditions are $`\stackrel{}{N}\stackrel{}{N}=1,\stackrel{}{U}\stackrel{}{U}=1`$. The extrinsic curvatures relative to the Gaussian normal coordinates are simply $`K_{\tau \tau }=N_{\tau ;\tau }=U^\mu U^\nu N_{\mu \nu }`$ and $`K_{xy}=N_{x;y}(x,y=\theta ,\varphi )`$. Evaluating the jump $`\gamma _j^i`$ in the extrinsic curvatures between different regions and employing Israelโs equation
$$\gamma _j^i\delta _j^i\mathrm{Tr}\gamma _{ij}=8\pi s_j^i,$$
(6)
we can obtain masses $`M_i`$ and surface pressures $`P_i(i=1,2)`$ for the inner and outer plates,
$`M_i`$ $`=`$ $`R\xi _i(1\sqrt{h_i(R)})`$ (7)
$`16\pi P_i`$ $`=`$ $`({\displaystyle \frac{\xi _ih_i^{}(R)}{\sqrt{h_i(R)}}}{\displaystyle \frac{2M_i}{R^2}})`$
where $`R`$ is the common radius of the two plates and $`\xi _i=(1)^i`$.
The shell serves merely as the working substance and the nature of the material in the shell is irrelevant, provided that the first law of thermodynamics,
$$\mathrm{d}S=\beta \mathrm{d}M+\beta P\mathrm{d}A\alpha \mathrm{d}N,$$
(8)
is satisfied (we drop the index $`i`$ for the moment). Above, $`\beta =1/T,\alpha =\mu /T`$ and $`\mu `$ is the chemical potential. $`N`$ here is introduced as the number of particles in the shell to make the representation of differential $`\mathrm{d}S`$ complete . Using the Gibbs-Duhem relation
$$S=\beta (M+PA)\alpha N$$
(9)
we have
$$n\mathrm{d}\alpha =\beta \mathrm{d}P+(\sigma +P)\mathrm{d}\beta $$
(10)
where $`n=N/A`$ and $`\sigma `$ is the surface mass density which satisfies $`\sigma ={\displaystyle \frac{M}{4\pi R^2}}`$. Substituting Eqs(5,7) into (10), we get
$$n\mathrm{d}\alpha =\sigma ^2\mathrm{d}(\frac{2\pi \sqrt{h}}{\sigma h^{}})$$
(11)
The functions $`n`$ and $`\alpha `$ can be chosen arbitrarily if they satisfy (11). For simplicity, we can choose plate materials fulfilling the state equation
$$n_i^{}=\sigma _i^2,\mathrm{and},\alpha _i^{}=\frac{2\pi \sqrt{h_i}}{\sigma _ih_i^{}}.$$
(12)
Considering the chemical potential $`\mu _i^{}=T_i\alpha _i^{}`$, we have
$$\mu _i^{}n_i^{}=\sigma _i/2$$
(13)
Substituting the above into Eq(9), we arrive at the entropy density $`s_i=S_i/A`$ of the plates as
$$s_i^{}=\beta _iP_i+\beta _i\sigma _i\alpha _i^{}n_i^{}=\beta _iP_i+\frac{\beta _i\sigma _i}{2}.$$
(14)
Using Eq(7) for the surface pressure and the local temperature of the plate Eq(5), we find
$$s_2^{}=1/4$$
(15)
For the outer plate, when it reaches the black hole horizon, its entropy is one quarter of its area in Planck units.
This result is also valid choosing general functions $`n,\alpha `$ which satisfy (11). The most general way of proceeding along these lines is
$$\alpha _i=g_i(\alpha _i^{}),n_i=n_i^{}/g_i^{}(\alpha _i^{})$$
(16)
and
$$\frac{\mu _in_i}{\sigma _i}=\frac{g_i}{2\alpha _i^{}g_i^{}}.$$
(17)
The most general expression for the entropy density of the plates is
$$s_i=\frac{1}{4}[\xi _i+\frac{4\pi \sqrt{h_i}}{h_i^{}}(2\sigma _i8\mu _in_i)]$$
(18)
When the outer plate approaches the horizon, $`h_20`$ as $`Rr_+`$. Therefore $`s_2=1/4`$ again. This means that in the anti-de Sitter RN NEBH, regardless of equations of states, the entropy of a shell made of materials approaches $`A/4`$ as the shell approaches its event horizon. This result is in agreement with for RN NEBH.
Now it is of interest to extend the above discussion to EBH cases. As stated in \[11-13\], two kinds of EBHs emerge due to two different treatments. The first kind of EBH obtained by Hawking et al is the original EBH with zero entropy, zero Euler characteristic and arbitrary imaginary time period $`\beta `$ because of its peculiar topology and no conical singularity for its spacetime . While the second kind of EBH proposed by Zaslavskii has entropy equaling $`A/4`$ and the same topology as that of NEBH \[8-10\]. We hope that using operational approach can give a deeper understanding of these two kinds of EBHsโ entropies.
In , attempts have been given to find the RN EBH entropy by using operational approach. However in their study, the authors only consider the first kind of RN EBH case with $`\beta `$ arbitrary and arrive at $`s_{EBH}^{}(1)=0`$ and $`s_{EBH}(1)={\displaystyle \frac{1}{2}}(1{\displaystyle \frac{\mu n}{\sigma }})`$ for different states of materials in the extremally charged spherical shell collapsing onto the hole. They concluded that entropy of RN EBH may depend on their prior history. This result is not valid for the second kind of RN EBH as we will show in the following. This kind of EBH, obtained by first taking the boundary limit and then the extreme limit in the grand canonical ensemble, has the same topology as that of NEBH, and still has a conical singularity in the spacetime. Therefore, $`\beta `$ cannot be arbitrary, being given by $`1/T`$, where $`T`$ is the nonzero local temperature . We keep this fact in mind and let a nonextreme shell collapse to black hole horizon first and make it become extreme afterwards as we have done for obtaining the second kind of EBH. We thus find, for the simplest choice of $`n,\alpha `$ as (22) in ,
$`s_{EBH}^{}(2)`$ $`=`$ $`\beta _2P_2`$
$`=`$ $`{\displaystyle \frac{4\pi V_2}{f_2^{}}}{\displaystyle \frac{1}{16\pi }}({\displaystyle \frac{\xi _2f_2^{}}{V_2}}{\displaystyle \frac{2M_2}{R^2}})=1/4,`$
where we first took the boundary limit $`Rr_+`$, which leads to $`V_2=\sqrt{f_2(R)}0`$.
However when $`\alpha ,n`$ fulfill the general formulas (29,30) of , we still have
$`s_{EBH}(2)`$ $`=`$ $`\beta _2P_2+\beta _2\sigma _2(1{\displaystyle \frac{\mu _2n_2}{\sigma _2}})`$
$`=`$ $`{\displaystyle \frac{4\pi V_2}{f_2^{}}}{\displaystyle \frac{1}{16\pi }}({\displaystyle \frac{\xi _2f_2^{}}{V_2}}{\displaystyle \frac{2M_2}{R^2}})+{\displaystyle \frac{4\pi V_2}{f_2^{}}}\sigma _2(1{\displaystyle \frac{\mu _2n_2}{\sigma _2}})=1/4`$
in the boundary limit, taking $`Rr_+`$ and $`V_20`$.
These results indicate that unlike the results obtained for the first kind of RN EBH, the operational approach leads to universal results for the entropy of the second kind of RN EBH. These results also hold in Anti-de Sitter RN EBH.
For the first kind of Anti-de Sitter RN EBH, because of their peculiar topology, it has been shown that its imaginary time period $`\beta `$ is arbitrary . Taking account of this fact, from Eq(14) and (7), we find for the extreme shell satisfying the equation of state (12)
$$s_{EBH}^{}(1)=\beta \frac{\xi _2}{16\pi }\frac{h_i^{}}{\sqrt{h_i}}=\beta \frac{\xi _2}{16\pi }2(\sqrt{h_i})^{}=0$$
(21)
We took $`Rr_+(h_i0)`$. It is the same as that in the RN first kind of EBH case. However it is worthy to point out that the entropy here decreases as the extreme shell approaches the black hole horizon, unlike the first kind of RN EBH where $`s_{EBH}^{}(1)=0`$ at all stages. This is because in the first kind of RN EBH, the surface pressure for the extremely charged shell material is always zero, while this does not hold for the first kind of anti-de Sitter RN EBH.
For general state satisfying (16,17),
$$s_{EBH}(1)=\frac{\beta }{4}(2\sigma _i8\mu _in_i)$$
(22)
when the extreme shell collapse to the black hole horizon. These results again support the argument that the entropy of the first kind of EBHs may depend on their prior history .
Now we turn to study entropy for the second kind of anti-de Sitter RN EBH. It has been shown that it has the same topology as that of the NEBH, therefore it has a conical singularity with $`\beta =1/T`$ . $`T`$ here is the local temperature $`T=T_H/[h(r_B)]^{1/2}`$ and $`T_H`$ is the Hawking temperature. For the second kind of EBH $`T`$ is nonzero though $`T_H=0`$ . And in the grand canonical ensemble actually only the local temperature $`T`$ has physical meaning, whereas $`T_H`$ can always be rescaled without changing observable quantities . Therefore, this kind of EBH can be achieved with no contradiction with the third law of thermodynamics. Let the nonextreme shell collapse to the black hole and make it extreme afterwards, which corresponds to the treatment of Zaslavskii by first adopting the boundary limit and then the extreme limit. We have
$$s_{EBH}^{}(2)=\{\frac{4\pi \sqrt{h_2}}{h_2^{}}[\frac{\xi _2h_2^{}}{16\pi \sqrt{h_2}}]|_{Rr_+}\}|_{extr}=1/4$$
(23)
for the equation of state of the shell material given by (12).
For the shell material satisfying the general equations of (16,17)
$$s_{EBH}(2)=\{\frac{4\pi \sqrt{h_2}}{h_2^{}}[\frac{\xi _2h_2^{}}{16\pi \sqrt{h_2}}+\frac{\sigma _2}{4}(2\frac{8\mu _2n_2}{\sigma _2})]_{Rr_+}\}|_{extr}=1/4,$$
(24)
by taking $`Rr_+`$ first ($`h_20`$ first). Therefore the entropies of this second kind of EBH are independent of their prior history as that of the second kind of RN EBH case.
At first sight, it seems hard to believe why the operational approach leads to drastically different results for two kinds of EBHs. Especially the changable results for entropy of the first kind of EBH concerning different equations of state of collapsing materials. We know that entropy is a function of state. We have shown in our previous papers \[11-13,18\] that although both of these two kinds of EBHs satisfy the extreme condition, their topological properties differ drastically, so we cannot treat them as the same state. This understanding may help us to understand the different entropies for two kinds of EBHs. But how can we explain different operational definitations of the entropy for the first kind of EBH?
As an example, let us first go over the issue of stability discussed for RN black hole . The heat capacity at constant electrostatic potential difference and cavity radius can be computed from $`C_{\varphi ,r_B}=\beta ({\displaystyle \frac{S}{\beta }})_{\varphi ,r_B}`$, which leads to $`C=4\pi r_B^2x^3(1x)/(3x^22xq^2),`$ where $`x=r_+/r_B,q=e/r_B`$ Eq(5.17) of . For the first kind of EBH obtained by Hawking et alโs treatment (adopting extreme condition at the very beginning) $`C=4\pi r_B^2x^3(1x)/2x(x1)<0`$. The negative sign here determine that this kind of EBH is locally unstable. In addition, it follows from Eq(5.9-5.11) of , that the second derivatives of the action $`I`$ with respect to entropy $`S`$ and the mean charge value $`<Q>`$ diverge by starting with the original EBH \[for example, $`^2I/<Q>^2=4\pi (1x)(1q^2/x^2)^1(1q^2/x)^1,q=x`$ for the original EBH\]. This means that fluctuations of the charge and entropy are infinite for the first kind of EBH. Because of this instability of the first kind of EBH, different states of materials collapses onto it will of course lead to different entropy results, while for the second kind of EBH, the entropy is uniform because of its stability, as shown in . The extension of the stability results to two kinds of anti-de Sitter RN EBHs is obvious.
In summary, we have extended the operational definition of RN black hole entropy to an interesting anti-de Sitter RN NEBH model and found that the entropy for this NEBH is described by the Bekenstein-Hawking formula as well. Extending the operational approach to two kinds of EBHs we suggested \[11-13\], we arrived at different results for their entropy. For the first kind of EBH, the entropy values depends on different material equations of state which collapse onto the hole, which is in agreement with the results of . However, for the second kind of EBH, uniform entropy emerges, regardless of the material equation of state. These different results can be attributed to the issues of stability for these two kinds of EBHs. Using the argument given in , it is easy to find that the first kind of EBH is unstable. Thus its entropy changes in case of different states of the material collapsing onto it. However the second kind of EBH is stable , and it has the same values of entropy, regardless of material states. Therefore, although two kinds of EBHs can be created, due to stability only the second kind of EBH can last for long in nature.
It is worth pointing out that all discussions for obtaining EBH in this paper are just from a theoretical viewpoint and refers to the mathematical treatment only. Physical realization for creating EBHs leads to the problem about how to satisfy the third law of thermodynamics. Although it is clear that quantum processes like evaporation, which typically involve the absorption of negative energy, can violate Nernstโs form of the third law, in all classical studies it holds \[4,21-23\]. The fact that the entropy of the first kind of EBH tends, as $`T0`$, to an absolute zero, ensures that the strong version of the third law holds in this kind of EBH. For the second kind of EBH obtained by first letting the nonextreme shell collapse to the black hole horizon and next making it become extreme does not upset the third law as well, because the final extreme state can be obtained at nonzero temperature. Therefore no challenge to the third law arise. Besides mathematical treatments, physical processes for obtaining EBH still need further studies. This is still an open question and we shall discuss it elsewhere.
ACKNOWLEDGMENT: This work was partically supported by Fundacรฃo de Amparo ร Pesquisa do Estado de Sรฃo Paulo (FAPESP) and Conselho Nacional de Desenvolvimento Cientรญfico e Tecnolรณgico (CNPQ). B. Wang would like to acknowledge the support given by Shanghai Science and Technology Commission.
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# Resonant decay of parity odd bubbles in hot hadronic matter
## 1 Introduction
Recently, Kharzeev, Pisarski and Tytgat presented the following idea, based on the nontrivial topological structure of the QCD-vacuum : Metastable states which act like regions with a non-vanishing QCD vacuum angle $`\theta `$ may be excited when matter undergoes the deconfining phase transition, provided it is of second order. In these false vacua, parity (and also CP) is spontaneously broken, a quality that could lead to experimental signatures like parity odd correlations of the produced particles.
In the present work, we investigate the production of $`\eta ^{}`$-particles during the decay of the CP-odd metastable states, concentrating on the amplification of the low momentum modes by parametric resonance when the zero mode rolls down from the false to the true vacuum. This mechanism plays an important role for particle production in inflationary cosmology and has also been investigated for the formation of Disoriented Chiral Condensates and for Axions . Here we discuss the efficiency of the amplification mechanism for the production of $`\eta `$ from the decaying metastable bubble. As a main result we derive the corresponding momentum spectrum showing non-thermal behaviour. The integrated density of particles produced is estimated to be $`0.7fm^3`$.
## 2 Metastable states in hot hadronic matter
The idea is based on the effective Lagrangian of the Witten-DiVecchia-Veneziano model
$`_{eff}`$ $`=`$ $`{\displaystyle \frac{f_\pi ^2}{4}}[tr(_\mu U_\mu U^+)+tr(MU+MU^+)`$ (1)
$`{\displaystyle \frac{a}{N_c}}(\theta {\displaystyle \frac{i}{2}}tr(\mathrm{ln}U\mathrm{ln}U^+))^2],`$
which describes the low-energy dynamics of the pseudoscalar mesons in the large $`N_c`$-limit of QCD. In the following, we consider the โreal worldโ-vacuum angle $`\theta 0`$. Formation and decay of a non-zero $`\theta `$-vacuum and its signatures are investigated in . The $`N_f\times N_f`$-matrix $`U`$ in Eq. (1) describes the meson fields, and the term containing the mass matrix $`M`$ provides the explicit soft breaking of chiral symmetry due to the quark masses. $`M`$ can be written as $`M_{ij}=\mu _i^2\delta _{ij}`$, where the diagonal elements are chosen as $`\mu _1^2\mu _2^2m_\pi ^2`$ and $`\mu _3^22m_K^2m_\pi ^2`$. With the parametrization $`U=\mathrm{exp}(i\varphi /f_\pi )`$, the matrix $`\varphi `$ representing the singlet and the octet meson fields reads
$$\varphi =\sqrt{\frac{2}{3}}\eta _1\mathrm{๐}+\left(\begin{array}{ccc}\pi ^0+\frac{\eta _8}{\sqrt{3}}& \sqrt{2}\pi ^+& \sqrt{2}K^+\\ \sqrt{2}\pi ^{}& \pi ^0+\frac{\eta _8}{\sqrt{3}}& \sqrt{2}K^0\\ \sqrt{2}K^{}& \sqrt{2}\overline{K^0}& \frac{2}{\sqrt{3}}\eta _8\end{array}\right).$$
(2)
The last term in the effective Lagrangian reflects the $`U_A(1)`$-anomaly by giving a mass to the singlet, even in the chiral limit of vanishing quark masses. The parameter $`a=2N_f\lambda _{YM}/f_\pi ^2`$ represents the topological susceptibility.
Finite temperature is introduced into the model via the temperature dependence of the parameters: For $`T=0`$, the value $`a=m_\eta ^2+m_\eta ^{}^22m_K^20.726GeV^2`$ can be determined from experimental results , and $`\mu ^2(\mu _1^2+\mu _2^2+\mu _3^2)/3=(m_\pi ^2+2m_K^2)/30.171GeV^2`$. The corresponding value of the pion decay constant is $`f_\pi 93MeV`$.
When the temperature approaches the temperature $`T_d`$ of the phase transition, $`a(T)`$ goes to zero, indicating the effective restoration of the $`U_A(1)`$-symmetry for $`T>T_d`$ with enhanced $`\eta `$ \- and $`\eta ^{}`$ \- production as a possible signature . We use the temperature of the deconfinement phase transition, $`T_d`$, following the assumption of that any other phase transition of this model also occurs at $`T_d`$. According to mean field estimates , $`a(T)(T_dT)`$ and $`\mu ^2(T)(T_dT)^{1/2}`$ near the phase transition.
In order to get an analytically tractable approximation that allows some insight into the physical processes, we consider only the singlet, which is the main component of the $`\eta ^{}`$-meson. Using $`N_c=3`$, the effective Lagrangian simplifies to
$$_s=\frac{1}{2}(_\mu \eta _1)(^\mu \eta _1)+\frac{3}{2}f_\pi ^2\mu ^2\mathrm{cos}\left(\sqrt{\frac{2}{3}}\frac{\eta _1}{f_\pi }\right)\frac{a}{2}\eta _1^2.$$
(3)
The singlet effective potential
$$V\left(\frac{\eta }{f}\right)\frac{V_s}{f^2\mu ^2}\left(\frac{\eta }{f}\right)=\mathrm{cos}\frac{\eta }{f}+\frac{a}{2\mu ^2}\frac{\eta ^2}{f^2},$$
(4)
where we have dropped the index $`1`$ in the notation of the singlet and introduced $`f\sqrt{3/2}f_\pi `$ for convenience, can be considered as a projection of the full effective potential. It takes three qualitatively different shapes, depending on the temperature, more specifically, on the value of $`a(T)/\mu ^2(T)`$. For large values of $`a/\mu ^2`$, corresponding to small $`T`$, the potential is a more or less deformed parabola (see Fig. 1a). For $`a/\mu ^2=0`$, i.e. $`T>T_d`$, the potential is periodic with minima at $`\eta /f=0,2\pi ,4\pi \mathrm{}`$, which are all equivalent true vacua (see Fig. 1d). And for a small, nonzero value of $`a/\mu ^2`$, corresponding to a temperature just below the phase transition, metastable states appear in the potential which are distinct from the true vacuum (see Fig. 1c). One recognizes that the false vacua are odd under parity transformation, but even under charge conjugation, and therefore odd under CP .
One can show that the effective potential exhibits metastable states also in the case of nonvanishing $`\pi ,\eta _8\text{and}K`$fields, provided that $`m_\pi 0`$ . But only the singlet is responsible for the formation of these states, in the sense that it is the only field which appears in the term of the effective potential which includes the topological susceptibility.
## 3 The decay of the metastable states
### 3.1 Dynamics of the model
Applying the model to the description of a heavy ion collision, we suppose that during the phase transition, the potential changes its shape from the parabola (Fig. 1a), corresponding to the low temperature phase, to that of the high temperature phase (Fig. 1d), and back to the low temperature phase, but this time including metastable states that have been created at a temperature just below the phase transition. After the phase transition, but still at high temperature, the $`\eta ^{}`$field may be trapped in a metastable state, forming a โCP-odd bubble inside the hadronic phaseโ in the language of . Considering the effective potential as a function of temperature, one recognizes that there exists a value $`a/\mu ^2(T_{sp})`$ for which the local minimum turns into a saddle point (see Fig. 1b). For the lowest false vacuum, this happens for $`a/\mu ^2(T_{sp})=0.217`$ at the value $`(\eta /f)_{sp}=4.493`$. Assuming the proportionality $`a(T)/\mu ^2(T)(T_dT)^{3/2}`$ to hold even far from $`T_d`$, the temperature corresponding to this special case of the potential can be estimated to be $`T_{sp}0.86T_d`$. Once the saddle point temperature is reached and the barrier has disappeared, the $`\eta ^{}`$field starts rolling down into the true vacuum and oscillating around it, while energy is transferred from the mean value of the field into its quantum fluctuations. The resulting growth of the occupation numbers of the quantum fluctuations is interpreted as particle production with respect to the true, CP even, vacuum, whereas the zero mode is damped. This process is investigated in more detail in the following.
A crucial assumption for this picture is that at least after the saddle point temperature has been reached, the effective potential changes slowly compared to the motion of the field, so that the oscillations can be considered as taking place in a static potential, frozen at $`T=T_{sp}`$. Accordingly, the ratio $`a/\mu ^2=0.217`$ is kept fixed during the evolution in time, and the value of $`f`$ as well. The solutions we find (see below) show that the zero mode reaches the true vacuum for the first time after $`12fm`$, whereas the relevant time scale to be considered for the particle production is of the order of $`10fm`$.
### 3.2 Evolution equations
To study the dynamics of the decay process, we decompose the field into its expectation value $`\phi (t)\eta (\stackrel{}{x},t)`$ and fluctuations $`\chi `$ about it,
$$\eta (\stackrel{}{x},t)=\phi (t)+\chi (\stackrel{}{x},t)$$
(5)
with $`\chi (\stackrel{}{x},t)=0`$, so that the โphysicalโ fluctuations are given by the usual expression $`(\chi \chi )^2=\chi ^2`$. With this decomposition, following the standard method as described in, e.g., , we obtain two evolution equations from the effective Lagrangian (3). The effect of the quantum fluctuations at lowest order is taken into account with a Hartree-type approximation that consists of the factorization $`\chi ^33\chi ^2\chi `$ and $`\chi ^2\chi ^2`$ and the self-consistency condition represented by Eq. (7) given below; terms of higher order in $`\chi `$ are neglected. Taking the expectation value $``$, we arrive at the equation for the zero mode
$$\frac{d^2}{d\tau ^2}\frac{\phi (\tau )}{f}+\left(1\frac{\chi ^2_\tau }{2f^2}\right)\mathrm{sin}\frac{\phi (\tau )}{f}+\frac{a}{\mu ^2}\frac{\phi (\tau )}{f}=0,$$
(6)
written for the dimensionless mean field $`\phi /f`$, where the dimensionless time variable $`\tau \mu t`$ is introduced. The expectation value of the fluctuations
$$\chi ^2_\tau \frac{d^3k}{(2\pi )^3}|\chi _k(\tau )|^2$$
(7)
is calculated from the mode functions $`\chi _k(\tau )`$ which are known from the Fourier expansion of the field. The integration is done over the interval of momenta for which the quantum fluctuations are enhanced.
We use the definition (7) for the back reaction, although it does not fulfill $`\chi ^2_0=0`$, because it reflects the background of the calculation, the closed-time-path technique which is usually applied for the description of out-of-equilibrium dynamics in field theory. In this framework, the fluctuations are defined as the two-point correlation functions or equal-time Greenโs functions
$$\chi ^2_\tau \frac{d^3k}{(2\pi )^3}G_k(\tau ,\tau ).$$
(8)
The mode functions satisfy the mode equation
$$\frac{d^2}{d\tau ^2}\chi _\kappa (\tau )+\omega _\kappa ^2(\tau )\chi _\kappa (\tau )=0,$$
(9)
written for the dimensionless mode functions $`\chi _\kappa \sqrt{\mu }\chi _k`$, where
$$\omega _\kappa ^2(\tau )=\kappa ^2+\left(1\frac{\chi ^2_\tau }{2f^2}\right)\mathrm{cos}\frac{\phi (\tau )}{f}+\frac{a}{\mu ^2}$$
(10)
is the time-dependent frequency squared with $`\kappa k/\mu `$. They are used for the calculation of the particle production: The induced particle number density
$$n(\tau )=\frac{d^3k}{(2\pi )^3}n_k(\tau )$$
(11)
is obtained from the spectral particle number density
$$n_\kappa (\tau )=\frac{\omega _\kappa }{2}\left(|\chi _\kappa (\tau )|^2+\frac{|\dot{\chi }_\kappa (\tau )|^2}{\omega _\kappa ^2}\right)\frac{1}{2},$$
(12)
which can be calculated from the mode functions. Since the unstable solutions of the mode equations contribute to the resonant particle production, the integration in (11) is carried out over the momenta corresponding to these amplified solutions, as in Eq. (7). For practical purposes, we use $`\overline{\omega }_\kappa ^2\kappa ^2+a/\mu ^2`$ in definition (12), because it keeps the particle number positive definite for all values of $`\tau `$ and minimizes the fluctuations in the particle number. Accordingly, the initial values of the mode functions are chosen as $`\chi _\kappa (0)=1/\sqrt{2\overline{\omega }_\kappa }`$ and $`\dot{\chi }_\kappa (0)=i\sqrt{\overline{\omega }_\kappa /2}`$ such that $`n_\kappa (0)=0`$. However, the results depend only weakly on the choice of $`\omega _\kappa `$ both in the initial conditions and in the definition of the particle number. The choice for the initial conditions of the zero mode is suggested by the model: $`\varphi `$ is โreleasedโ right below the saddle point $`\phi _{sp}/f=4.493`$ with zero velocity, i.e. $`\dot{\phi }(0)=0`$. For the numerical results shown below, we used the starting value $`\phi (0)/f=4.2`$.
### 3.3 Solutions and results
Considering the coupled system (6) and (9) with (10), one recognizes that the expectation value transfers energy to the mode functions via a time dependent frequency, and the mode functions in turn modify the equation of motion for the zero mode. A natural first step in the solution of this system is to neglect the fluctuation terms in both equations, obtaining the โclassicalโ evolution equations with $`\chi ^2_\tau =0`$ where the mode functions do not react back on the zero mode. If the variable frequency
$$\omega _{\kappa ,cl}^2=\kappa ^2+\mathrm{cos}\frac{\phi (\tau )}{f}+\frac{a}{\mu ^2}$$
(13)
in the mode equation depends periodically on time, according to Floquetโs theorem it has quasiperiodic solutions of the form
$$\chi _\kappa ^{cl}(\tau )=\mathrm{exp}(\mu _\kappa \tau )P(\tau ),$$
(14)
where the characteristic exponent $`\mu _\kappa `$ depends on $`\omega _{\kappa ,cl}`$, and $`P(\tau )`$ is a periodic function of the same period as $`\omega _{\kappa ,cl}`$, usually normalized to have unit amplitude. For certain values of the parameter $`\omega _{\kappa ,cl}`$, there exist so called resonance bands that contain the exponentially growing (or resonant) solutions with a real characteristic exponent $`\mu _\kappa `$. These resonance bands are known, for example, for the solutions of the Mathieu- or the Lamรฉ-equation.
An analytic solution of the classical system can be found with a Sine-Gordon-equation , which consists essentially of setting $`a/\mu ^20`$ in our case. The corresponding approximate zero mode equation is given by
$$\frac{\ddot{\phi }_{cl}(\tau )}{f}+\rho \mathrm{sin}\frac{\phi _{cl}(\tau )}{\rho f}=0,$$
(15)
where $`\rho =1.43`$ represents the radius of the potential and is needed to adjust it to the original potential. It has the solution
$$\frac{\phi _{cl}(\tau )}{f}=2\rho \mathrm{arctan}\left(\sqrt{\frac{ฯต}{2ฯต}}cn(\tau ,m)\right)$$
(16)
with $`ฯต1\mathrm{cos}(\phi (0)/\rho f)`$, and $`m=ฯต/2`$ is the square of the modulus of the Jacobian elliptic function $`cn`$. The resulting mode equation is a Lamรฉ-equation with the resonance band $`0\kappa ^2ฯต/2`$. The resulting $`\mu _\kappa `$ is shown in Fig. 2 in comparison with the result of the numerical calculation of $`\mu _\kappa `$ for the solution of the classical mode equation. In the range of the single resonance band of the Lamรฉ-equation, $`0\kappa ^20.97`$, the numerical calculation yields three resonance bands, and at least two for $`\kappa ^2>1`$. (This structure reminds of the Mathieu-equation which has more than one resonance band.) The maximum value $`\mu _{\kappa ,max}0.42`$ is found around $`\kappa ^20.1`$, and the maximum value for the Sine-Gordon-solution, $`\mu _{\kappa ,max}^{sg}0.5`$, yields an acceptable approximation .
The momentum spectrum of the produced particles reflects the structure of the resonance bands. When the general expression for the mode functions (14) is inserted into Eq. (12), it is obvious that the particle production is essentially determined by $`\mu _\kappa `$:
$$\mathrm{ln}n_\kappa (\tau )\mathrm{ln}\frac{\omega _\kappa (\tau )}{2}2\mu _\kappa \tau .$$
(17)
In the classical approximation discussed above, the back reaction of the fluctuations on the zero mode (and also on the fluctuations themselves) is ignored. This implies that the approximation is not energy conserving, which can immediately be seen from the zero mode solution which is not damped although it is supposed to transfer its energy to the mode functions, which indeed grow exponentially. Since the mode functions are not decoupled from the zero mode, at least the other direction of energy transfer is reflected by the approximation: The larger the initial value $`\phi (0)`$, the more energy can be transferred, and the larger the mode functions and the particle number grow.
The evolution equations including the back reaction are solved numerically. For the evaluation, we use $`f(T_{sp})f(0)`$ and $`\mu (T_{sp})676MeV`$. Because the value of $`\mu `$ in our approximation is only fixed by the value of $`T_{sp}`$, which depends on our restriction to the singlet case and may not correspond to the real physical situation, we also try $`\mu (T=0)=412MeV`$ as the other limit of the range of possible values for $`\mu (T)`$.
Although the inclusion of the back reaction destroys the parametric resonance, it does not suppress the particle production completely. The particle number (Fig. 3) reaches its asymptotic value for $`t_{end}10fm`$. At about the same time, the value of the back reaction term $`\chi ^2_\tau /2f^2`$ grows larger than 1, indicating that the approximations made become inappropriate for the description of the dynamics for larger times (cf. Eq. (6)). One recognizes from Fig. 3 that the choice of $`\mu (T)`$ has some influence on $`t_{end}`$; the asymptotic number density of the produced particles lies between $`0.7fm^3`$ and $`0.9fm^3`$. The absolute number of particles can be calculated from this density if the size of the CP-odd bubbles is known. For example, for a domain with radius $`5fm`$, as suggested in , we estimate a yield of 90-100 particles from our model.
A measure for the size of correlated domains of particles is the correlation length $`\xi `$, which can be calculated from the two-point correlation function
$$D(x,\tau )\frac{d\kappa \kappa ^2}{2\pi ^2}j_0(\kappa x)|\chi _\kappa (\tau )|^2,$$
(18)
where $`x=\mu r`$ denotes the dimensionless length. For large $`x`$,
$$D(x,\tau )\mathrm{exp}\left(\frac{x^2}{\xi ^2(\tau )}\right),$$
(19)
from which we estimate the correlation length for $`t_{end}`$. As a lower limit for the domain radius, we obtain $`\xi 2.4fm`$ for $`\mu (T=0)`$ and $`\xi 1.3fm`$ for $`\mu (T_{sp})`$.
The momentum spectrum of the produced particles at $`t_{end}10fm`$ is shown in Fig. 4. Its maximum is located at about the value of $`\kappa `$ as the maximum of $`\mu _\kappa `$ in Fig. 2, which had been calculated for the system without back reaction. The spectrum deviates significantly from a thermal Bose-Einstein distribution.
R.B. thanks the E. Schrรถdinger Institute, Vienna, for support during the workshop โQuantization, Generalized BRS Cohomologies and Anomaliesโ, where this work has been encouraged from discussions with M.H.G. Tytgat. We thank especially D. Kharzeev for useful remarks and suggestions. This research is supported in part by DFG, contract Ka 1198/4-1.
## Figure Captions
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# limit-from๐ฟ-invariant for Quasi-periodic Oscillations and Physical Parameters of 4U 0614+09 binary
## 1 Introduction
This Letter contains further verification of the Two Oscillator (TO) model for kHz QPOโs in NS binaries and the related low frequency features using the observational results for 4U 0614+09 from van Straaten et al. (2000) (hereafter VS00). In a series of the papers (Osherovich & Titarchuk 1999a, Titarchuk & Osherovich 1999, Osherovich & Titarchuk 1999b, Titarchuk, Osherovich & Kuznetsov 1999, hereafter OT99a, TO99, OT99b and TOK respectively) the authors offered the TO model which explained the variation of the frequency difference between two kHz QPOs suggesting that the upper one, namely $`\nu _h`$ is an upper hybrid frequency of the Keplerian oscillator under the influence of the Coriolis force and the lower kHz QPO is the Keplerian frequency $`\nu _K`$. This identification is a principal difference between the TO model and other QPO models presented so far (e.g. Stella & Vietri 1998 Miller, Lamb, Cook 1998, Kaaret, Ford & Chen 1997) since in other models, the highest kHz QPO is identified as a Keplerian frequency.
According to the TO model, the Keplerian oscillator has two branches characterized by high frequency $`\nu _h`$ ($``$ 1 kHz) and by low frequency $`\nu _L`$ ($`50100`$ Hz). The frequency $`\nu _L`$ depends strongly on the angle $`\delta `$ between the normal to the neutron star disk and $`๐`$the angular velocity of the magnetosphere surrounding the neutron star. In the lower part of the QPO spectrum ($``$ 10 Hz), the second oscillator of the TO model describes the physics of the viscous transition layer, namely, radial viscous oscillations with frequency $`\nu _V`$ (previously called โextra noise componentโ) and the diffusive process in the transition region (the innermost part of the disk) which is characterized by the break frequency $`\nu _b`$. According to the TO model, all frequencies (namely $`\nu _h`$, $`\nu _L`$, $`\nu _b`$ and $`\nu _V`$) have specific dependences on $`\nu _\mathrm{K}`$. Correlations of $`\nu _b`$ and $`\nu _V`$ with kHz frequencies are fitted in TO99 by the theoretical curves using the dimensionless parameter, $`a_K`$:
$$a_\mathrm{K}=m(x_0/3)^{3/2}(\nu _0/363\mathrm{Hz})$$
(1)
where normalized mass $`m=M/M_{}`$, radius of the inner edge of disk $`x_0=R_0/R_\mathrm{S}`$ in Schwarzchild radius ($`R_\mathrm{S}`$ units) and $`\nu _0`$ is the spin of the disk inner edge.
Van Straaten et al. (2000) have found for 4U 0614+09 that dependences of $`\nu _V`$ and $`\nu _b`$ on $`\nu _\mathrm{K}`$ do not follow the theoretical curves which fit the data points for 4U1728-34 (TO99) and Sco X-1 (TOK). In this Letter, we revisit this issue and show that the observational correlations for source 4U 0614+09 are related to the theoretical ones for the same parameter $`a_\mathrm{K}=1.03`$ (as for 4U 1728-34 and Sco X-1) but with a spin of the inner edge of accretion disk higher than 363 Hz. Such a high spin leads to the high compactness of the central object because
$$R_0=9\mathrm{km}\times m^{1/3}\left(\frac{a_\mathrm{K}}{\nu _0/363\mathrm{Hz}}\right)^{2/3}.$$
(2)
We check predictions of the TO model regarding $`\nu _V`$ vs $`\nu _\mathrm{K}`$ dependence as well as relation between $`\nu _b`$ and $`\nu _\mathrm{K}`$ for all available observations (see ยง3).
In OT99b, the authors demonstrated for the source 4U 1702-42 that the inferred angle $`\delta `$, (see Eq 2, 5 there)
$$\delta =\mathrm{arcsin}\left[(\nu _h^2\nu _\mathrm{K}^2)^{1/2}(\nu _L\nu _h/\nu _\mathrm{K})\right]$$
(3)
depends only, on $`\nu _\mathrm{K}`$, $`\nu _h`$ and $`\nu _L`$ and stays the same ($`3.9^o\pm 0.2^o`$) over significant range of $`\nu _\mathrm{K}`$ (650-900 Hz). In Formula (3), $`\delta `$ is assumed to be small. A general formulation for any $`\delta `$ is presented, in OT99a with conjecture that eventually $`\delta `$ as an invariant will be found for all 20 sources known to have QPOs. Using series of observations where $`\nu _h`$, $`\nu _\mathrm{K}`$, $`\nu _L`$ are detected simultaneously (van der Klis et al. 1997; Markwardt, Strohmayer & Swank 1999; van Straaten et al. 2000) we check that the angle $`\delta `$ is a true invariant for three specific sources: Sco X-1, 4U 1702-42, 4U 0614+09 (see ยง2). We will also present the magnetospheric rotational profile for 4U 0614+09 inferred from VS00 data in ยง4. Discussion and summary follow in the last section.
2. $`\delta `$Invariant and Verification of Two-Oscillator Model
We have used the frequencies measured for sources 4U 0614+09, Sco X-1 and 4U 1702-42 where the low and high kHz QPO peaks $`\nu _\mathrm{K}`$ and $`\nu _h`$ and HBO frequencies $`\nu _L`$ are measured simultaneously. The resulting values of $`\delta `$ calculated from Eq. (3) are shown in Figure 1. Indeed, for each of these sources, the $`\delta `$values show little variation with $`\nu _\mathrm{K}`$, $`\nu _h`$, $`\nu _L`$. They are $`3^o.9\pm 0.2`$ and $`5^o.5\pm 0.5`$ for 4U 1702-42 and Sco X-1 respectively. For the source 4U 0614+09 the point at $`\nu _\mathrm{K}=821.6`$ Hz is not included since the $`\nu _\mathrm{K}`$ peak is very broad with FWHM= 205 Hz, which is almost order of magnitude larger than that for any other points for the same star according to Table 1 in VS00. The angle $`\delta `$ obtained for the source 4U 0614+09 is $`15^o.6\pm 0^o.5`$ and is significantly larger than for other stars. The existence of the $`\delta `$invariant predicted by the TO model (OT99b) is a challenge for any other QPO model. It is important to note that all frequencies included in this $`\delta `$relation are observed frequencies and thus the relation Eq. (3) is a model independent invariant.
3. Break Frequency vs kHz QPO Frequency Correlations
Further tests of the TO model can be done using a comparison of the observed correlation of break and kHz QPO frequencies (VS00) with the theoretical dependences derived in Titarchuk, Lapidus & Muslimov (1998, hereafter TLM) and TO99. In Fig. (2) (upper panel), we present the theoretical curves calculated using Eq. (9) in TO99 for different values of $`\nu _{0,363}=\nu _0/363`$ Hz. For the source 4U 0614+09 we found the Lebesgueโs measure (for definition, see TOK) is better than 15% for $`a_\mathrm{K}=1.03`$ and $`\nu _{0,363}>2`$ which is at least 3 times less than that for $`\nu _{0,363}=1`$ (see upper panel in Fig. 2). We fitted the data points by the curves with $`\nu _{0,363}>3`$ ($`\nu _0>1089`$ Hz) and we have found that the quality of those fits is improved only slightly. Thus it is difficult to get the constraints on $`\nu _0`$ better than $`\nu _0>800`$ Hz. Recently Bhattacharyya et al. (2000) using the full General Relativity technique argued that the NS in Cyg X-2 rapidly rotates with the frequencies of order of 1 kHz which is very close to that we obtained here for 4U 0614+09.
We approximate the solution of Eq (9) of TO99 by a polynomial for $`a_\mathrm{K}=1.03`$ and $`\nu _{0,363}=1`$, namely,
$$\nu _b=C_bP_4(\nu _\mathrm{K})=C_b(B_1\nu _\mathrm{K}+B_2\nu _\mathrm{K}^2+B_3\nu _\mathrm{K}^3+B_4\nu _\mathrm{K}^4)$$
(4)
where $`B_1=1.68\times 10^4`$, $`B_2=1.36\times 10^5`$ Hz<sup>-1</sup>, $`B_3=2.47\times 10^8`$ Hz<sup>-2</sup>, $`B_4=4.45\times 10^{11}`$ Hz<sup>-3</sup>. We fit the observational break frequencies by the theoretical curve which is a function of parameter $`\nu _{0,363}`$. For a given $`\nu _{0,363}`$ that curve can be obtained by the scaling the argument $`x`$ of the polynomial $`P_4(x)`$:
$$\nu _b(\nu _\mathrm{K})=C_bP_4(\nu _\mathrm{K}/\nu _{0,363}).$$
(5)
The normalization of the curve is controlled by the constant $`C_b`$ which reflects the properties of the specific source. For 4U 1728-34 the observed frequencies are fitted by the curve with $`a_\mathrm{K}=1.03`$ (TO99) for which the constant $`C_b=1`$ with $`R_0=11`$ km for 1.4 solar masses. For 4U 0614+09 we fit the data points by the curves with the parameters $`\nu _{0,363}=1,2,3,`$ (i.e. $`\nu _0=`$ 363 Hz, 726 Hz and 1089 Hz) for which the corresponding constants $`C_b=2,22.8,55.4`$ are chosen respectively. Using Eq (2), one can estimate mass-radius (MR) relations for a given value of $`\nu _{0,363}`$. For $`\nu _{0,363}=`$ 2 and 3 the radius $`R_06.4`$ and 4.9 km respectively, for the standard NS mass, 1.4 solar masses. There is definitely an indication of higher compactness for 4U 0614+09. It is true that there is a possibility to elaborate a much more stringent condition on the MR relation using Eq. (2) and fits of the data in Fig. 2 which may lead to small masses and radii for the star $``$ $`m<1`$ and $`R<9`$ km. In that case such a low mass compact star can be easily explained in terms of strange star models (e.g. Bombaci 1997). In fact, under constraints imposed by the TO model Li et al. 1999 pointed out the possibility that 4U 1728-34 is a strange star rather than a neutron star. Thus our arguments of the high compactness in 4U 0614+09 support a strange star concept. But the final conclusion on MR constraints cannot be worked out without the full relativistic treatment which is out of the scope of this paper.
According to the low panel in Fig. 2, $`\nu _h`$ is not lower than 400 Hz for any value of $`\nu _\mathrm{K}`$. Because
$$\nu _h=[\nu _\mathrm{K}^2+4(\mathrm{\Omega }/2\pi )^2]^{1/2},$$
(6)
the rotational magnetospheric frequencies $`\nu _{mag}=(\mathrm{\Omega }/2\pi )`$ should have an asymptote for large radii (or small values $`\nu _\mathrm{K}<2\nu _{mag}`$) which equals approximately 200 Hz. The detailed profile of $`\nu _{mag}`$ constructed below confirms this expectation.
4. Inferred Rotational Frequency Profile of the NS Magnetosphere
From the observed kHz frequencies (VS00, Table 1), $`\nu _h`$ and $`\nu _\mathrm{K}`$, the profile of $`\nu _{mag}=\mathrm{\Omega }(\nu _K)/2\pi `$ has been calculated according to formula (6) and modeled using theoretically inferred magnetic multipole structure of differentially rotating magnetosphere (OT99a)
$$\nu _{mag}=C_0+C_1\nu _K^{4/3}+C_2\nu _K^{8/3}+C_3\nu _K^4$$
(7)
where $`C_2=2(C_1C_3)^{1/2}`$. The constants $`C_0\nu _{mag}^0=190`$ Hz, $`C_1=6.52\times 10^2`$ Hz<sup>-1/3</sup>, $`C_2=1.19\times 10^5`$ Hz<sup>-5/3</sup> and $`C_3=0.91\times 10^9`$ Hz<sup>-3</sup> have been obtained by the least-squares fit with $`\chi ^2=29.8/25`$ (OT99a). The resulting fit is shown in the lower panel of Fig. 3 by solid line. The $`\nu _{mag}`$ profiles for Sco X-1 (which is very similar to those for 4U 1608-52 and 4U 1728-34) and that for 4U 0614+09 differ in the sign of curvature: it is negative for the profile of Sco X-1 profile and positive for that of 4U 0614+09. The $`\nu _{mag}`$ profiles for Sco X-1, 4U 1608-52 (as well as for 4U 1728-34) according to analytical formula (14) in OT99a have a maximum at large $`\nu _\mathrm{K}`$ (small radii) and also have a minimum at small $`\nu _\mathrm{K}`$ (large radii). The adjustment of the Keplerian disk to the sub-Keplerian rotation of the NS in 4U 0614+09 occurs at much larger distances from the central source (3.6-7 NS radii) than that for 4U 1728-34 where the adjustment radii are within 1.6-3.4 NS radii. Thus the blobs (formed at the adjustment radius as result of the super$``$ Keplerian motion \[TLM\]) presumably probe the different parts of the magnetosphere: in the 4U 0614+09 case, it is the outer part of the magnetosphere, near the minimum of the rotational frequency while in the 4U 1728-34 case, it is a part which is closer to the inner part, near the maximum. Thus one should make an expansion of $`\nu _{mag}`$ over the Chandrasekhar potential $`A`$ near the minimum in the former case and near the maximum in the latter case (see OT99a for details). We may try to predict $`\nu _{mag}`$ for much smaller $`\nu _\mathrm{K}`$ (around 150 Hz) than those frequencies $`\nu _\mathrm{K}`$ (between 400 and 850 Hz) where we have constructed the profile (7). Using VS00 data points for 8 August with $`\nu _\mathrm{K}=153.4\pm 5.6`$ Hz and $`\nu _h=449.4\pm 19.5`$ Hz and Eq. (6), we find $`\nu _{mag}=211.\pm 23.`$ Hz which is in good agreement with the value $`\nu _{mag}=235`$ Hz obtained from the analytical expression, Eq. (7).
As shown in Fig. 3 (the upper panel) the difference in $`\nu _{mag}`$ profiles for Sco X-1 and for 4U 0614+09 results in different profiles for $`\mathrm{\Delta }\nu =\nu _h\nu _\mathrm{K}`$: namely, $`\mathrm{\Delta }\nu `$ monotonically decreases with $`\nu _\mathrm{K}`$ for Sco X-1 but has a minimum for 4U 0614+09. According to VS00 $`\mathrm{\Delta }\nu =311.8\pm 1.8`$ Hz except one point (not included in Fig. 3) which deviates 11 $`\sigma `$ from this trend. We will return to this point in another paper.
5. Break Frequency vs Low Lorentzian QPO Frequency Correlation
According to solutions presented in TO99 and TOK, the break frequency $`\nu _b`$ correlates with the low Lorentzian frequency $`\nu _V`$ (which we interpreted as a viscous frequency in the TO model) through power law relation
$$\nu _b0.041\nu _V^{1.61}.$$
(8)
It was confirmed in TOK using the data of Wijnands & van der Klis (1999) that the same index approximately is valid for the wide class of NS and BH objects. In Figure 4 we demonstrate that the observed correlation of $`\nu _b`$ vs $`\nu _V`$ in 4U 0614+09 is described well by the same power law dependence and the quality of this fit is similar to that of 4U 1728-34, i.e. the Lebesgueโs measure is less than 20%. Equations (5) and (8) give us theoretical dependences of $`\nu _V`$ on $`\nu _\mathrm{K}`$, assuming $`\nu _{0,363}=3`$. For August 8, 1996 the theoretical value $`\nu _V9`$ Hz with inherited 20% accuracy comes close to the observed $`\nu _V=11.2\pm 1.4`$ Hz as given in VS00.
6. Discussion and Conclusions
In this Letter we have verified predictions of the two-oscillator (TO) model for three sources. The unifying characteristic of spectra for both oscillators which share the common boundary at the outer edge of the viscous transition layer is the strong dependence of all frequencies in the model on $`\nu _\mathrm{K}`$. The angle $`\delta `$ as a global parameter describes the inclination of the magnetospheric equator to the equatorial plane of the disk. Measured locally for different radial distances (therefore different $`\nu _K`$), $`\delta `$ may vary considerably if the observed oscillations do not correspond to the predicted low Keplerian branch. The constancy of $`\delta `$ shown in Figure 1 argues in favor of the TO model. The parameter $`\delta `$ is critical in evaluating the differences between the spectra of different sources and we suggest that the higher $`\nu _L`$ for 4U 0614+09 and 4U 1728-34 (than that for Sco X-1) are mainly a result of a larger angles $`\delta `$.
We consider the angle $`\delta `$ to be a fundamental geometric parameter for neutron stars and may be related to the evolution and prehistory of these systems. During the NS life time, the NS is believed to spin up gaining more angular momentum from the companion and the axis of magnetosphere may change with respect the rotational axis of the star (or the disk) (Bhattacharya 1995). The X-ray bursters are found in old population (age of order $`10^9`$ years) concentrated towards the galactic center (Lewin, van Paradijs & Taam 1995) and thus it is not surprising that in these systems the angle $`\delta `$ between spin axis and the magnetic axis is small presumably as a result of long evolution. The absence of X-ray pulsations from the LMXBs and the fact that type I X-ray bursters never occur in systems which show pulsations (Lewin, van Paradijs & Taam 1995) probably can be explained by the small angles $`\delta `$. In contrast, in X-ray pulsars which are mostly observed in the High Mass Binaries (HMXBs), $`\delta `$ is not small: the mean value of $`\delta `$ is $`27^0`$ (Leahy 1991) and the age of these systems is much smaller, of order of $`10^7`$ years.
The authors acknowledge discussions with Joe Fainberg, Sergey Kuznetsov and Craig Markwardt. We are grateful to Steve van Straaten for the data which enabled us to make comparisons with the data. We are also grateful the referee for the fruitful suggestions.
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# ROSAT/ASCA Observations of a Serendipitous BL Lac Object PKS 2316-423: The Variable High-Energy Tail of Synchrotron Radiation
## 1. Introduction
BL Lac objects โ a subclass of active galactic nuclei (AGN) โ show variable non-thermal emissions from radio to UV/X-rays, and even to $`\gamma `$-rays over different timescales (e.g., Urry and Padovani 1995; Kollgaard 1994), which are commonly attributed to be synchrotron and inverse Compton radiation from plasma in a relativistic jet oriented at a small angle with our line of sight. As such, they represent a fortuitous natural laboratory to study the physics of jets, and ultimately the mechanisms of energy extraction from the central black holes, a fundamental goal of extragalactic astrophysics.
Earlier studies of BL Lac objects have shown that the systematic differences between radio and X-ray selected BL Lac objects (RBLs vs XBLs) are consistent with orientation differences (Kollgaard et al. 1996; Ghisellini et al. 1993). Meanwhile, BL Lac objects have been reclassified in a more physical way as โlow energyโ and โhigh energyโ peaked BL Lac objects (LBLs vs HBLs) based on the peak frequency of synchrotron radiation (e.g. Giommi and Padovani 1994), rather than just obervational selection. In general, RBLs and XBLs tend to be LBLs and HBLs, respectively. However, recent evidence has shown that some of the differences between HBLs and LBLs cannot be accounted for by differences in orientation alone (Sambruna, Maraschi & Urry 1996). The alternate explanation, that these classes are merely opposite extremes of the sourceโs spectral energy distribution (SED), also fails to explain some HBL-LBL differences (Kollgaard et al. 1996; Stocke 1996). Instead the modern thinking is that the HBL-LBL dichotomy represents two extremes of a continuum in either luminosity and viewing angle (Georganopoulos & Marscher 1998) or luminosity and peak frequency (Fossati et al. 1998).
Recent studies from deeper and larger X-ray surveys have indeed shown that BL Lac objects tend to exhibit more continuous distributions of properties (Sambruna et al. 1996; Scarpa and Falomo 1997; Perlman et al. 1998; Laurent-Muehleisen et al. 1999) rather than disparate ones. This has resulted in an important role for those objects with intermediate SEDs between HBL and LBL, namely intermediate BL Lac objects (IBLs), in revealing BL Lac mysteries.
In this paper, we present the X-ray spectral analysis (ROSAT and ASCA archival data) and the SED of a serendipitous BL Lac object, PKS 2316-423. It is a southern radio source at $`z=0.0549`$, and was formerly classified as a BL Lac candidate on the base of its featureless non-thermal radio/optical continuum (Crawford & Fabian 1994; Padovani & Giommi 1996). We noticed this object as it has been the brightest contaminating source to the nearby narrow-line X-ray galaxy, NGC 7582 (Xue et al. 1998; Schachter et al. 1998) in most of its historical X-ray records.
The ROSAT(PSPC) and ASCA satellites observed this object as a serendipitous source in April 1993 and November 1994 respectively. These observations, though non-simultaneous, extend our knowledge of the sourceโs SED properties to the X-ray domain (0.1โ10 keV), but also provide a good opportunity for X-ray spectroscopic studies, which turn out to be very important for its classification.
In section 2 we describe the observations and data reduction, and then present the spectral analysis in section 3. We construct the source SED in section 4 on the base of our newly X-ray flux measurements as well as the published photometry data. The results are summarized and discussed in section 5. Throughout this paper, H$`{}_{0}{}^{}=50`$ km s<sup>-1</sup>$`\mathrm{Mpc}^1`$ and $`q_0=0.5`$ are assumed. All errors reported below are quoted at the 90% confidence level for one interesting parameter (i.e. $`\mathrm{\Delta }\chi ^2`$$`=2.7`$).
## 2. Observations and data reduction
To investigate the X-ray spectroscopic properties of PKS 2316-423, we collected the relevant data available in the NASA/GSFC data archive. This resulted in one ASCA, one ROSAT PSPC and two HRI pointed observations (see Table 1). All data were reduced using XSELECT within FTOOLS package version 4.2.
PKS 2316-423 is 17 arcmin away from the center of the ROSAT(PSPC) pointed observation, which was centered on NGC 7582. The source was somewhat affected by scattering and vignetting from the PSPC superstructure. We tried to extract the source counts using a larger chip region (with a radius as large as 3.2 arcmin) to account for former effect. The latter, which was estimated a $``$9 % effect at 1 keV, was corrected using a FTOOLSโs task, pcarf. The background was estimated from an annular source-free region. The source spectrum was also corrected for deadtime, and finally regrouped to at least 20 counts per channel. No significant variations in the source count rate were detected over the $``$ 7200-s ROSAT observation spanning a little more than one day.
Both the PKS 2316-423 and NGC 7582 were clearly detected in the ASCA observation with the two GIS and SIS instruments. The former instruments were operated in normal PH mode; the later were operated in 2 CCD mode (along the diagonal) and carefully oriented to place the two objects near the axis.
The ASCA data were selected from the intervals of high and medium telemetry rates for both GIS and SIS. The SIS data were screened using the following criteria: a) The data was not taken in the region of the South Atlantic Anomaly, b) the angle between the field of view and the edge of the bright and dark earth exceeded 25 and 5, respectively, and c) the cutoff rigidity was greater than 4 GeV $`c^1`$. After these selections, we also deleted data if d) there were any spurious events or the dark frame error was abnormal. For GIS data, no bright earth angle and cutoff rigidity criteria were applied. More details concerning the performance and instrumentation of ASCA have been reported in separate papers (ASCA satellite: Tanaka et al. 1994; SIS: Burke et al. 1991; GIS: Ohashi et al. 1996).
Source counts of ASCA observation were extracted from circulars regions of radius $``$ 3 arcmin and $``$ 5 arcmin for the SIS and GIS respectively. The background was extracted from the source free regions on the same detectors. No evident variations in the source count rate were detected either over the $``$19,000-s ASCA observation spanning about half a day.
## 3. Spectral analysis
Since no evident variations in the source count rate were detected during either observations, the time-averaged spectra from both satellites were used for spectral analysis. The spectra of both ROSAT PSPC and ASCA were grouped so that each energy channel contains at least 20 counts allowing chi-squared minimization techniques. Spectral analysis have been performed using XSPEC(V10) program.<sup>1</sup><sup>1</sup>1 http://legacy.gsfc.nasa.gov/docs/xanadu/xspec/u\_manual.html.
### 3.1. The ROSAT Data
A simple power-law fit to the ROSAT PSPC data in the range of 0.1โ2 keV gives photon index $`\mathrm{\Gamma }=2.0\pm 0.2`$ and absorption column density of $`(1.41\pm 0.5)\times 10^{20}`$ $`\mathrm{cm}^2`$ (Figure 1), which is consistent with the Galactic value N$`{}_{\mathrm{Hgal}}{}^{}=2.0\times 10^{20}\mathrm{cm}^2`$ (Stark et al. 1992). This model describes the data reasonably well (Table 2). The fluxes corrected for the Galactic absorption are $`4.20_{0.28}^{+0.25}\times 10^{12}`$ $`\mathrm{ergs}\mathrm{cm}^2\mathrm{s}^1`$ in the 0.1โ2.4 keV band and $`1.83_{0.12}^{+0.11}\times 10^{12}`$ $`\mathrm{ergs}\mathrm{cm}^2\mathrm{s}^1`$ in the 0.5โ2 keV band, estimated from the best-fit power-law model. The inferred intrinsic luminosity is $`5.7\times 10^{43}\mathrm{ergs}\mathrm{s}^1`$ in the 0.1โ2.4 keV band, which is similar to that of other non-quasar AGN. The source was observed twice with ROSAT HRI in June 1992 and May 1993 respectively, the obtained count rates correspond to fluxes of $`(5.9\pm 0.3)\times 10^{12}`$ $`\mathrm{ergs}\mathrm{cm}^2\mathrm{s}^1`$ and $`(6.1\pm 0.3)\times 10^{12}`$ $`\mathrm{ergs}\mathrm{cm}^2\mathrm{s}^1`$ in the observed energy range 0.1- 2.4 keV, assuming a power-low spectrum identical to that of PSPC. The PSPC observation was made between the two HRI observations (Table 1), and recorded a relatively lower flux level. These factors suggest the source is variable and thus there might be a non-thermal origin for the X-ray flux.
### 3.2. The ASCA Data
The spectra of the four ASCA instruments were fitted simultaneously in the range of 0.5โ10 keV for SIS(0/1) and 0.7โ10 keV for GIS(2/3). A simple power-law model gives an acceptable fit to the whole dataset (Table 2). However, comparing the results with the ROSAT data, the ASCA spectra suggest that the source spectral slope with phonton index of $`\mathrm{\Gamma }=2.4\pm 0.2`$ in the 0.5โ10 keV broad-band is significantly steeper than that of ROSAT PSPC spectrum in the 0.1โ2 keV band. In addition, the ASCA data requires an absorption column density of $`8.8_{5.3}^{+5.9}\times 10^{20}`$ $`\mathrm{cm}^2`$ which is significantly in excess of the Galactic value. This is clear evidence of spectral variability, since the two observations were made 1.5 years apart. There are two possibe explanations for this result: either a variable absorbing column (which, as described below, we believe to be unlikely), or a spectral flattening at soft energies due to a convex continuum. However, the first case seems unlikely, since that very few BL Lacs show evidence of significant cold absorbing gas in excess of line-of-sight Galactic column density (Perlman et al. 1996b; Urry et al. 1996). The recent detections of several BL Lacs by the Extreme Ultraviolet Explorer (EUVE) (Marshall, Fruscione, & Carone 1995; Fruscione 1996) is further evidence that these objects do not have significant intrinsic absoption.
For the latter case, the โexcessโ absorption seen in the ASCA spectrum might be an artifact falsely introduced in the fitting process. To test this idea, we next fitted the ASCA data using a broken power-law with a bound absorption at the Galactic value. Fitting with this model is notably improved ($`\mathrm{\Delta }\chi ^2`$$`=6.8`$ for two additional parameters, $`P_F>95\%`$) compared to the fit of the single power-law model with the same bound absorption (Figure 2). The final fit yields two powerlaw components with a break-point at $`2.1_{0.7}^{+0.9}`$ keV. The lower-energy component is flatter with a photon index of $`\mathrm{\Gamma }=2.0_{0.2}^{+0.4}`$; the higher-energy component is steeper with $`\mathrm{\Gamma }=2.6_{0.3}^{+0.3}`$. Thawing the absorption parameter in the broken power-law model produces little change of chi-squared value, and the resulted absorption column density is still consistent with the Galactic value (Table 3).
It should be noted that although the above two models cannot be statistically discerned, consider that the spectral-flattening effect is by far the more likely physical explanation of the โexcessโ absorption seen in the ASCA spectra, hereafter, we would refer the broken power-law model as the best fit to the ASCA data.
Figure 3 shows confidence plots of the powerlaw spectral components versus the corresponding absorption column density for both the ASCA and ROSAT data. The ASCA observation described by the best-fitted model showed that the source brightness decreased by 33% in the 0.1โ2.4 keV range compared with ROSAT PSPC data.
In summary, the ASCA observation of PKS 2316-423 in November 1994 is best-fitted by a broken power-law model, in which a relatively flatter component at lower energies below $`2.1`$ keV is very consistent with the shape of the ROSAT/PSPC spectrum observed in May 1993. This indicates that the broad-band X-ray spectral slope remained constant at lower-energies, even though the source brightness evidently decreased; however, the steeper component above $`2.1`$ keV does suggest the X-ray spectrum became softer at higher energies with the decreasing of the source brightness. As shown in the next section, the ROSAT and ASCA observations actually revealed a variable high-energy tail of the source synchrotron radiation.
## 4. Spectral Energy Distribution
The X-ray spectral analysis of PKS 2316-423 in the last section has shown its HBL-type properties. Moreover, even though non-simultaneous, the composite SED from the literature could present more insights on this object through comparison with the properties of the most recent complete BL Lac object samples. We plot in Figure 4 non-simultaneous radio, optical, UV and X-ray data, from both the space and ground-based observations, assuming an โaverageโ SED of the source. It is clear that the composite SED from radio to X-ray is likely from only one radiation component (i.e., synchrotron emission). Its optical and ultraviolet radiation appear to be a continuation of the radio synchrotron spectrum; the X-ray data are likely from a common emission origin as the lower energy parts and represents a high energy tail of the synchrotron spectrum.
First, we derive some spectral parameters from the SED of PKS 2316-423. Following the general definition, we get the K-corrected two point spectral indices: radio-optical spectral index $`\alpha _{ro}=0.56`$, optical-X-ray one $`\alpha _{ox}=1.18/1.26`$ and radio-X-ray one $`\alpha _{rx}=0.78/0.82`$ (note that hereafter the two values of one parameter is due to the different X-ray fluxes discussed in the last section). Because of the spectral variability shown evidently in the hard X-ray band ( $`>2`$ keV), the composite X-ray/optical spectral, $`\alpha _{xox}=\alpha _{ox}\alpha _x`$ also varied, with a value of $`0.18/0.26`$ for $`\alpha _x=1.0`$<sup>2</sup><sup>2</sup>2Note that the parameter used here refers to the energy spectral index, which is related to the photon spectral index with $`\alpha _x=\mathrm{\Gamma }1`$., and $`\alpha _{xox}=0.42/0.34`$ for $`\alpha _x=1.6`$, respectively. Another important spectral parameter is the X-ray to radio logarithmic flux density ratio, and we get log$`S_X/S_r6.36/6.29`$.
Furthermore, it is clearly seen from Figure 4 that the radio to X-ray SED of PKS 2316-423 likely peaks in the EUV/soft-X-ray band. Following Sambruna et al. (1996), we performed a parabolic fit to the radio through optical/UV/X-ray SED of the source to determine the value of the peak frequency. We obtain $`\nu _p=7.3\times 10^{15}`$ Hz. The same fit yields an estimate of the luminosity at the peak frequency $`L_p=1.5\times 10^{43}\mathrm{ergs}\mathrm{s}^1`$ and of the integrated radio-to-X-ray synchrotron luminosity of $`L_{\mathrm{syn}}=2.1\times 10^{44}\mathrm{ergs}\mathrm{s}^1`$. In the next section we will discuss the general SED properties of BL Lac objects and the attribute of PKS 2316-423.
For comparison, we plotted in Figure 4 the EGRET sensitivity threshold as an upper limit to the GeV flux (marked by an arrow), since the source was never detected at $`\gamma `$-rays. It is shown that the source is dominated by a synchrotron process in the a broad energy range below 10 keV; above which, we know nothing about the inverse-Compton emission which would undoubtedly dominate the source emission.
## 5. Discussions
PKS 2316-423 was first serendipitously detected as a marginally extended (1.5 arcmin) X-ray source in the Einstein Medium-Sensitivity Survey (EMSS; Gioia et al. 1984; Stocke et al. 1991). The ROSAT HRI observation made by Crawford and Fabian (1994) in 1992 June confirmed the result of Einstein and further revealed a second weak X-ray source 10 arcsec (15 kpc) south of the galaxy which may be related to the central source. They suggested that PKS 2316-423 is a BL Lac object based on its probably non-thermal radio/optical/X-ray radiation, and has an โoff-sightโ jet component. If this is verified, then it will provide an interesting test case for both the host galaxies of BL Lac objects and their environment. Unfortunately, the later ROSAT HRI observation made in 1993 May did not go deep enough compared to the former observation, thus the misaligned scenario of this BL Lac object can hardly be further confirmed or constrained.
The multiple X-ray observations show that PKS 2316-423 was variable on timescales of weeks to years. The flux in the 0.1โ2.4 keV band increased by $``$50 % from ROSAT PSPC to HRI observations in a 13 day interval in May 1993, and dropped $``$73 % 1.5 years later when observed with ASCA in November 1994. These factors suggest the source is variable and thus consistent with the expected characteristics of a high-energy tail from a dominant synchrotron emission. The source, however, did not show any short term variability during both ROSAT PSPC and ASCA observations spanned over $``$ 1 day. More X-ray observations will be necessary.
X-ray spectral analysis provides more insights into the nature of the sourceโs emission. A single power-law model fit to the ROSAT PSPC and ASCA data shows significant discrepancies. However, the two datasets can be reconciled if the ASCA spectrum is fit with a broken power-law model. This model indicates that, below an energy-break point, $`2.1`$ keV, there is a relatively flatter component ($`\mathrm{\Gamma }=2.0_{0.2}^{+0.4}`$) to be consistent with the PSPC spectrum; and a steeper component with $`\mathrm{\Gamma }=2.6_{0.3}^{+0.3}`$ above the energy-break point, this is one of the general X-ray spectral properties of a HBL (Sambruna et al. 1996). This kind of intrinsically downward curved X-ray spectrum can be easily interpreted as due to synchrotron losses from a relativistic plasma. The high-energy peaked SED of the source indicates that PKS 2316-423 is a HBL-type BL Lac object, with the X-ray spectrum being the high energy tail of synchrotron emission over a wide wavelength range, with correlated flux and spectral variability. With the decrease in source brightness from ROSAT PSPC to ASCA observations, the X-ray spectrum in 0.1โ10 keV band became softer in the harder X-ray band.
The recent studies have greatly modified our view of BL Lac objects. DXRBS (Perlman et al. 1998), REX (Caccianiga et al. 1999), and RGB (Laurent-Muehleisen et al. 1999) surveys have shown that BL Lac objects exhibit a continuous distribution of properties rather than two distinct classes: HBLs vs LBLs (or XBLs vs RBLs). These findings show that the distribution of the SED parameters of BL Lac objects peaks where the empty region between the two extreme subclasses were previously seen. Therefore, most BL Lac objects should exhibit intermediate SED properties. The previously observed bimodal distribution is primarily due to observational selection effects (e.g. Laurent-Muehleisen et al. 1999).
Here we compare the SED properties of PKS 2316-423 with the general properties of BL Lacs found recently. When we put PKS 2316-423 ($`\alpha _{ro}=0.56`$ and $`\alpha _{ox}=1.18/1.26`$) on the $`\alpha _{ro}`$ vs $`\alpha _{ox}`$ color-color diagram, we find that it is a somewhat intermediate BL Lac object. As we know, $`\alpha _{XOX}`$ can more precisely measure spectral changes from optical to soft X-ray bands. If $`\alpha _{XOX}\genfrac{}{}{0pt}{}{<}{}\mathrm{\hspace{0.17em}0}`$, then the X-rays lie along a powerlaw or steepening synchrotron continuum. A positive value of $`\alpha _{XOX}`$ represents a concave spectrum and is likely caused by a hard inverse-Compton component. The values of $`\alpha _{XOX}`$ for PKS 2316-423 depend on the X-ray spectral slopes in different variability states, being 0.18/0.26 and -0.42/-0.34 for $`\alpha _x=1.0`$, in the high state of the ROSAT/PSPC observation, and 1.6, in the low state of the ASCA observation respectively. We find these values should just locate in the intermediate range of the $`\alpha _{XOX}`$ distribution of recent BL Lac samples (e.g., Perlman et al. 1998; Caccianiga et al. 1999; Laurent-Muehleisen et al. 1999). Interestingly, the value of $`\alpha _{XOX}`$ at different epoch was in opposite symbol, this might indicate that, in the different variability states of PKS 2316-423, the spectral change from optical to X-ray could show either synchrotron or inverse-Compton characteristics, just as would be expected for an IBL.
Previous studies show a clear bimodality in the ratio of the X-ray to radio flux densities of HBLs and LBLs at log$`S_X/Sr5.5`$ (e.g. Padovani and Giommi 1995; Perlman et al. 1996a; Brinkmann et al. 1996). The value of log$`S_X/Sr`$ for PKS 2316-423 is -6.36/-6.29, smaller than this sharp division value, seems in contradiction with its non-LBL attribute; however, the value is consistent with the distribution of the flux ratios of the intermediate BL Lacs found in DXRBS and RGB samples, which show no such dichotomy between HBL- and LBL-like SEDs (Perlman et al. 1998; Laurent-Muehleisen et al. 1999).
The importance of the frequency at which the synchrotron radiation peaks is that it provides a powerful diagnostic for the physical condition of the emitting region. Recent studies showed that among BL Lacs the synchrotron peak frequencies are inversely correlated with their luminosities (Fossati et al. 1998). Due to its low peak luminosity, PKS 2316-423 should be located near the bottom right end of Figure 7c of Fossati et al. (1998), with a peak frequency of around $`10^{18}`$ Hz. However the parabolic fit to the SED of PKS 2316-423 just gives $`\nu _p10^{16}`$ Hz. Consider the given ROSAT photon index of 2.0 (see section 3), which equates to a flat spectrum in $`\nu F_\nu `$ vs $`\nu `$ space within the soft X-ray range, suggesting that the peak frequency of PKS 2316-423 might be underestimated. This could be caused by imposing that the X-ray point smoothly connects to the lower energy data, as would happen in a parabolic fit. If a cubic fit is adopted, the synchrotron peak of PKS 2316-423 could be shifted toward a value of $`10^{17}`$ Hz, which is basically follows the general trend of inverse correlation of the synchrotron peak frequency versus the luminosity among blazars (Fossati et al. 1998). In this sense, however, the source also behaves like a typical HBL.
In summary, the averaged SED properties of PKS 2316-423 indicates that it is an IBL, like those many objects found in DXRBS (Perlman et al. 1998), REX (Caccianiga et al. 1999) and RGB surveys (Laurent-Muehleisen et al. 1999). However, the source also shows some typical X-ray spectral properties of a HBL. This in fact supports a peopleโs general idea โ that in X-rays, an intermediate BL Lac object can show either synchrotron or inverse-Compton characteristics in different variability states. Though the most compelling test of the idea awaits a large sample of such studies.
We would like to thank the anonymous referee for his/her insightful comments that led to a number of significant improvements in this paper. This research has made use of the NASA/IPAC Extra-galactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. S.J.X. acknowledges the financial support from the Chinese Post-Doctoral Program.
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# Encomplexing the writhe
## 1. Introduction
This paper is a detailed version of my preprint , which was written about five years ago. Here I do not discuss results that have appeared since then. I plan to survey them soon in another paper. The subject is now evolving into a real algebraic knot theory.
This paper is dedicated to the memory of my teacher Vladimir Abramovich Rokhlin. It was V. A. Rokhlin, who suggested to me, a long time ago, in 1977, to develop a theory of real algebraic knots. He suggested this as a topic for my second dissertation (after PhD, like habilitation). Following this suggestion, I moved then from knot theory and low-dimensional topology to the topology of real algebraic varieties. However, in the topology of real algebraic varieties, problems on spatial surfaces and plane curves were more pressing than problems on spatial curves, and my second dissertation defended in 1983 was devoted to the constructions of real algebraic plane curves and spatial surfaces with prescribed topology.
The change in the topic occured mainly because I managed to obtain decent results in another direction, on plane curves. There was also a less respectable reason: I failed to relate the traditional techniques of classical knot theory to real algebraic knots. One of the obstacles was a phenomenon which became the initial point of this paper. A large part of the traditional techniques in knot theory uses plane knot diagrams, i.e., projections of knots to the plane. The projection of an algebraic curve is algebraic, and one could try to apply results on plane real algebraic curves. However, the projection contains extra real points, which do not correspond to real points of the knot. These points are discussed below. In the seventies they ruined my weak attempts to study real algebraic knots. Now they allow us to detect crucial differences between topological and real algebraic knots.
I am grateful to Alan Durfee, Tobias Ekholm, and V. M. Kharlamov for stimulating conversations.
The lengthy informal introduction, which follows, is intended to explain the matter prior to going into details. I cannot resist the temptation to write in the style of popular mathematics and apologize to the reader whom this style may irritate.
### 1.1. Knot theory and algebraic geometry
In classical knot theory, by a link one means a smooth closed $`1`$-dimensional submanifold of the $`3`$-dimensional sphere $`S^3`$, i.e., the union of several disjoint circles smoothly embedded into $`S^3`$. A link may be equipped with various additional structures such as orientation or framing and considered up to various equivalence relations like smooth (or ambient) isotopy, PL-isotopy, cobordism or homotopy. See, e.g., or .
In algebraic geometry classical links naturally appear as links of singular points of complex plane algebraic curves. Given a singular point $`p`$ of a complex plane algebraic curve $`C`$, the intersection of $`C`$ with the boundary of a sufficiently small ball centered at $`p`$ is called the link of the singularity. It provides a base for a fruitful interaction between topology and algebraic geometry with a long history and lots of important results.
Another obvious opportunity for interaction between algebraic geometry and knot theory is based on the fact that a classical link may emerge as the set of real points of a real algebraic curve. This opportunity was completely ignored, besides that a number of times it was proved that any classical link is approximated by (and hence isotopic to) the set of real points of a real algebraic curve. There are two natural directions in which algebraic geometry and knot theory may interact in the study of real algebraic links: first, the study of relationships between invariants which are provided by link theory and algebraic geometry, second, developing a theory parallel to the classical link theory, but taking into account the algebraic nature of the objects. From the viewpoint of this second direction, it is more natural to consider real algebraic links up to isotopy consisting of real algebraic links, which belong to the same continuous family of algebraic curves, rather than up to smooth isotopy in the class of classical links. I call an isotopy of the former kind a rigid isotopy, following the terminology established by Rokhlin in a similar study of real algebraic plane projective curves and the likes (see, e.g., the survey ). Of course, there is a forgetting functor: any real algebraic link can be regarded as a classical link and a rigid isotopy as a smooth isotopy. It is interesting to see how much is lost under that transition.
In this paper I point out a real algebraic link invariant which is lost. It is unexpectedly simple. In an obvious sense it is a nontrivial Vassiliev invariant of degree $`1`$ on the class of real algebraic knots (recall that a knot is a link consisting of one component). In classical knot theory the lowest degree of a nontrivial Vassiliev knot invariant is $`2`$. Thus there is an essential difference between classical knot theory and the theory of real algebraic knots. Of course this difference has a simple topological explanation: a real algebraic link is more complicated topologically, besides its set of real points contains the set of complex points invariant under the complex conjugation and a rigid isotopy induces an equivariant smooth isotopy of this set.
The invariant of real algebraic links which is defined below is very similar to the self-linking number of a framed knot. In I call it also the self-linking number. Its definition looks like a replacement of an elementary definition of the writhe of a knot diagram, but taking into consideration the imaginary part of the knot.
### 1.2. The word โencomplexโ
Here I propose to change this name (i.e., self-linking number) to encomplexed writhe, and, in general, since many other characteristics can also be enhanced in a similar way, I suggest a new verb encomplex for similar enhancements by taking into consideration additional imaginary ingredients. This would agree with the general usage of the prefix en- which is described in the Oxford Dictionary of Current English as follows: โen- prefix $`\mathrm{}`$ forming verbs $`\mathrm{}`$ 1 from nouns, meaning โput into or onโ (engulf; entrust; embed), 2 from nouns or adjectives, meaning โbring into the condition ofโ (enslave) $`\mathrm{}`$โ.
The word complexification does not seem to be appropriate for what we do here with the writhe. A complexification of the writhe should be a complex counterpart for the writhe, it should be a characteristic of complex objects, while our enhancement of the writhe is defined only for real objects possessing complexification.
### 1.3. Self-linking and writhe of nonalgebraic knots
The linking number is a well-known numerical characteristic of a pair of disjoint oriented circles embedded in three-dimensional Euclidean space. Roughly speaking, it measures how many times one of the circles runs around the other. It is one of the most classical topological invariants, introduced in the nineteenth century by Gauss .
In the classical theory, a self-linking number of a knot is defined if the knot is equipped with an additional structure like a framing or just a vector field nowhere tangent to the knot.<sup>1</sup><sup>1</sup>1A framing is a pair of normal vector fields on the knot orthogonal to each other. There is an obvious construction that makes a framing from a nontangent vector field and establishes a one-to-one correspondence between homotopy classes of framings and nontangent vector fields. The vector fields are more flexible and relevant to the case. The self-linking number is the linking number of the knot oriented somehow and its copy obtained by a small shift in the direction specified by the vector field. It does not depend on the orientation, since reversing the orientation of the knot is compensated by reversing the induced orientation of its shifted copy. Of course, the self-linking number depends on the homotopy class of the vector field.
A knot has no natural preferable homotopy class of framings, which would allow us to speak about a self-linking number of the knot without a special care on the choice of the framing.<sup>2</sup><sup>2</sup>2Moreover, the self-linking number is used to define a natural class of framings: namely, the framings with self-linking number zero. Some framings appear naturally in geometric situations. For example, if one fixes a generic projection of a knot to a plane, the vector field of directions of the projection appears. The corresponding self-linking number is called the writhe of the knot. However, it depends on the choice of the projection and changes under isotopy.
The linking number is a Vassiliev invariant of order $`1`$ of two-component oriented links. This means that it changes by a constant (in fact, by $`2`$) when the link experiences a homotopy with the generic appearance of an intersection point of the components. Whether the linking number increases or decreases depends only on the local picture of orientations near the double point: when it passes from through to , the linking number increases by $`2`$. Generalities on Vassiliev invariants see, e.g., in .
In a sense the linking number is the only Vassiliev invariant of degree $`1`$ of two-component oriented links: any Vassiliev invariant of degree $`1`$ of two-component oriented links is a linear function of the linking number. Similarly, the self-linking number is a Vassiliev invariant of degree $`1`$ of framed knots (it changes by $`2`$ when the knot experiences a homotopy with a generic appearance of a self-intersection point) and it is the only Vassiliev of degree $`1`$ of framed knots in the same sense. The necessity of a framing for the definition of self-linking number can now be formulated more rigorously: only constants are Vassiliev invariants of degree $`1`$ of (non-framed) knots.
The diagrammatical definition of the writhe, which is imitated below, runs as follows: for each crossing point of the knot projection, one defines a local writhe equal to $`+1`$ if near the point the knot diagram looks like and $`1`$ if it looks like . Then one sums the local writhes over all double points of the projection. The sum is the writhe.
A continuous change of the projection may cause the vanishing of a crossing point. This happens under the first Reidemeister move shown in the left hand half of Figure 1. This move changes the writhe by $`\pm 1`$.
### 1.4. Algebraicity enhances the writhe
If a link is algebraic, then its projection to a plane is algebraic, too. A generic projection has only ordinary double points and the total number of its complex double points is constant.<sup>3</sup><sup>3</sup>3Here by a generic projection we mean a projection from a generic point. When one says that a generic projection has some properties, this means that for an open everywhere dense set of points the projection from any point of this set has these properties. The whole set of undesirable points is closed nowhere dense although it depends on the properties under consideration. A proof is an easy exercise either on Sardโs Lemma, or Bertiniโs Theorem. The number of real double points can vary, but only by an even number. A real double point cannot turn alone into an imaginary one, as it seems to happen under the first Reidemeister move. Under an algebraic version of the first Reidemeister move, the double point stays in the real domain, but becomes solitary, like the only real point of the curve $`x^2+y^2=0`$. The algebraic version of the first Reidemeister move is shown in the right hand half of Figure 1.
It is not difficult to prove that the family of spatial curves that realizes this move can be transformed by a local diffeomorphism to the family of affine curves defined by the following system of equations
$$\{\begin{array}{cc}\hfill xz+y& =0,\hfill \\ \hfill x+z^2+\tau & =0,\hfill \end{array}$$
where $`\tau `$ is the parameter of the deformation. These are rational curves, admitting a rational parametrization
$$\{\begin{array}{cc}\hfill x& =t^2\tau ,\hfill \\ \hfill y& =t(t^2+\tau ),\hfill \\ \hfill z& =t.\hfill \end{array}$$
The projection corresponds to the standard projection $`(x,y,z)(x,y)`$ to the coordinate $`xy`$-plane. It maps these curves to the family of affine plane rational cubic curves defined by $`y^2+x^2(\tau +x)=0`$ with $`\tau `$.
A solitary double point of the projection is not the image of any real point of the link. It is the image of two imaginary complex conjugate points of the complexification of the link. The preimage of the point in the 3-space under the projection is a real line. It is disjoint from the real part of the link, but intersects its complexification in a couple of complex conjugate imaginary points.
In the model of the first Reidemeister move above, $`(0,0)`$ is the double point of the projection for each $`\tau 0`$. If $`\tau <0`$, it is a usual crossing point. Its preimage consists of two real points $`(0,0,\sqrt{\tau })`$ and $`(0,0,\sqrt{\tau })`$. If $`\tau >0`$, it is a solitary double point. Its preimage consists of two imaginary conjugate points $`(0,0,i\sqrt{\tau })`$ and $`(0,0,i\sqrt{\tau })`$, which lie on a real line $`x=y=0`$
Below, in Section 2.2, with any solitary double point of the projection, a local writhe equal to $`\pm 1`$ is associated. This is done in such a way that the local writhe of the crossing point vanishing in the first Reidemeister move is equal to the local writhe of the new-born solitary double point. In the case of an algebraic knot, the sum of local writhes of all double points, both solitary and crossings, does not depend on the choice of projection and is invariant under rigid isotopy. This sum is the encomplexed writhe.
### 1.5. Encomplexed writhe for nonoriented and semi-oriented links
A construction similar to the construction of the encomplexed writhe number of an algebraic knot can be applied to an algebraic link. However in this case there are two versions of the construction.
In the first of these, we define an encomplexed writhe number generalizing the encomplexed writhe number defined above for knots. We consider a link diagram and the sum of local writhes at solitary double points and crossing points where the branches belong the same connected component of the set of real points. At these crossing point, to define a local writhe, we need orientations of the branches. As above, we choose an orientation on each of the components. If we make another choice, at a crossing point for which the branches belong the same component, either both orientations change or none. Hence the local writhe numbers at crossing points of this kind do not depend on the choice. In Section 2 below, we prove that the whole sum of local writhes over crossing points of this kind and solitary double points does not depend on the projection and is invariant under rigid isotopy. We call this sum the encomplexed writhe number of the link $`A`$ and denote by $`w(A)`$.
In the second version of the construction, we consider a real algebraic link which is equipped with an orientation of the set of real points, use these orientations to define local writhe numbers at all crossing points and sum the local writhe numbers over all crossing points and all solitary double points. The result is called the encomplexed writhe number of an oriented real algebraic link. This encomplexed writhe number does not change when the orientation reverses. An orientation considered up to reversing is called a semi-orientation. Thus the encomplexed writhe number depends only on the semi-orientation of the link.
The (semi-)orientation may be an artificial extra structure, but it may also appear in a natural way, say, as a complex orientation, if the set of real points divides the set of real points, see . In fact, the complex orientation is defined up to reversing, so it is indeed a semi-orientation. Another important class of semi-oriented algebraic links appears as transversal intersections of two real algebraic surfaces of degrees $`p`$ and $`q`$ with $`pqmod2`$.
The encomplexed writhe number of (semi-)oriented real algebraic link differs from the encomplexed writhe number of the same link without orientation by the sum of all pairwise linking numbers of the components multiplied by $`2`$: let $`A`$ be a real algebraic link, let $`\overline{A}`$ be the same link equipped with an orientation of its set of real points and $`\overline{A}_1,\mathrm{},\overline{A}_n`$ the (oriented) connected components of this set, then
$$w(\overline{A})=w(A)+2\underset{1ijn}{}\mathrm{lk}(\overline{A}_i,\overline{A}_j).$$
### 1.6. Encomplexed writhe and framings
In the case of a knot, the encomplexed writhe number defines a natural class of framings, since homotopy classes of framings are enumerated by their self-linking numbers and we can choose the framing having the self-linking number equal to the algebraic encomplexed writhe number. I do not know any direct construction of this framing. Moreover, there seems to be a reason for the absence of such a construction. In the case of links, the construction above gives a single number, while framings are enumerated by sequences of numbers with entries corresponding to components.
## 2. Real algebraic projective links
Let $`A`$ be a nonsingular real algebraic curve in $`3`$-dimensional projective space. Then the set $`A`$ of its real points is a smooth closed $`1`$-dimensional submanifold of $`P^3`$, i.e., a smooth projective link. The set $`A`$ of its complex points is a smooth complex $`1`$-dimensional submanifold of $`P^3`$.
Let $`c`$ be a point of $`P^3`$. Consider the projection $`p_c:P^3cP^2`$ from $`c`$. Assume that $`c`$ is such that the restriction to $`A`$ of $`p_c`$ is generic. This means that it is an immersion without triple points and at each double point the images of the branches have distinct tangent lines. It follows from well-known theorems that those $`c`$โs for which this is the case form an open dense subset of $`P^3`$ (in fact, it is the complement of a $`2`$-dimensional subvariety).
The real part $`p_c(A)P^2`$ of the image consists of the image $`p_c(A)`$ of the real part and, maybe, several solitary points, which are double points of $`p_c(A)`$.
### 2.1. The local writhe of a crossing
There is a purely topological construction which assigns a local writhe equal to $`\pm 1`$ to a crossing belonging to the image of only one component of $`A`$. This construction is well-known in the case of classical knots. Here is its projective version. I borrow it from Drobotukhinaโs paper on the generalization of Kauffman brackets to links in projective space.
Let $`K`$ be a smooth connected one-dimensional submanifold of $`P^3`$, and $`c`$ be a point of $`P^3K`$. Let $`x`$ be a generic double point of the projection $`p_c(K)P^2`$ and $`LP^3`$ be the line which is the preimage of $`x`$ under the projection. Denote by $`a`$ and $`b`$ the points of $`LP^3`$. The points $`a`$ and $`b`$ divide the line $`L`$ into two segments. Choose one of them and denote it by $`S`$. Choose an orientation of $`K`$. Let $`v`$ and $`w`$ be tangent vectors of $`K`$ at $`a`$ and $`b`$ respectively directed along the selected orientation of $`K`$.
Let $`l`$ be a vector tangent to $`L`$ at $`a`$ and directed inside $`S`$. Let $`w^{}`$ be a vector at $`a`$ such that it is tangent to the plane containing $`L`$ and $`w`$ and is directed to the same side of $`S`$ as $`w`$ (in an affine part of the plane containing $`S`$ and $`w`$). See Figure 2. The triple $`v`$, $`l`$, $`w^{}`$ is a base of the tangent space $`T_aP^3`$. Define the local writhe of $`x`$ to be the value taken by the orientation of $`P^3`$ on this frame.
The construction of the local writhe of $`x`$ contains several choices. Here is a proof that the result does not depend on them.
We have chosen an orientation of $`K`$. Had the opposite orientation been selected, then $`v`$ and $`w^{}`$ would be replaced by the opposite vectors $`v`$ and $`w^{}`$. This would not change the result, since $`v`$, $`l`$, $`w^{}`$ defines the same orientation as $`v`$, $`l`$, $`w^{}`$.
We have chosen the segment $`S`$. If the other half of $`L`$ was selected, then $`l`$ and $`w^{}`$ would be replaced by the opposite vectors. But $`v`$, $`l`$, $`w^{}`$ defines the same orientation as $`v`$, $`l`$, $`w^{}`$.
The construction depends on the order of points $`a`$ and $`b`$. The other choice (with the same choice of the orientation of $`K`$ and segment $`S`$) gives a triple of vectors at $`b`$. It can be moved continuously without degeneration along $`S`$ into the triple $`w^{}`$, $`l`$, $`v`$, which defines the same orientation as $`v`$, $`l`$, $`w^{}`$. โ
### 2.2. Local writhe of a solitary double point
Let $`A`$, $`c`$, and $`p_c`$ be as in the beginning of Section 2 and let $`sP^2`$ be a solitary double point of $`p_c`$. Here is a construction assigning $`\pm 1`$ to $`s`$. I will also call the result a local writhe of $`s`$.
Denote the preimage of $`s`$ under $`p_c`$ by $`L`$. This is a real line in $`P^3`$ connecting $`c`$ and $`s`$. It intersects $`A`$ in two imaginary complex conjugate points, say, $`a`$ and $`b`$. Since $`a`$ and $`b`$ are conjugate, they belong to different components of $`LL`$.
Choose one of the common points of $`A`$ and $`L`$, say, $`a`$. The natural orientation of the component of $`LL`$ defined by the complex structure of $`L`$ induces an orientation on $`L`$ as on the boundary of its closure. The image under $`p_c`$ of the local branch of $`A`$ passing through $`a`$ intersects the plane of the projection $`P^2`$ transversally at $`s`$. Take the local orientation of the plane of projection such that the local intersection number of the plane and the image of the branch of $`A`$ is $`+1`$.
Thus the choice of one of two points of $`AL`$ defines an orientation of $`L`$ and a local orientation of the plane of projection $`P^2`$ (we can speak only of a local orientation of $`P^2`$, since the whole $`P^2`$ is not orientable). The plane of projection intersects<sup>4</sup><sup>4</sup>4We may think on the plane of projection as embedded into $`P^3`$. If you would like to think on it as on the set of lines of $`P^3`$ passing through $`c`$, please, identify it in a natural way with any real projective plane contained in $`P^3`$ and disjoint from $`c`$. All such embeddings $`P^2P^3`$ are isotopic. transversally $`L`$ in $`s`$. The local orientation of the plane, the orientation of $`L`$ and the orientation of the ambient $`P^3`$ determine the intersection number. This is the local writhe.
It does not depend on the choice of $`a`$. Indeed, if one chooses $`b`$ instead, then both the orientation of $`L`$ and the local orientation of $`P^2`$ would be reversed. The orientation of $`L`$ would be reversed, because $`L`$ inherits opposite orientations from the different halves of $`LL`$. The local orientation of $`P^2`$ would be reversed, because the complex conjugation involution $`conj:P^2P^2`$ preserves the complex orientation of $`P^2`$, preserves $`P^2`$ (point-wise) and maps one of the branches of $`p_c(A)`$ at $`s`$ to the other reversing its complex orientation.
### 2.3. Encomplexed writhe and its invariance
Now for any real algebraic projective link $`A`$, choose a point $`cP^3`$ such that the projection of $`A`$ from $`c`$ is generic and sum the writhes of all crossing points of the projection belonging to the image of only one component of $`A`$ and the writhes of all solitary double points. This sum is called the encomplexed writhe number of $`A`$.
I have to show that it does not depend on the choice of projection. The proof given below proves more: the sum is invariant under rigid isotopy of $`A`$. By rigid isotopy we mean an isotopy consisting of nonsingular real algebraic curves. The effect of a movement of $`c`$ on the projection can be achieved by a rigid isotopy defined by a path in the group of projective transformations of $`P^3`$. Therefore the following theorem implies both the independence of the encomplexed writhe number from the choice of projection and its invariance under rigid isotopy.
###### 2.A Theorem.
For any two rigidly isotopic real algebraic projective links $`A_1`$ and $`A_2`$ whose projections from the same point $`cP^3`$ are generic, the encomplexed writhe numbers of $`A_1`$ and $`A_2`$ defined via $`c`$ are equal.
This theorem is proved in Section 2.5.
###### 2.B Corollary 1.
The encomplexed writhe number of a real algebraic projective link does not depend on the choice of the projection involved in its definition.
###### Proof of 2.B.
A projection depends only on the center from which it is done. The effect on the projection of a movement of the center can be achieved by a rigid isotopy defined by a path in the group of projective transformations of $`P^3`$. โ
Thus the encomplexed writhe number is a characteristic of a real algebraic link.
###### 2.C Corollary 2.
The encomplexed writhe number of a real algebraic projective link is invariant under rigid isotopy.โ
### 2.4. Algebraic counterparts of Reidemeister moves
As in the purely topological situation of an isotopy of a classical link, a generic rigid isotopy of a real algebraic link may be decomposed into a composition of rigid isotopies, each of which involves a single local standard move of the projection. There are $`5`$ local standard moves. They are similar to the Reidemeister moves. The first of these $`5`$ moves is shown in the right hand half of Figure 1. The other moves are shown in Figure 3. The first two of these coincide with the second and third Reidemeister moves. The fourth move is similar to the second Reidemeister move: also two double points of projection come to each other and disappear. However the double points are solitary. The fifth move is similar to the third Reidemeister move: a triple point also appears for a moment. But at this triple point only one branch is real, the other two are imaginary conjugate to each other. In this move a solitary double point traverses a real branch.
### 2.5. Reduction of Theorem 2.A to Lemmas
To prove Theorem 2.A, first replace the rigid isotopy by a generic one and then decompose the latter into local moves described above, in Section 2.4. Only in the first, fourth and fifth moves solitary double points are involved. The invariance under the second and the third move follows from the well-known fact of knot theory that the topological writhe is invariant under the second and third Reidemeister moves. Cf. . Thus the following three lemmas imply Theorem 2.A.
###### 2.D Lemma.
In the fifth move the writhe of the solitary point does not change.
###### 2.E Lemma.
In the fourth move the writhes of the vanishing solitary points are opposite.
###### 2.F Lemma.
In the first move the writhe of vanishing crossing point is equal to the writhe of the new-born solitary point.
### 2.6. Proof of Lemmas 2.D and 2.E
Proof of Lemma 2.F is postponed to Section 2.7. Note that although Lemma 2.F is the most difficult to prove, it is the least significant: here its only role is to justify the choice of sign made in the definition of local writhe in solitary double point of the projection. It is clear that the writhes of vanishing double points involved in the first move are related, and if they were opposite to each other, then the definition of the encomplexed writhe number should be changed, but would not be destroyed irrecoverably.
###### Proof of Lemma 2.D.
This is obvious. Indeed, the real branch of the projection does not interact with the imaginary branches, it just passes through their intersection point. โ
###### Proof of Lemma 2.E.
At the moment of the fourth move take a small ball $`B`$ in the complex projective plane centered in the solitary self-tangency point of the projection of the curve. Its intersection with the projection of the complex point set of the curve consists of two smoothly embedded disks tangent to each other and to the disk $`BP^2`$. Under the move each of the disks experiences a diffeotopy. Before and after the move the intersection the curve with $`B`$ is the union of the two disks meeting each other transversally in two points, but before the move the disks do not intersect $`P^2`$, while after the move they intersect $`P^2`$ in their common points.
To calculate the writhe at both vanishing solitary double points, let us select the same imaginary branch of the projection of the curve passing through the points. This means that we select one of the disks described above. The sum of the local intersection numbers of this disk (equipped with the complex orientation) and $`BP^2`$ (equipped with some orientation) is zero since under the fourth move the intersection disappears, while in the boundary of $`B`$ no intersection happens.
Therefore the local orientations of the projective plane in the vanishing solitary double points defined by this branch define opposite orientations of $`BP^2`$. (Recall that the local orientations are distinguished by the condition that the local intersection numbers are positive.)
On the other hand, under the move the preimages of the vanishing solitary double points come to each other up to coincidence at the moment of the move and their orientations defined by the choice of the same imaginary branch are carried to the same orientation of the preimage of the point of solitary self-tangency. Indeed, the preimages are real lines and points of intersection of their complexifications with the selected imaginary branch of the curve also come to the same position. Therefore the halves of the complexifications containing the points come to coincidence, as well as the orientations defined by the halves on the real lines.
It follows that the intersection numbers of $`B`$ with the preimages of the vanishing solitary double points equipped with these orientations are equal. Since the local orientations of the projective plane in the vanishing solitary double points define distinct orientations of $`BP^2`$, the writhes are opposite to each other.โ
### 2.7. Proof of Lemma 2.F
It is sufficient to consider the model family of curves described in Section 1.4. Recall that the curves of this family are defined by the following system of equations
$$\{\begin{array}{cc}\hfill xz+y& =0,\hfill \\ \hfill x+z^2+\tau & =0,\hfill \end{array}$$
where $`\tau `$ is the parameter of the deformation. These curves admit a rational parametrization
$$\{\begin{array}{cc}\hfill x& =t^2\tau ,\hfill \\ \hfill y& =t(t^2+\tau ),\hfill \\ \hfill z& =t.\hfill \end{array}$$
The projection corresponds to the standard projection $`(x,y,z)(x,y)`$ to the coordinate $`xy`$-plane. It maps these curves to the family of affine plane rational cubic curves defined by $`y^2+x^2(\tau +x)=0`$ with $`\tau `$.
We must prove that the local writhe at $`(0,0)`$ for $`\tau <0`$ coincides with the local writhe at $`(0,0)`$ for $`\tau >0`$.
Let us calculate the local writhe for $`\tau <0`$. Denote $`\sqrt{\tau }`$ by $`\rho `$. The preimage of $`(0,0)`$ consists of points $`a=(0,0,\rho )`$ and $`b=(0,0,\rho )`$ corresponding to the values $`\rho `$ and $`\rho `$ of $`t`$, respectively, see Figure 4. The tangent vectors to the curve at these points are $`v=(2\rho ,2\rho ^2,1)`$ and $`w=(2\rho ,2\rho ^2,1)`$. The vector $`l`$ connecting $`a`$ and $`b`$ is $`(0,0,2\rho )`$. By definition, the writhe is the value taken by the orientation of $`P^3`$ on the frame $`v`$, $`l`$, $`w^{}`$. This value is equal to the value of this orientation on the frame $`(1,0,0)`$, $`(0,1,0)`$, $`(0,0,1)`$ multiplied by the sign of
$$det\left(\begin{array}{ccc}2\rho & 2\rho ^2& 1\\ 0& 0& 2\rho \\ 2\rho & 2\rho ^2& 1\end{array}\right)=16\rho ^4<0.$$
Let us calculate the local writhe for $`\tau >0`$. Denote $`\sqrt{\tau }`$ by $`\rho `$. The preimage of $`(0,0)`$ consists of points $`a^{}=(0,0,i\rho )`$ and $`b^{}=(0,0,i\rho )`$ corresponding to the values $`i\rho `$ and $`i\rho `$ of $`t`$. Choose the branch which passes through $`a^{}`$. It belongs to the upper half of the line $`x=y=0`$, which induces the positive orientation of the real part directed along $`(0,0,1)`$. At $`a^{}`$ the branch of the curve has tangent vector $`v=(2i\rho ,2\rho ^2,1)`$ and the real basis consisting of $`v`$ and $`iv=(2\rho ,2i\rho ^2,i)`$ positively oriented with respect to the complex orientation of this branch. The projection maps this basis to the positively oriented basis $`(2i\rho ,2\rho ^2)`$, $`(2\rho ,2i\rho ^2)`$ of the projection of the branch. The intersection number of this projection and $`^2`$ in $`^2`$ is the sign of
$$det\left(\begin{array}{cccc}0& 2\rho & 2\rho ^2& 0\\ 2\rho & 0& 0& 2\rho ^2\\ 1& 0& 0& 0\\ 0& 0& 1& 0\end{array}\right)=4\rho ^3<0.$$
Hence the orientation of $`^2`$ such that its local intersection number with the selected branch of the projection does not coincide with the orientation defined by the standard basis. The intersection number of the line $`x=y=0`$ with the standard orientation and the $`xy`$-plane with the standard orientation is the value of the orientation of the ambient space $`^3`$ taken on the standard basis $`(1,0,0)`$, $`(0,1,0)`$, $`(0,0,1)`$. Therefore the local writhe is opposite to this value. โ
###### Remark.
There is a more conceptual proof of Lemma 2.F. It is based on a local version of the Rokhlin Complex Orientation Formula, see and . In fact, the original proof was done in that way. However, the Complex Orientation Formula is more complicated than the calculation above.
### 2.8. Encomplexed writhe of an algebraic link as a Vassiliev invariant of degree one
To speak about Vassiliev invariants, we need to fix a connected family of curves, in which links under consideration comprise the complement to a hypersurface. In the case of classical knots one could include all knots in such a family by adjoining knots with self-intersections and other singularities. A singular knot is a right equivalence class of a smooth map of the circle to the space (recall that two maps from a circle are right equivalent if one of them is a composition of a self-diffeomorphism of the circle with the other one).
In the case of real algebraic knots, such a family including all real algebraic knots does not exist. Even the space of complex curves in the three-dimensional projective space consists of infinitely many components. It is impossible to change the homology class realized by the set of complex points of an algebraic curve in $`P^3`$ by a continuous deformation. Recall that the homology class belongs to the group $`H_2(P^3)=`$ and is a positive multiple $`d[P^1]`$ of the natural generator of $`[P^1]H_2(P^3)`$ realized by a line. The coefficient $`d`$ is called the order of the curve. The genus is another numerical characteristic of a complex curve which takes the same value for all nonsingular curves in any irreducible family. As is well known, the nonsingular complex curves of given order and genus in three-dimensional projective space are parametrized by a finite union of quasi-projective varieties. For each of these varieties, one can try to build a separate theory of Vassiliev invariants on a class of nonsingular real algebraic curves whose complexifications are parametrized by points of this variety. (A similar phenomenon takes place in topology: links with different numbers of components cannot be included into a reasonable connected family, and therefore for each number of components there is a separate theory of Vassiliev invariants.)
Among the varieties of algebraic curves in three-dimensional projective space, there are two special families: for each natural number $`d`$ there is an irreducible variety of rational curves of order $`d`$ (recall that a an algebraic curve is called rational if it admits an algebraic parametrization by a line), and for each pair of natural numbers $`p`$ and $`q`$ there is an irreducible variety of curves which can be presented as intersection of surfaces of degrees $`p`$ and $`q`$.
In the class of real algebraic rational curves of order $`d`$, singular curves comprise a discriminant hypersurface in which a generic point is a rational curve such that it has exactly one singular point and this point is an ordinary double point. An ordinary double point may be of one of the following two types: either it is an intersection point of two real branches, or two imaginary conjugate branches.
Any two real algebraic rational nonsingular curves of order $`d`$ can be connected by a path in the space of real rational curves of degree $`d`$ that intersects the discriminant hypersurface only transversally at a finite number of generic points. Such a path can be regarded as a deformation of a curve to the other one. When it intersects the discriminant hypersurface at a point, which is a curve with singularity on real branches, the set of real points of the curve behaves as in classical knot theory: two pieces of the set of real points come to each other and pass through each other. As in classical knot theory, at the moment of intersection, the generic projection of the curve experiences an isotopy. Nothing happens besides that one crossing point becomes for a moment the image of a double point and then changes back into a crossing point, but with the opposite writhe. When the path intersects the discriminant hypersurface at a point, which is a curve with singularity on imaginary branches, two complex conjugate imaginary branches pass through each other. At the moment of passing, they intersect in a real isolated double point. At this moment the set of real points of a generic projection experiences an isotopy. No event happens besides that a solitary double point becomes for a moment the image of a solitary real double point of the curve and then changes back into an ordinary solitary double point of the projection (which is not the image of a real point of the knot), but with the opposite writhe number.
It is clear that the encomplexed writhe number of an algebraic curve changes under a modification of each of these kinds by $`\pm 2`$, with the sign depending only on the local structure of the modification near the double point. This means that the encomplexed writhe number on the family of real rational curves under consideration is a Vassiliev invariant of degree $`1`$.
This is true also for any space of nonsingular real algebraic curves that can be included into a connected family of real algebraic curves by adjoining a hypersurface, penetration through which at a generic point looks as in the family of rational curves described above.
There are many families of this kind besides the families of rational knots. However, in many families of algebraic curves a transversal penetration through the discriminant hypersurface in a generic point looks differently. In particular, for intersections of two surfaces it is a Morse modification of the real part of the curve. At the moment, the old double points of the projection, both solitary and crossing, do not change. An additional double point appears just for a moment. However the division of crossing points to self-crossing points of a single component and crossing points of different components may change. Therefore the encomplexed writhe number changes in a complicated way. If the degrees of the surfaces defining the curve are of the same parity, the real part of the curve has a natural semi-orientation. The Morse modification respects this semi-orientation. Therefore the encomplexed writhe number of the semi-oriented curve does not change.
###### 2.G Theorem.
The encomplexed writhe number of any nonsingular semi-oriented real algebraic link which is a transversal intersection of two real algebraic surfaces whose degrees are of the same parity is zero.
###### Proof.
Any two nonsingular real curves of the type under consideration can be connected by a path as above. Hence their self-linking numbers coincide. On the other hand, it is easy to construct, for any pair of natural numbers $`p`$ and $`q`$ of the same parity, a pair of nonsingular real algebraic surfaces of degrees $`p`$ and $`q`$ transversal to each other in three-dimensional projective space such that their intersection has zero self-linking number. โ
In contrast to this vanishing result, one can prove that the encomplexed writhe number of a real algebraic rational knots of degree $`d`$ can take any value in the interval between $`(d1)(d2)/2`$ and $`(d1)(d2)/2`$ including these limits and congruent to them modulo $`2`$.โ
## 3. Generalizations
### 3.1. The case of an algebraic link with imaginary singularities
The same construction may be applied to real algebraic curves in $`P^3`$ having singular imaginary points, but no real singularities. In the construction we can eliminate projections from the points such that some singular point is projected from them to a real point. Indeed, for any imaginary point there exists only one real line passing through it (the line connecting the point with its complex conjugate), thus we have to exclude a finite number of real lines.
This gives a generalization of encomplexed writhe numbers with the same properties: it is invariant with respect to rigid isotopies (i.e., isotopies made of curves from this class), and is multiplied by $`1`$ under a mirror reflection.
### 3.2. Real algebraic links in the sphere
The construction of this paper can be applied to algebraic links in the sphere $`S^3`$. Although from the viewpoint of knot theory this is the most classical case, from the viewpoint of algebraic geometry the case of curves in the projective space is simpler. The three-dimensional sphere $`S^3`$ is a real algebraic variety. It is a quadric in four-dimensional real affine space. The stereographic projection is a birational isomorphism of $`S^3`$ onto $`P^3`$. It defines a diffeomorphism between the complement of the center of the projection in $`S^3`$ and a real affine space.
Given a real algebraic link in $`S^3`$, one may choose a real point of $`S^3`$ from the complement of the link and project the link from this point to an affine space. Then include the affine space into the projective space and apply the construction above. The image has no real singular points, therefore we can use the result of the previous section.
This construction blows up the center of projection, making a real projective plane out of it, and maps the complement to the center of the projection in the set of real points of the sphere isomorphically onto the complement of the projective plane. In the imaginary domain, it contracts each generatrix of the cone which is the intersection of the sphere with its tangent plane at the center of projection. The image of the cone is an imaginary quadric curve contained in the projective plane which appeared as the result of blowing up of the central point.
### 3.3. Other generalizations
It is difficult to survey all possible generalizations. Here I indicate only two directions.
First, consider the most straightforward generalization. Let $`L`$ be a nonsingular real algebraic $`(2k1)`$-dimensional subvariety in the projective space of dimension $`4k1`$. Its generic projection to $`P^{4k2}`$ has only ordinary double points. At each double point either both branches of the image are real or they are imaginary complex conjugate. If the set of real points is orientable, then one can repeat everything with obvious changes and obtain a definition of a numeric invariant generalizing the encomplexed writhe number defined above.
Let $`M`$ be a nonsingular three-dimensional real algebraic variety with oriented set of real points equipped with a real algebraic fibration over a real algebraic surface $`F`$ with fiber a projective line. There is a construction which assigns to a real algebraic link (i.e., a nonsingular real algebraic curve in $`M`$) with a generic projection to $`F`$ an integer, which is invariant under rigid isotopy, is multiplied by $`1`$ under the orientation reversal in $`M`$ and is a Vassiliev invariant of degree $`1`$. This construction is similar to the one presented above, but uses, instead of the projection to $`P^2`$, an algebraic version of Turaevโs shadow descriptions of links .
### 3.4. Not only writhe can be encomplexed
Here we discuss only one example. However it can be easily generalized. Consider immersions of the sphere $`S^{2n}`$ to $`^{4n}`$. Up to regular homotopy (i.e., a homotopy consisting of immersions whose differentials also comprise a homotopy), an immersion $`S^{2n}^{4n}`$ is defined by its Smale invariant , which is an element of $`\pi _{2n}(V_{4n,2n})=`$. For a generic immersion, it can be expressed as the sum of local self-intersection numbers over all double points of the immersion, see .
Let us encomplex the Smale invariant. For this, first, we have to consider a real algebraic counterpart for the notion of generic immersion $`S^{2n}^{4n}`$. The identification is defined via the universal covering $`^{4n}(S^1)^{4n}`$. Replace Euclidean space $`^{4n}`$ by torus $`(S^1)^{4n}`$, which has the advantage of being compact. The classification of immersions $`S^{2n}(S^1)^{4n}`$ up to regular homotopy coincides with the Smale classification of immersions $`S^{2n}^{4n}`$. The sphere $`S^{2n}`$ is the real part of a quadric projective hypersurface. The torus $`(S^1)^{4n}`$ is the real part of a complex Abelian variety. Consider real regular maps of the quadric to the Abelian variety. A generic map defines an immersion both for the complex and real parts. The only singularities are transversal double points. Double points in the real part of the target variety are of two kinds. At a double point of the first kind two sheets of the image of $`S^{2n}`$ meet. At a double point of the second kind the images two complex conjugate sheets of the complexification of $`S^{2n}`$ meet. The Smale invariant is the sum of the local intersection numbers over the double points of the first kind. One can extend the definition of the local intersection number to the double points of the second kind in such a way that the total sum of the local intersection numbers over double points of both kinds would be invariant under continuous deformations of regular maps.
This total sum is the encomplexed Smale invariant. Notice that it is, in a sense, more invariant than the original Smale invariant. The Smale invariant may change under homotopy, it is invariant only under regular homotopy. The encomplexed Smale invariant does not change under a homotopy in the class of regular maps, which corresponds to the class of all continuous maps.
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# 0.1 Introduction.
## 0.1 Introduction.
Let $`X`$ be a smooth projective curve of genus $`g2`$ defined over an algebraically closed field $`KI.`$ In this paper we study the rank $`r`$ stable vector bundles, E, on $`X`$ such that for some integer $`k`$ with $`0<k<r`$ $`E`$ has a โlargeโ family of subbundles with rank $`k`$ and maximal degree. For positive integers $`r,`$ $`d`$ let $`M(X;r,d)`$ be the moduli space of stable vector bundles on $`X`$ of rank $`r`$ and degree $`d.`$ It is well known that $`M(X;r,d)`$ is smooth and irreducible. For a positive integer $`k`$ with $`0<k<r,`$ let $`M_k(E)`$ be the set of all rank $`k`$ subbundles of $`E`$ with maximal degree. Being a Quot-scheme, $`M_k(E)`$ has a natural scheme-structure. For the intent of this paper we will only need to consider its reduced structure. Indeed we are interested in finding a stable vector bundle $`E`$ such that $`M_k(E)`$ has an irreducible component with prescribed dimension. Since every element in $`M_k(E)`$ has maximal degree, the scheme $`M_k(E)`$ is complete. Hence by \[MS\], pp. 254-255, we have $`\text{dim}(M_k(E))k(rk)`$ for every rank $`r`$ vector bundle $`E.`$ Fixed $`x`$ with $`xk(rk),`$ it is very easy to find a decomposable rank $`r`$ vector bundle $`E`$ such that $`M_k(E)`$ has an irreducible component of dimension $`x.`$ But we are interested in stable vector bundles which are indecomposable. Hence using extensions of a line bundle by a decomposable rank $`r1`$ bundle we will prove in section $`2`$ the following result:
###### Theorem 0.1.1
Fix integers $`g,`$ $`r,`$ $`k`$ with $`2gr+1,`$ $`0<k<r`$; if $`2kr,`$ then assume $`x(k1)(r2k+1);`$ if $`2k>r,`$ then assume $`x(rk1)(2kr+1).`$ Let $`X`$ be a smooth projective curve of genus $`g.`$ Then there exists a stable vector bundle $`E`$ on $`X`$ such that $`M_k(E)`$ has an irreducible component of dimension $`x.`$
The proof of Theorem 0.1.1 is quite simple but even if we tried we were not able to produce larger families of maximal degree subbundles. The bound on the dimension $`x:=\text{dim}(M_k(E))`$ seems to be quite good, (see Proposition 0.3.6). The dimension of $`M_k(E)`$ is known when $`E`$ is a general element of $`M(X;r,d)`$ (see Remark 0.2.1 and Proposition 0.2.2). Classically the picture was clear for a rank $`2`$ stable vector bundle $`E:`$ either $`\text{dim}(M_1(E))=0`$ or $`\text{dim}(M_1(E))=1`$ (see the introduction of \[LN\] and references therein). In fact the situation is described by one invariant, called degree of stability, $`s(E).`$ It is known that $`0<s(E)g`$ and $`s(E)\mathrm{deg}(E)(2)`$ (\[Na\]). Furthermore, for $`E`$ general in its moduli space we have $`s(E)=g`$ if $`gd`$ is even and $`s(E)=g1`$ if $`gd`$ is odd. Maruyama proved two main facts: if $`s(E)=g,`$ then $`\text{dim}(M_1(E))=1`$ and if $`s(E)<g`$ then dim$`(M_1(E))=0`$. H. Lange and M.S. Narasimhan produced examples of stable rank 2 vector bundles with $`\text{dim}(M_1(E))=0`$ and $`s(E)<g`$ (see \[LN\], Prop. 3.3. and sections $`5,`$ $`6`$ and $`7).`$ Indeed taking $`f:XY`$ a multiple covering of curve $`Y`$ of genus $`g^{}2`$ they were able to produce examples of curves $`X`$ of genus $`g`$ big enough to obtain a stable rank $`2`$ vector bundle, $`E,`$ on $`X`$ with $`s(E)<g`$ and $`\text{dim}(M_1(E))=1,`$ by pulling back a stable vector bundle, $`F,`$ on $`Y`$ with $`s(F)=g^{}`$ (see \[LN\], Prop. 7.3). In \[Bu\] D. Butler proved some kind of reverse question: if $`E`$ is a stable vector bundle of rank $`2`$ with $`\text{dim}(M_1(E))=1`$ and $`s(E)(2s(E)1)<g`$ then there is a covering $`f:XY`$ and a stable vector bundle on $`Y,`$ $`F,`$ with $`R\text{Pic}(X)`$ with $`ARf^{}(B)`$ and $`\text{dim}(M_1(F))=1.`$ In higher rank the situation is more complicated (see Remark 0.2.1). In particular the stability condition for a rank $`r`$ vector bundle, $`E,`$ is controlled by $`r1`$ invariants called degrees of stability (or Langeโs invariants):
$$s_k(E)=k\mathrm{deg}(E)r\underset{\begin{array}{c}HE\hfill \\ \text{ }\text{rk }H=k\hfill \end{array}}{\mathrm{min}}\mathrm{deg}(H).$$
In section $`3`$ we give a partial generalization to higher rank of a theorem of D. C. Butler (see Theorem 0.3.4) which gives how restrictive is to have โmany and very spreadโ maximal degree subbundles. This is the key motivation of our paper: Theorem 0.3.4 and Proposition 0.3.6 show the existence of a rank $`r`$ stable vector bundle, $`E,`$ with a low value of $`s_k(E)`$ and large dimension of $`M_k(E).`$
This research was partially supported by MURST. Both authors are members of the VBAC Research group of Europroj.
## 0.2 Proof of Theorem 0.1.1
Before proving Theorem 0.1.1 we need the following remark
###### Remark 0.2.1
Assume char$`KI=0.`$ Fix some integers $`g,`$ $`r,`$ $`k,`$ $`a,`$ $`b`$ with $`g3,`$ $`r2,`$ $`0<k<r`$ and $`kba(rk)>0.`$ Let $`X`$ be a smooth projective curves of genus $`g.`$ Let $`A`$ be a general member of $`M(X;k,a),`$ $`B`$ a general member of $`M(X;rk,b)`$ and $`E`$ a general extension of $`B`$ by $`A.`$ If $`kba(rk)<k(rk)(g1),`$ by \[RT\], Thm.0.1, $`E`$ is stable (see also \[BL\] for several special cases). Furthermore, by a result of A. Hirschovitz (\[Hi\]) a general member of $`M(X;r,a+b)`$ is an extension of a general $`BM(X;rk,b)`$ by a general $`AM(X;k,a)`$ if and only if $`kba(rk)k(rk)(g1).`$ As remarked in the introductions of \[RT\] and \[BL\] (\[BL\] eq. (D)), the stability of such an $`E`$ implies dim$`(M_k(E))=\mathrm{max}\{sk(rk)(g1),0\}.`$ In fact $`M_k(E)`$ turns out to be the fiber of a morphism, $`\varphi ,`$ between the parameter space of stable extensions of stable vector bundles and the moduli space $`M(X;r,d);`$ this allows to estimate the dimension of $`M_k(E).`$ In particular if $`s=k(rk)g`$ then $`\text{dim}(M_k(E))=k(rk)`$ which by \[MS\], pp 254-255, it is the maximum admissible dimension of $`M_k(E).`$
If char$`KI=0`$ there exists a first weak version of theorem 0.1.1:
###### Proposition 0.2.2
Assume char$`KI=0`$. Fix integers $`r,`$ $`k,`$ $`x`$ with $`0<k<r,`$ $`0xk(rk)`$ and $`x`$ divisible by the highest common divisor, $`u,`$ of $`k`$ and $`r.`$ Let $`X`$ be a smooth curve of genus $`g3.`$ Then there exists an integer $`d`$ such that for a general $`EM(X;d,r)`$ the algebraic set $`M_k(E)`$ has an irreducible component of dimension $`x`$ and every irreducible component of $`M_k(E)`$ ha dimension at most $`x.`$
Proof. Since $`u`$ divides $`x,`$ there exists an integer $`d`$ with $`0d<r.`$ Moreover there exists an unique integer $`a`$ satisfying $`\frac{da}{rk}g\frac{a}{k}\frac{da}{rk}g+1.`$ Hence as pointed out in 0.2.1 we have $`\text{dim}(M_k(E))=x=\mathrm{max}\{sk(rk)(g1),0\}`$ with $`s=(da)ka(rk).`$
Proof of 0.1.1 Since the cases $`k=1`$ and $`k=r1`$ are covered by Proposition 0.2.2, when char$`KI=0`$ and $`g3,`$ we may assume $`k2`$ and $`rk2.`$ Furthermore, $`M_k(E)M_{rk}(E^{})`$ for every rank $`r`$ vector bundle $`E.`$ Therefore taking, if necessary, the dual bundle, we may assume $`2kr.`$ If char$`KI>0`$ or $`g=2`$ and $`k=1`$ or $`k=r1`$ proceed as in the last part of case 2) below. Hence from now on we may assume $`42kr.`$ Since $`x(k1)((rk)(k1))`$ we can find two integers $`y`$ and $`t`$ with $`0<2tyrk,`$ $`tk1`$ and $`t(y1t)xt(yt).`$ Set $`e:=xt(y1t)`$. Then $`0e<t`$ and if $`y=rk`$ then $`e=0.`$ Therefore $`y+e+1r1.`$ Take a general $`(rey1)`$-ple $`(M,R_1,\mathrm{},R_{rey1})\text{Pic}^0(X)\times \mathrm{}\times \text{Pic}^0(X)`$ and $`L\text{Pic}^1(X)`$ with $`h^0(X,L)=0.`$ Set $`F:=๐ช_X^yM^{(e+1)}(_{1irey1}R_i)`$ (notice that $`y+e+1r1)`$. By construction $`F`$ is a semi-stable vector bundle with rk $`F=r1`$ and $`\mathrm{deg}F=0.`$ Let $`E`$ be a general extension of $`L`$ by $`F.`$
Claim. $`E`$ has no proper subsheaf with positive degree and every degree $`0`$ subsheaf of $`E`$ is a subsheaf of $`F.`$
Here we assume the Claim. Hence $`E`$ is stable. Choose some integers $`u,`$ $`v`$ with $`0uy,`$ $`0ve+1`$ and $`0kuvrey2.`$ Let $`I`$ any subset of $`\{1,\mathrm{},rey2\}`$ with card$`(I)=kuv.`$ Call $`T(u,v,I)`$ the following family of rank $`k`$ subbundles of $`F`$ with degree $`0:`$ $`AT(u,v,I)`$ if and only if $`AA_1A_2A_3`$ where $`A_1`$ subsheaf of $`๐ช^y`$ isomorphic to $`๐ช^u`$ $`A_2`$ is a subsheaf of $`M^{(e+1)}`$ isomorphic to $`M^v`$ and $`A_3_{iI}R_i.`$ Since $`F`$ is polystable and no two among the degree $`0`$ line bundles $`๐ช_X,`$ $`M`$ and $`R_i,`$ $`1irye2,`$ are isomorphic, then $`T(u,v,I)`$ is an irreducible component of $`M_k(E)`$ with $`\text{dim}(T(u,v,I))=u(yu)+(e+1v)v.`$ Varying $`u,`$ $`v`$ and $`I`$ we obtain in this way all the irreducible components of $`M_k(F).`$ By the second part of the Claim, these are the irreducible components of $`M_k(E).`$ When $`u=t`$ and $`v=1`$ by the definition of $`e`$ we get $`\text{dim}(T(t,1,I))=x.`$ Hence to prove 0.1.1 it is sufficient to prove the Claim.
Proof of the Claim. We move the line bundles $`M`$ and $`R_i`$ $`1irey2,`$ in Pic$`{}_{}{}^{0}(X).`$ By the semicontinuity of the Langeโs invariants $`s_k`$ (\[La\], Lemma 1.3), it is sufficient to prove the Claim for the following general extension,
$`0๐ช_X^{(r1)}GL0`$ (1)
.
Since $`h^0(X,L)=0,`$ we have $`h^0(X,G)=r1.`$ In particular the subsheaf $`๐ช_X^{(r1)}`$ is the subsheaf spanned by $`H^0(X,G).`$ Hence it is uniquely determinated by $`G`$ and sent into itself by any endomorphism of $`G.`$ Therefore $`G`$ fits in a unique way into 1, up to an element of Aut$`(G).`$ Since $`\chi (L^{})=g`$ and by our assumptions on $`g`$ and $`r,`$ $`G`$ contains no factor isomorphic to $`๐ช_X.`$ In order to obtain a contradiction we assume the existence of a proper subsheaf $`B`$ of $`G`$ with $`\mathrm{deg}(B)0,`$ and if $`\mathrm{deg}B=0`$ we suppose that $`B`$ is not a direct factor of $`๐ช_X^{(r1)}.`$ Taking $`h:=`$ rk $`B`$ minimal among all the ranks of such subbundles, we may assume $`B`$ stable. Taking $`\mathrm{deg}(B)`$ maximum among all the degrees of all such rank $`h`$ subbundles we may assume $`B`$ saturated in $`G.`$ Since $`B`$ is not contained in $`๐ช_X^{(r1)},`$ the map $`\pi :BL`$ induced by the surjection $`j:GL`$ in 1 is not zero. Set $`B^{}:\text{Ker }(\pi ),`$ $`L^{}=\text{Im }(\pi )`$ and $`w:=h^0(X,B^{}).`$ Since $`B^{}`$ is a subsheaf of $`๐ช_X^{(r1)},`$ we have $`B^{}B^{\prime \prime }๐ช_X^{(w)}`$ with $`h^0(X,B^{\prime \prime })=0.`$ Since $`B^{}`$ is spanned, $`det(A^{})`$ is spanned. Thus if $`\mathrm{deg}(B^{})=\mathrm{deg}(det(B^{})0,`$ $`X`$ has a degree $`\mathrm{deg}(B^{})`$ pencil. By our assumption on the degree of $`B`$ we have $`\mathrm{deg}(B^{})\mathrm{deg}(L^{})\mathrm{deg}(L)=1.`$ Since $`g>0`$ there is no degree $`\mathrm{deg}(B^{})`$ pencil on $`X.`$ Hence a contradiction. Thus $`\mathrm{deg}(B^{})=0,`$ that is $`w=h1`$ and $`B^{}๐ช_X^{(h1)}.`$ At this point we distinguish two cases:
Case 1) Here we assume $`L\simeq ฬธL^{},`$ that is the existence of a positive divisor $`D`$ with $`L^{}=L(D).`$ Since $`\mathrm{deg}(L^{})\mathrm{deg}(L)1=0,`$ $`\mu (B)0`$ and $`B`$ is stable, we obtain a contradiction, unless $`h=1,`$ $`BL^{}`$ and $`w=0.`$ In this case we have $`L^{}L(P)`$ for some $`PW`$ and $`F`$ a positive elementary transformation of $`๐ช_X^{(r1)}L(P)`$ supported in $`P.`$ Hence the set of all such bundles $`G`$ depends at most on $`r`$ parameters. Since $`\text{dim}(\text{Ext}^1(L,๐ช_X^{(r1)})))=(r1)g`$ by the Riemann-Roch Theorem and any such $`G`$ fits, up to a multiplicative constant, in a unique exact sequence 1, we get a contradiction concluding the proof in Case 1).
Case 2) Here we assume $`LL^{}.`$ Then since $`B^{}๐ช_X^{(w)}`$ as direct factor of $`๐ช_X^{(r1)}`$ we get $`G/B๐ช_X^{(r1w)}=๐ช_X^{(rh)}.`$ Hence $`G/B`$ is isomorphic to a direct factor of $`G.`$ But $`G`$ cannot have any trivial factor which is a contradiction and the theorem is proved.
###### Remark 0.2.3
The proof of 0.1.1 shows the existence of a vector bundle $`EM(X;r,1)`$ such that $`M_k(E)`$ has an irreducible component $`t`$ of dimension $`x`$ and such that every $`BT`$ is a direct sum of line bundles of degree $`0.`$
###### Remark 0.2.4
Let $`TM_k(E)`$ be an irreducible subvariety such that there is a subbundle $`F`$ of $`E`$ containing every $`BT.`$ By \[MS\], pp. 254-255, it follows $`\text{dim}(T)k(rk).`$ In the proof of Theorem 0.1.1 we have constructed a vector bundle $`E`$ which has a subbundle $`F`$ with exactly this property.
We repeat here the description of the irreducible components of $`M_k(E)`$ for the stable bundle, $`E,`$ obtained in the proof of Theorem 0.1.1. First choose integers $`u,`$ $`v`$ with $`0uy,`$ $`0ve+1,`$ $`0kuvrey2.`$ Then choose any subset, $`I,`$ of $`\{1,\mathrm{},rey2\}`$ with card$`(I)=kuv.`$ For any such data $`(u,v,I)`$ there is an irreducible component, $`T(u,v,I)`$ of $`M_k(E)`$ and every irreducible component of $`M_k(E)`$ arises in this way. Furthermore, we have $`\text{dim}(T(u,v,I))=u(yu)+(e+1v)v.`$
## 0.3 Maximally spread families and multiple covering curves
In this section we will give a partial generalization of a result of D. C. Butler, \[Bu\]. As in \[Bu\] we will use a result of Accola (\[Ac\]) which is valid in characteristic zero. Therefore we assume char$`KI=0.`$ Let $`X`$ be a smooth projective curve of genus $`g2.`$ Fix two integers $`k,`$ $`r`$ with $`0<k<r`$ and set $`m:=GCD(k,rk),`$ $`v:=\frac{rk}{m}`$ and $`w:=\frac{k}{m}.`$ Let $`E`$ be a rank $`r`$ vector bundle on $`X`$ and $`:=\{H_t\}_{tT}`$ be a flat family of saturated rank $`k`$ subbundles of $`E`$ parameterized by an irreducible complete variety $`T.`$ For every $`tT`$ set $`G_t:=E/H_t.`$ For all pairs $`(x,y)T^2`$ the composition of the inclusion $`i_x:H_xE`$ with the surjection $`j_x:EG_y`$ gives a map $`\varphi (x,y):H_xG_y`$ such that $`\varphi (x,y)=0`$ if and only if $`H_x`$ and $`H_y`$ are isomorphic subsheaf of $`E.`$ More generally, for all $`(x(1),\mathrm{},x(v),y(1),\mathrm{},y(w))T^{v+w}`$ we have a map $`\mathrm{\Phi }((x(1),\mathrm{},x(v),y(1),\mathrm{},y(w))):H_{x(1)}\mathrm{}H_{x(v)}G_{y(1)}\mathrm{}G_{y(w)}.`$ Notice that $`H_{x(1)}\mathrm{}H_{x(v)}`$ and $`G_{y(1)}\mathrm{}G_{y(w)}`$ have the same rank $`\frac{k(rk)}{m}.`$
###### Definition 0.3.1
The family $``$ is called maximal spread if for general $`(x(1),\mathrm{},x(v),y(1),\mathrm{},y(w))T^{v+w}`$ the map $`\mathrm{\Phi }((x(1),\mathrm{},x(v),y(1),\mathrm{},y(w)))`$ is invertible at a general point of $`X.`$
###### Remark 0.3.2
If $`r=2k`$ maximally spread means that for general $`(x(1),y(1))T^2`$ the map $`H_{x(1)}G_{y(1)}`$ is an injective map of sheaves, which is a condition that may be satisfied.
By definition a maximal spread family $``$ induces an inclusion of sheaves of $`H_{x(1)}\mathrm{}H_{x(v)}`$ in $`G_{y(1)}\mathrm{}G_{y(w)}.`$ If $``$ is maximal spread then the map
$$det(\mathrm{\Phi }((x(1),\mathrm{},x(v),y(1),\mathrm{},y(w)))):det(H_{x(1)}\mathrm{}H_{x(v)})det(G_{y(1)}\mathrm{}G_{y(w)})$$
is an inclusion. Therefore there is an effective divisor, $`Z((x(1),\mathrm{},x(v),y(1),\mathrm{},y(w))),`$ associated to a line bundle isomorphic to $`det(H_{x(1)}\mathrm{}H_{x(v)})^{}det(G_{y(1)}\mathrm{}G_{y(w)}).`$ Hence
$$\begin{array}{c}\mathrm{deg}(Z((x(1),\mathrm{},x(v),y(1),\mathrm{},y(w))))=w(\mathrm{deg}(G_t))v(\mathrm{deg}(H_t))=\hfill \\ \\ =w(\mathrm{deg}(E)\mathrm{deg}(H_t))v(\mathrm{deg}(H_t)=\frac{(k(\mathrm{deg}(E)r(\mathrm{deg}(H_t))}{m}.\hfill \end{array}$$
Hence if $`H_t`$ is maximal (that is has maximum degree among rank $`k`$ subbundles of $`E)`$ then $`\mathrm{deg}(Z((x(1),\mathrm{},x(v),y(1),\mathrm{},y(w))))=\frac{s_k(E)}{m}.`$ The divisor $`Z((x(1),\mathrm{},x(v),y(1),\mathrm{},y(w)))`$ depends symmetrically on the variables $`x(i)T,`$ $`1iv,`$ and $`y(j)T,`$ $`1jw.`$ Notice that we have defined the divisors $`Z((x(1),\mathrm{},x(v),y(1),\mathrm{},y(w)))`$ in a general open set of $`T^{v+w}.`$ Since $`T`$ is complete the set of effective divisors $`Z((x(1),\mathrm{},x(v),y(1),\mathrm{},y(w)))`$ has limits for all $`(x(1),\mathrm{},x(v),y(1),\mathrm{},y(w))T^{v+w}.`$ These limits are not unique, but this does not effect our computation. In particular for every $`xT,`$ we may find $`Z(x,\mathrm{},x,x,\mathrm{},x)`$ an effective divisor such that $`๐ช(Z(x,\mathrm{},x,x,\mathrm{},x))det(H_x)^vdet(G_x)^w.`$
###### Remark 0.3.3
Notice that for every $`(x(1),\mathrm{},x(v),y(1),\mathrm{},y(w))T^{v+w}`$ the divisor
$$(v+w)Z((x(1),\mathrm{},x(v),y(1),\mathrm{},y(w)))$$
and the divisor
$$\underset{1iv}{}Z((x(i),\mathrm{},x(i),x(i),\mathrm{},x(i)))+\underset{0jw}{}Z((y(j),\mathrm{},y(j),y(j),\mathrm{},y(j)))$$
are associated to the same line bundle
$$det(H_{x(1)}\mathrm{}H_{x(v)})^{}det(G_{y(1)}\mathrm{}G_{y(w)})^{(v+w)}$$
and therefore they are linearly equivalent. Call $`L((x(1),\mathrm{},x(v),y(1),\mathrm{},y(w)))`$ the subsheaf of $`det(H_{x(1)}\mathrm{}H_{x(v)})^{}det(G_{y(1)}\mathrm{}G_{y(w)})^{(v+w)}`$ spanned by $`H^0(X,det(H_{x(1)}\mathrm{}H_{x(v)})^{}det(G_{y(1)}\mathrm{}G_{y(w)})^{(v+w)}).`$ We believe that the two families of line bundles $`\{det(H_{x(1)}\mathrm{}H_{x(v)})^{}det(G_{y(1)}\mathrm{}G_{y(w)})\}`$ and $`\{L((x(1),\mathrm{},x(v),y(1),\mathrm{},y(w)))|(x(1),\mathrm{},x(v),y(1),\mathrm{},y(w))T^{v+w}\}`$ give more information on the geometry of $`E`$ then $`s_k(E)`$ (even in the case in which $`M_k(E)`$ is finite).
###### Theorem 0.3.4
Assume char$`KI=0.`$ Let $`X`$ be a smooth projective curve of genus $`g2`$ and $`EM(X;r,d),`$ $`r2,`$ such that $`M_k(E)`$ has a maximal spread family, $`T,`$ and such that $`s_k(E)(s_k(E)m)<m^2g,`$ where $`m:=GCD(k,r).`$ Then there exist a smooth curve $`C`$ and a morphism $`\pi :XC`$ with $`\mathrm{deg}(\pi )>1.`$
###### Remark 0.3.5
As one can easily see we are going to prove more then what is stated in the Theorem 0.3.4. In fact we are going to prove that there exists a family of line bundles $`R(x(1),\mathrm{},x(v),y(1),\mathrm{},y(w))\text{ Pic }(C)`$ such that $`\pi ^{}(R(x(1),\mathrm{},x(v),y(1),\mathrm{},y(w)))det(H_{x(1)}\mathrm{}H_{x(v)})^{}det(G_{y(1)}\mathrm{}G_{y(w)}).`$ If the rank of $`E`$ is $`2`$ the existence of this family (with $`w=v=1)`$ allows to construct a rank $`2`$ stable vector bundle $`F`$ on $`C`$ whose pull-back is $`E`$ and whose family of maximal degree linebundles is the pull-back of the one of $`E,`$ up to a twist by a line bundle, $`A,`$ on $`C,`$ (see \[Bu\]).
Proof. Set $`v:=\frac{rk}{m}`$ and $`w:=\frac{k}{m}`$ and take general $`(x(1),\mathrm{},x(v),y(1),\mathrm{},y(w))T^{v+w}.`$ By Remark 0.3.3 we have
$$h^0(det(H_{x(1)}\mathrm{}H_{x(v)})^{}det(G_{y(1)}\mathrm{}G_{y(w)})^{(v+w)})2.$$
As in Remark 0.3.3 consider the line bundles $`L((x(1),\mathrm{},x(v),y(1),\mathrm{},y(w)));`$ they form an infinite family of spanned non-trivial line bundles with degreee at most $`\frac{s_k(E)}{m}.`$ Since $`\frac{s_k(E)}{m}(\frac{s_k(E)}{m}1)<g,`$ we can apply a result of Accola (see \[Ac\], Th. 4.3, or \[Bu\], Lemma 1.2) finding a non-trivial covering $`\pi :XC`$ and $`R(x(1),\mathrm{},x(v),y(1),\mathrm{},y(w))\text{ Pic }(C)`$ with $`\pi ^{}(R(x(1),\mathrm{},x(v),y(1),\mathrm{},y(w)))det(H_{x(1)}\mathrm{}H_{x(v)})^{}det(G_{y(1)}\mathrm{}G_{y(w)}).`$
To explain the notion of maximally spread family, we prove the following easy result
###### Proposition 0.3.6
(any char$`KI)`$ Let $`X`$ be a smooth projective curve of genus $`g2.`$ Fis integers $`r,`$ $`k`$ with $`0<k<r`$ and a rank $`r`$ vector bundle $`E`$ on $`X.`$ Let $`TM_k(E)`$ be an irreducible projective family with dim$`(T)>k(r1k).`$ Then $`T`$ is maximally spread. Furthermore, for every $`PX`$ the union of the subspaces $`H_t{}_{|_{\{P\}}}{}^{}E_{|_{\{P\}}}`$ is not contained in a lower dimensional vector subspace of $`E_{|_{\{P\}}}.`$
Proof. Fix $`PX.`$ By the proof of Proposition of pg 254 in \[MS\], the map
$$\pi :T\text{Grass }(rk,E_{|_{\{P\}}})$$
sending $`H_t,`$ $`tT,`$ into the $`(rk)`$dimensionl vector space $`E_{|_{\{P\}}}/H_t_{|_{\{P\}}}`$ is finite. Since $`\text{dim}(T)>k(rk)=\text{Grass }(rk,E_{|\{P\}})`$ the union of all subspaces $`H_t_{|_{\{P\}}}`$ for $`tT`$ cannot be contained in a hyperplane of $`E_{|_{\{P\}}}.`$
E. Ballico, Universitร di Trento
Dip. di Matematica - 38050 Povo (TN) - Italy
e-mail: ballico@alpha.science.unitn.it
E. Russo, Universitร di Trento
Dip. di Matematica - 38050 Povo (TN) - Italy
e-mail: russo@degiorgi.science.unitn.it
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# Synthetic post-Asymptotic Giant Branch evolution: basic models and applications to disk populations
## 1 Introduction
The post-asymptotic giant branch (post-AGB) phase of evolution, when stars leave the red giant region to eventually become white dwarf (WD) remnants, is characterized by a series of intriguing events and diversities. The beginning of the post-AGB phase is conventionally set at the time of cessation of the high mass loss rate episode $`\dot{M}10^5M_{}\mathrm{yr}^1`$ which has almost completely stripped the star of its hydrogen-rich envelope. After this superwind (SW) quenching, the central star shrinks in radius at almost constant luminosity, thus getting hotter and hotter. Initially it is highly obscured by the fossil superwind around it, but in $`1001000`$ yr this circumstellar material becomes optically thin and the star returns to be observable also at optical wavelengths. The so-called proto-planetary nebula phase (proto-PN) begins at SW quenching and ceases when the star becomes hot enough ($`T_{\mathrm{eff}}10,000`$ K) to start photo-ionizing the materials that were ejected at relatively low velocity ($`1020`$ km s<sup>-1</sup>) during the superwind phase. As cascade recombination and free-free emission make the nebular material shining in the optical, the object has turned into a planetary nebula (PN). The hot star is now emanating a very fast (several 1000 km s<sup>-1</sup>), radiatively-driven wind which dynamically interacts with the fossil superwind, occasionally shaping it in extravagant forms. As nebular expansion continues, the luminosity and surface brightness of the nebula drop as a result of its decreasing emission measure and increasing transparency to the ionizing photons. A time inevitably comes that its surface brightness becomes too dim for the nebulosity to be noticed by terrestrial astronomers, and the still hot post-AGB star is now slowly evolving towards its final WD configuration.
Occasionally, however, before becoming a WD a final thermal pulse of the still active helium burning shell injects enough energy into the outer layers to cause a dramatic expansion of the post-AGB star, and sending it back towards the red giant region for a while. Depending on the details of this metamorphosis, this residual hydrogen-rich envelope can be engulfed by the underlying convective helium shell, with its hydrogen being then diluted and burned, and the energy thus released causing further expansion. A hydrogen-deficient carbon star \- a class of objects including RCrB stars - is thus formed, but even when the ingestion of the hydrogen envelope does not take place, modest wind mass loss suffice to remove virtually all this residual envelope, thus exposing a bare core with a helium, carbon and oxygen atmosphere emanating a Wolf-Rayet like spectrum. Therefore, depending on whether or not post-AGB stars miss to experience such a final thermal pulse, stars enter the WD stage with or without a hydrogen dominated atmosphere, a condition that certainly has to do with the observed dichotomy of WDs into the DAs and the non-DAs groups (respectively with and without hydrogen). Life and times of AGB and post-AGB stars have been extensively reviewed by Iben & Renzini (1983, IR) and Iben (1993, 1995).
Besides its interest per se, the post-AGB phase in stellar evolution has also a number of interesting connections to other important astrophysical issues. For example the PN composition provides an essential tool to study the previous nucleosynthesis which has taken place during the AGB phase and before, thus offering important input for modeling the chemical evolution of galaxies. Moreover, PNs are used as tracers of stellar populations, e.g., for determining the stellar death rate of galaxies, and are being used as standard candles in extragalactic distance determinations. Finally, after neutron stars, post-AGB stars are the hottest objects a stellar population can generate, with the most massive of them coming close to one million K, though for a short while. Their contribution to the UV and soft X-ray part of the spectrum is therefore of potential importance, especially in old populations.
The broad scenario sketched above for the post-AGB evolution has now gathered fairly general consensus - at least for most of its aspects. This picture has gradually emerged thanks to an extended set of stellar evolutionary calculations, starting with the pioneering work by Paczyลski (1971), and continuing with a series of other relevant studies (which include: Hรคrm & Schwarzschild 1975; Schรถnberner 1979, 1983; Iben et al. 1983; Iben 1984, 1985, 1987; Iben & Tutukov 1984; Iben & MacDonald 1985, 1986; Wood & Faulkner 1986; Vassiliadis & Wood 1994; and Blรถcker 1995). Still on the theoretical side, several of the evolutionary transformations mentioned above have been first sketched in a scattered series of papers (Renzini 1979, 1981a, 1981b, 1982, 1983, 1989, 1990, hereafter P1 through P7), along with a number of conceptual tools which may be useful for their understanding. This paper is now an attempt at presenting these ideas in an orderly and systematic form. To do so we have proceeded with the construction of simulated post-AGB populations, as this technique offers an opportunity to see in a rather straightforward way what are the observable effects of the various theoretical ingredients, and therefore to estimate the relative uncertainties when - proceeding in the opposite direction - one attempts to use evolutionary theory to infer astrophysical quantities from the observational data.
The paper is organized as follows: In ยง2 we describe our procedure for the construction of post-AGB synthetic populations, in a way that makes them easily reproducible. We therein also discuss the role of the parameters that are at variation in our population synthesis. In ยง3 we present our results: synthetic post-AGB populations are plotted within the HR diagrams, and $`M_\mathrm{V}`$-age ( $`M_\mathrm{V}`$-PN radius) distributions; mass distributions are also shown. Our conclusions are summarized in ยง4, where we also draft future applications. In this paper we deal with disk populations with solar composition. Future work will extend the parameter space to other populations as well.
## 2 Synthetic post-AGB populations: ingredients
### 2.1 Interpolation of the evolutionary tracks
The stellar models used in this paper are interpolation of Vassiliadis and Wood (1994, hereafter VW)โs hydrogen burning evolutionary tracks, calculated for post-AGB masses between 0.569 and 0.9 $`M_{}`$, and solar composition (from here onward, M indicates the post-AGB stellar mass). To extend the baseline of our population synthesis, we extrapolate to lower and higher masses by using, as templates, the log $`T_{\mathrm{eff}}`$-log $`L`$tracks for $`M=0.546`$ and 1.2 $`M_{}`$ respectively from Schรถnberner (1983) and Paczรญynski (1971). We choose the VW database for its homogeneity and wide mass range. The VW data set includes also some lower metallicity tracks, which will be used in a future paper for simulating the Magellanic Cloud post-AGB populations.
In Figure 1 and 2 we plot effective temperature and luminosity versus time, as read directly from the VW tracks<sup>1</sup><sup>1</sup>1In the VW database, all log $`T_{\mathrm{eff}}`$-log $`t`$evolutionary tracks start at t=0, corresponding to log $`T_{\mathrm{eff}}`$=10,000 K. The curves in Fig. 1 start from the second evolutionary data point for clarity.. These tracks do not support straightforward interpolation, since they cross one another in several points. We divide each track into parts roughly representing physical phases, and then we interpolate the appropriate normalized functions for each phase.
The effective temperature-age curves have been interpolated as follows. The six log $`T_{\mathrm{eff}}`$-log $`t`$tracks have been divided into four parts, which hereafter are called (a) H-burning phase, (b) quenching phase, (c) cooling phase, and (d) white dwarf phase. In Table 1 we list the characteristics of the VW models. Column (1) give the phase, column (2) the mass of the evolutionary models, column (3) the characteristic time of each phase and mass. Note that each phase starts at the end of the previous one, and that the H-burning phase starts at t=0. Table 1 also indicates the values of log $`T_{\mathrm{eff}}`$and log $`L`$at the end of each phase, in columns (4) and (5). From here after, the capital letters H, Q, C, and W will flag the physical variables at the end respectively of the hydrogen burning, the quenching, the cooling, and the white dwarf phases (e.g., $`t_\mathrm{H}`$(0.569 $`M_{}`$)=$`3.215\times 10^4`$ yr). For each phase, which correspond to a well defined age interval, we have studied the best way to normalize the log $`T_{\mathrm{eff}}`$-log $`t`$function to eliminate the crossing of the tracks of contiguous masses.
In Figure 3 we show the temperature evolution on the hydrogen burning phase.<sup>2</sup><sup>2</sup>2It is understood that all temperatures are in Kelvins, all masses in $`M_{}`$, the stellar luminosities in $`L_{}`$, and the ages in years, unless otherwise noted. Here log $`T_{\mathrm{eff}}`$has been plotted against the normalized age
$$t_{\mathrm{nn}}=t/t_\mathrm{H};$$
$`(1)`$
the dots represent the normalized evolutionary data points. The interpolation of the normalized function is straightforward in this phase. Similarly, Figures 4, 5, and 6 show respectively the log $`T_{\mathrm{eff}}`$curves for the quenching, the cooling, and the white dwarf phases. In the quenching phase (Fig. 4), it is possible to interpolate directly log $`T_{\mathrm{eff}}`$versus $`t_{\mathrm{nn}}`$. The normalized age in the quenching phase is:
$$t_{\mathrm{nn}}=(tt_\mathrm{H})/(t_\mathrm{Q}t_\mathrm{H}).$$
$`(2)`$
In the cooling phase (Fig. 5), we use the interpolating function:
$$T_{\mathrm{nn}}=\mathrm{log}T_{\mathrm{eff}}/\mathrm{log}T_{\mathrm{eff},\mathrm{Q}}$$
$`(3)`$
with normalized age:
$$t_{\mathrm{nn}}=(tt_\mathrm{Q})/(t_\mathrm{C}t_\mathrm{Q}).$$
$`(4)`$
In the last branch of the evolutionary track, the white dwarf phase (Fig. 6), we normalize the effective temperature as:
$$T_{\mathrm{nn}}=\mathrm{log}T_{\mathrm{eff}}\mathrm{log}T_{\mathrm{eff},\mathrm{C}},$$
$`(5)`$
where the normalized age assumes the values of:
$$t_{\mathrm{nn}}=(tt_\mathrm{C})/(t_\mathrm{W}t_\mathrm{C})$$
$`(6)`$
(note that $`T_{\mathrm{nn}}`$ of Eqs. 3 and 5 do not correspond to an actual model temperature).
Let us explore the luminosity tracks, for which we use a a similar interpolation approach. The luminosities against normalized ages are plotted, for phase H, in Figure 7. The luminosity interpolation of this phase is evidently straightforward. The quenching phase is very well interpolated by a straight line, calibrated for each mass as
$$L_{\mathrm{nn}}=A+B[(tt_\mathrm{H})/(t_\mathrm{Q}t_\mathrm{H})],$$
$`(7)`$
with $`A=0.041`$ and $`B=1.02`$. The cooling phase normalized luminosity is plotted in Figure 8 for tracks corresponding to masses between 0.569 and 0.754 $`M_{}`$. For larger masses, a better interpolation is achieved by using the power function, with exponent $`\alpha `$=0.23 for $`0.754M0.8`$ $`M_{}`$, and $`\alpha `$=0.26 for $`M0.8`$ $`M_{}`$. Figure 9 shows the white dwarf phase of the luminosity (vs. time) curves, where log $`L`$has been plotted against the normalized time,
$$t_{\mathrm{nn}}=(tt_\mathrm{C})/(t_\mathrm{W}t_\mathrm{C}).$$
$`(8)`$
Even in this case, the interpolation is straightforward.
After breaking the log $`T_{\mathrm{eff}}`$and log $`L`$tracks into phases, and constructing the normalized functions as described above, our procedure includes the following steps, for each given mass and age. (a) We identify the mass-interval to be used. For example, if $`M=0.8`$$`M_{}`$, we will use the two enclosing tracks for the interpolation, corresponding to $`0.754`$ and $`0.9`$ $`M_{}`$. (b) We identify the phase of the evolution to consider, given the age and mass. To do so, we evaluate $`t_\mathrm{H}`$, $`t_\mathrm{Q}`$, $`t_\mathrm{C}`$, $`t_\mathrm{W}`$ for the given mass. Let us suppose, in our example, that $`t=400`$ years. We obtain respectively 775, 859, and 29,689 yr for the H, Q, and C characteristic times. This means that at $`t=400`$ our 0.8 $`M_{}`$ star is in the H-burning phase. (c) We calculate the normalized time, $`t_{\mathrm{nn}}`$. In our example, $`t_{\mathrm{nn}}=0.52`$. (d) On the appropriate tracks (in our example, Figs. 3 and 7) we read off log $`T_{\mathrm{eff}}`$and log $`L`$in the appropriate range of $`t`$ and $`M`$. In our example, in Fig. 3 we draw the vertical line $`x=0.52`$ and read off the log $`T_{\mathrm{eff}}`$values from the two upper tracks, the four points immediately before and after the vertical line. (c) We interpolate, on the higher and lower mass tracks, and find the log $`T_{\mathrm{eff}}`$and log $`L`$that intersect the $`x=t_{\mathrm{nn}}`$ vertical line; (4) we interpolate vertically versus log M, where M is the given mass.
When evaluated on the masses of the original VW evolutionary tracks, the synthetic temperature and luminosity curves are indistinguishable from those plotted in Figures 1 and 2. We have studied the interpolation offsets in detail. The relative errors in the linear temperatures, (T<sub>eff,syn</sub>-T<sub>eff,VW</sub>)/T<sub>eff,VW</sub> <sup>3</sup><sup>3</sup>3 The subscripts syn and VW refer respectively to the synthetic and the VW evolutionary tracks., result to be collimated between -0.005 and 0.005 in most evolutionary phases, with the exception of the very fast quenching phase, where the linear errors on the temperature are rather between -0.02 and 0.02 for the lower 5 mass tracks, and between -0.02 and 0.1 for the 0.9 track. The relative errors in the luminosities, (L<sub>syn</sub>-L<sub>VW</sub>)/L<sub>VW</sub>, are between -0.01 and 0.01 in most evolutionary phases, with the exception of the quenching phase, where the linear errors on the luminosity are between -0.15 and 0.15 for the lower 5 mass tracks, and between -0.15 and 0.7 for the 0.9 track.
When using the interpolated values au lieu of the evolutionary ones, we obtain a perfect substitution. In the quenching phase, the substitution would displace the temperatures by a marginal fraction, and the misplacement of the luminosity curves would propagate only in a few percent error on log L. In all other cases, the substitution reflects in virtually no errors. We conclude that the interpolations are excellent and can be reliably used as evolutionary tracks for all stellar masses between 0.569 and 0.9 $`M_{}`$.
We obtain extrapolations to higher and lower masses using the log $`T_{\mathrm{eff}}`$-log $`L`$templates of 0.546 and 1.2 solar masses from Schรถnberner (1983) and Paczรฌnski (1971). Only the shape of the log $`T_{\mathrm{eff}}`$-log $`L`$relation is used here, not the actual evolution of the physical parameters, which are instead extrapolated directly from the 0.596 and 0.9 $`M_{}`$ tracks by VW, for homogeneity.
### 2.2 The initial mass function
The synthetic distribution of post-AGB stars (e.g., on the HR diagram) depends on the adopted IMF and history of star formation. Note that a given distribution of post-AGB stars obtained with a certain IMF and star formation (SR) history can be reproduced also with a different IMF provided a properly tuned SF history is adopted. We have renounced to play with two independent parameters - such as the slope of the IMF and the e-folding time of the SF rate - and have exclusively explored the effect of changing slope of the actual initial mass distribution $`\psi (M_\mathrm{i})M_\mathrm{i}^{(1+x)}`$. Therefore $`x`$ is effectively the slope of the IMF only if a constant star formation rate is assumed. This may not be a bad approximation for the Galactic Disk (Scalo 1998). In our simulations it is supposed that only stars in the mass range $`0.85M_\mathrm{i}<9M_{}`$ experience the thermal pulses on the AGB (TP-AGB) and the post-AGB phases. The upper limit of the mass domain has been chosen in the light of the carbon ignition limits described in Iben (1995). We include in our investigation the classic Salpeterโs (1955) IMF, with constant index for each mass range. Furthermore, we inspect the effects of a IMF with variable exponent, such as the one by Miller & Scalo (1979), the low-mass extension by Kroupa et al. (1991), and the IMF recently proposed by Scalo (1998). In all three cases, the function index changes for masses smaller or larger than $`M_\mathrm{i}`$= 1$`M_{}`$.
### 2.3 The initial mass-final mass relation
The initial mass-final mass relation (IMFMR) may have a major impact on the resulting synthetic distributions. To illustrate the case we have constructed simulations adopting four different IMFMRs, either from theoretical calculations or from observations: (1) the old theoretical Renzini & Voli (1981) relation as analytically approximated by IR <sup>4</sup><sup>4</sup>4In this paper we flag this IMFMR with โIRโ.; (2) a more recent theoretical IMFMR introduced by Ciotti et al. (1991), which is an attempt at incorporating new important results of evolutionary calculations; (3) the classic empirical IMFMR proposed by Weidemann (1987) on the basis of WD masses observed in open clusters; and (4) a new empirical relation presented by Herwig (1997) that includes new observations of cluster white dwarfs in the Pleiades, the Hyades, and NGC 3451. Figure 10 shows the four IMFMRs (note that in this paper we use $`M`$=$`M_\mathrm{f}`$).
A brief justification for the most recent theoretical relation (Ciotti et al. 1991) is appropriate. The old theoretical IMFMR (IR) was obtained interpolating on an insufficient grid of stellar models, and assumed a universal core mass-luminosity relation. Lattanzio (1989) pointed out that the straight line approximation for the core mass at the first AGB thermal pulse (as a function of $`M_\mathrm{i}`$, and for $`M_\mathrm{i}<3M_{}`$) was much steeper than resulting from actual evolutionary calculations. Furthermore, Blรถcker & Schรถnberner (1991) have found a major breakdown of the core mass-luminosity relation (a key ingredient in theoretical IMFMRs) for those stellar models in which the so-called envelope burning process is activated, i.e., for $`M_\mathrm{i}3M_{}`$ (cf. Renzini & Voli 1981). In practice the former effect results in a flattening the IMFMR below $`3M_{}`$, the second one in a flattening of the IMFMR above $`3M_{}`$. The new theoretical IMFMR incorporates these findings, while keeping the same parameterization of mass loss processes (wind and superwind) as adopted by Renzini & Voli, with $`\eta =0.5`$ and $`b=1`$.
On the empirical side, the new relation by Herwig (1997) is not very different than the new theoretical relation. In addition, we should mention that new data points by Jeffries (1997) of white dwarfs in NGC 2516 also lie very close to Herwigโs relation. Furthermore, it should be noted that more recent calculations presented in conferences (e.g., Lattanzio & Forestini 1999; Blรถcker 1999) show that the envelope burning occurs for M$`>`$4 $`M_{}`$, bringing the theoretical MIMFR hardly distinguishable from Herwigโs (1997).
### 2.4 The transition time and the post-AGB envelope mass
At the TP-AGB phase, the red giant envelope is ejected by the superwind. The subsequent evolution of the star to PNN depends in a substantial way on the amount of envelope mass left on the star at the SW quenching, $`M_\mathrm{e}^\mathrm{R}`$. No matter how the transition between the AGB and the PN illumination occurs, the remnant envelope mass at the end of the SW plays a central role. This is just unfortunate when trying to simulate this evolutionary phase: in fact, $`M_\mathrm{e}^\mathrm{R}`$ is not defined by stellar evolution (given the hydrodynamical nature of the superwind), and only approximations or guesses can be made about its entity and its possible relation to the physical parameters. $`M_\mathrm{e}^\mathrm{R}`$ can indeed be considered a free parameter, in the sense that there is no theoretical nor observational constraint to fit it, nor to give it an exact dependence on any nebular or stellar parameter.
In principle, the higher the core mass, the higher the stellar luminosity, and the lower the envelope mass left on the AGB star. But this also depends on the thermal pulses that occur at the TP-AGB, and it is hard to know from first principle how the superwind ejection changes the stellar structure. We examine later in this section how we choose $`M_\mathrm{e}^\mathrm{R}`$ for our simulation.
To understand the transition between the AGB and the PN phases, we should introduce the timescales involved in the different phases. Immediately after the superwind quenching, the mass loss due to stellar wind dominates, until the star detaches itself from the AGB. The time scale at which this phase occurs is the wind time scale, and can be written as:
$$t_\mathrm{w}=\frac{\mathrm{\Delta }M_\mathrm{e}}{\dot{M}},$$
$`(9)`$
where $`\mathrm{\Delta }M_\mathrm{e}`$ is the difference between the residual envelope mass at the superwind quenching, $`M_\mathrm{e}^\mathrm{R}`$, and the envelope mass at the detachment from the AGB (i.e., at a later phase), $`M_\mathrm{e}^\mathrm{D}`$. The wind mass loss rate (MLR) is taken from Reimers (1975), evaluated at the AGB temperature of 5000 K. Reimersโ approximation seems reasonable in the considered evolutionary phase, that is, after the superwind quenching and before the onset of the radiation-driven wind (see Blรถcker 1995). In effect, the MLR right after the detachment from the AGB declines as the stellar temperature increases (Blรถcker 1995); but in the present application, i.e., to evaluate the transition time, this choice does not affect the results.
For larger temperatures, in the blueward runaway, the mass consumption occurs at a nuclear timescale, following the Equation:
$$t_\mathrm{n}=\frac{X_\mathrm{e}\mathrm{\Delta }M_\mathrm{e}E_\mathrm{H}}{L}.$$
$`(10)`$
$`\mathrm{\Delta }M_\mathrm{e}`$ is now the difference between the smaller among $`M_\mathrm{e}^\mathrm{R}`$ and $`M_\mathrm{e}^\mathrm{D}`$, and the amount of envelope mass remaining at the illumination of the nebula, $`M_\mathrm{e}^\mathrm{N}`$; the luminosity in the equation is the plateau luminosity <sup>5</sup><sup>5</sup>5The plateau luminosity, i.e., the luminosity in the early post-AGB phase, is $`L_{\mathrm{PL}}=56694\times (M0.5)`$ for VW models.; and the variables $`X_\mathrm{e}`$ and $`E_\mathrm{H}`$ are the envelope hydrogen abundance and the energy released by the nuclear burning of one gram of hydrogen, respectively.
Finally, the transition time will scale with the shorter among the wind and the nuclear time scales. We set the reciprocal of the transition time to be the sum of the reciprocals of the wind and the nuclear times.
The transition occurs on a thermal timescale in the case in which $`M_\mathrm{e}^\mathrm{R}`$$``$ $`M_\mathrm{e}^\mathrm{N}`$, thus in the case in which either the nuclear or the wind time is zero. The formula to be used then is:
$$t_{\mathrm{th}}=\frac{GMM_\mathrm{e}^\mathrm{R}}{LR},$$
$`(11)`$
where L and R are the stellar luminosity and radius (see also P6).
It is worth recalling that the synthetic tracks obtained in $`\mathrm{\S }`$2.1 have their zero age points at log $`T_{\mathrm{eff}}`$=4.0, which are approximately reached when the transition has been completed. At 10,000 K, the star is able to ionize hydrogen, although the complete nebular transparency at the optical wavelengths occurs generally at slightly higher temperatures (Kรคufl et al. 1993). In all events, we use the synthetic tracks from their natural zero point.
Evolutionary models for post-AGB stars predict the amount of mass available in the stellar envelope as a function of the effective temperature (Schรถnberner 1983, Paczynski 1971). We use Schรถnbernerโs models (M$`<`$ 0.65) to obtain $`M_\mathrm{e}^\mathrm{D}`$(envelope mass at the detachment) and $`M_\mathrm{e}^\mathrm{N}`$(envelope mass at nebular illumination, when T<sub>eff</sub>=10,000 K) as functions of M, by linear interpolation. For M$`>`$0.65 $`M_{}`$ we obtain $`M_\mathrm{e}^\mathrm{D}`$ and $`M_\mathrm{e}^\mathrm{N}`$ by scaling Paczynskiโs 1.2 $`M_{}`$ model to Schรถnbernerโs values.
As discussed at the beginning of this section, $`M_\mathrm{e}^\mathrm{R}`$ is quite unconstrainted, and its value affords unpredictable variations. We set ourselves to explore a wide parameter space with three different assumptions for $`M_\mathrm{e}^\mathrm{R}`$.
First, since the evolutionary models have shown that the envelope mass at several evolutionary phases in the post-AGB is inversly proportional to the core mass, we could envision the possibility that also at the superwind quenching $`M_\mathrm{e}^\mathrm{R}`$ has an inverse dependence on M. We parametrize this case following the usual models as guidelines, by taking $`M_\mathrm{e}^\mathrm{R}`$ as the envelope mass at the onset of the post-AGB models. Figure 11 illustrates how the characteristic timescales vary if $`M_\mathrm{e}^\mathrm{R}`$=f(M). The transition time peaks for the very small mass models, and then it declines for larger masses.
Second, we chose a constant $`M_\mathrm{e}^\mathrm{R}`$, independent of mass. In Figure 12 and 13 we analyze the consequences of assuming two values of $`M_\mathrm{e}^\mathrm{R}`$, $`M_\mathrm{e}^\mathrm{R}`$=3$`\times 10^4`$ and $`M_\mathrm{e}^\mathrm{R}`$=5$`\times 10^3`$. In the first case (Fig.12), the transition time follows up close the thermal time scale, then the nuclear time scale for m$`>`$0.85 $`M_{}`$, to decline for larger masses. In the second case (Fig. 13), the transition time rises to almost 6000 yr for M$``$0.55 $`M_{}`$. These two cases with constant $`M_\mathrm{e}^\mathrm{R}`$ are illustrated only for the sake of showing extreme cases, but it is very unrealistic that all stars end the unstable SW phase with exactly equal $`M_\mathrm{e}^\mathrm{R}`$.
Third, we use random values of $`M_\mathrm{e}^\mathrm{R}`$ for our calculations. In Figure 14 we show the transition, wind, and nuclear time scales for random $`M_\mathrm{e}^\mathrm{R}`$, with maximum equal to $`M_\mathrm{e}^\mathrm{R}`$=$`5\times 10^3`$. The transition time peaks at about 0.6 solar masses, then declines. The same simulation with different maximum $`M_\mathrm{e}^\mathrm{R}`$ gives similar results, except for the vertical scale. Naturally, giving the randomness of the extraction, the population plotted in Figure 14 is one of the infinite possible extractions with random $`M_\mathrm{e}^\mathrm{R}`$ lesser than $`5\times 10^3`$ $`M_{}`$.
### 2.5 The final helium-shell flash
During the post-AGB phase hydrogen is burned in a shell, and therefore the mass of the helium buffer zone between the C-O core and the hydrogen shell keeps increasing. In some cases the increase of the buffer mass can lead to a last thermal pulse (also called flash) of the helium burning shell (Schรถnberner 1979). This happens if the star left the AGB with a sufficiently massive buffer zone, so that its further increase during the post-AGB allows to reach the critical value for a last flash to erupt. This chance is therefore related to the phase in the AGB thermal pulse cycles at which the SW envelope ejection takes place. Detailed model calculations show that as a response to the final flash stars undergo an extended loop in the HR diagram. The model star would expand back to the AGB very rapidly, then it would evolve again as a post-AGB star, powered by helium burning; the duration of the He-burning phase is about three times longer than the fading time of hydrogen-burning post-AGB stars (Iben 1984). These calculations also show that during the power down phase the rate of luminosity decline is nearly constant, as opposed to the case of hydrogen burning post-AGB stars in which an abrupt drop follows the plateau phase.
To incorporate the effect of the final helium-shell flash (FF) we have proceeded as follows: we assume that such stars return to $`(L,T_{\mathrm{eff}})=(L_{\mathrm{PL}},10,000`$ K) upon a final flash, and their luminosity fades then linearly with time to $`L_\mathrm{F}`$ in a time $`3\times t_\mathrm{F}`$ (where $`t_\mathrm{F}`$ is the total fading time for the H-burning star, see $`\mathrm{\S }`$ 2.7.2), and $`L_\mathrm{F}`$ is the luminosity of the corresponding H-burner, i.e.:
$$L=L_{\mathrm{PL}}\frac{L_{\mathrm{PL}}L_\mathrm{F}}{3\times t_\mathrm{F}}t^{}.$$
$`(12)`$
We then locate the star on the log $`L`$\- log $`T_{\mathrm{eff}}`$curves as if the evolution were three times as slow as H-powered evolution. This provides a fairly good approximation of the behavior of the models constructed by Iben (1984). The fraction of stars experiencing a final flash is not strictly determined by theory. Iben has calculated the probability that a star ignites helium in various post-AGB phases (Table 2, Iben 1984), obtaining different guesses for FF occurrence in 10 to 21 $`\%`$ of all stars in this mass range, and up to about 60 $`\%`$ when he includes stars that leave the AGB burning helium. Our parametrization is comparable with Ibenโs prediction. In fact, if we were to chose that 60 percent of all stars (in this mass range) have the chance to experience a final flash, then about 12 $`\%`$ of the PNNi are in He-burning phase as observed in a synthetic HR diagram.
### 2.6 The duration of the planetary nebula phase
Following common wisdom, the duration of the PN phase is $`30,000`$ yr (e.g., Phillips 1989, and references therein). This is derived from the size of the largest observed PNs ($`0.7`$ pc) coupled with a typical nebular expansion velocity of $``$25 km s<sup>-1</sup>. In practice this assumes that all PNs remain visible as such until the expanding nebula has reached the maximum observed size, or, equivalently, that all PNs evolve in nearly the same way. On the other hand, the mere fact that PNs are produced by precursor stars in such a very extended range of initial masses ($`0.85M_\mathrm{i}9M_{}`$) makes most unlikely that all PNs have the same lifetime. Simple arguments suggest that the time $`t_{\mathrm{max}}`$ after the cessation of the superwind during which a PN remains observable scales as $`t_{\mathrm{max}}M_{\mathrm{PN}}^{2/5}SB_{\mathrm{min}}^{1/5}`$, where $`M_{\mathrm{PN}}`$ is the nebular mass and $`SB_{\mathrm{min}}`$ is the minimum surface brightness for a PN to be detected (P2,P5). In the adopted mass loss parameterization, $`M_{\mathrm{PN}}`$ corresponds to the mass ejected during the AGB superwind phase, and ranges from $`0.02M_{}`$ for $`M_\mathrm{i}=0.85M_{}`$, to over $`1M_{}`$ for $`M_\mathrm{i}=8M_{}`$. This factor of $`50`$ in PN mass therefore translates into at least a factor $`5`$ in $`t_{\mathrm{max}}`$. To get the actual duration of the PN phase we must subtract to $`t_{\mathrm{max}}`$ the AGB to PN transition time, during which the object is a proto-PN, i.e., $`t_{\mathrm{PN}}=t_{\mathrm{max}}t_{\mathrm{tr}}`$. If $`t_{\mathrm{tr}}>t_{\mathrm{max}}`$, then when the central star has reached 10,000 K the nebular material has already dispersed, no observable PN is produced, and one has a so-called lazy post-AGB remnant (P2). To explore the effects of this subtle interplay between central star and (partially decoupled) nebular evolution, we have explored the case of mass dependent $`t_{\mathrm{max}}`$, scaled with $`M_{\mathrm{PN}}`$. The nebular mass was derived by imposing its dependence on the plateau luminosity, since it is at the reach of a critical luminosity that the shell is ejected. The parametrization of $`M_{\mathrm{PN}}`$ and $`t_{\mathrm{max}}`$ are described in the Appendix (Eqs. A1 and A2). Simulations of post-AGB populations with constant maximum PN ages will be also shown in this paper.
### 2.7 The Montecarlo procedure
Our population synthesis code starts with the option for a time limited or luminosity limited sample. The former option is to explore post-AGB populations within a fixed time range, (we are talking here about the stellar evolutionary time, not the duration of the PN phase, described in $`\mathrm{\S }`$2.6). This option is used to deduce the mass distribution of the PNN as derived from the synthetic population analysis. The latter option, the luminosity-limited sample, is used to compare directly synthetic and observed diagrams. All stars in the synthetic population would have L$`L_\mathrm{F}`$ (we assume $`L_\mathrm{F}`$=1.0 $`L_{}`$)<sup>6</sup><sup>6</sup>6Naturally, the luminosity limit can be varied accordingly to the observed situation we may want to reproduce..
To build the synthetic population, we proceed as follows:
#### 2.7.1 Time limited sample
1) First, a value of the initial mass $`M_\mathrm{i}`$ is extracted, following the distribution:
$$\psi (M_\mathrm{i})M_\mathrm{i}^{(1+x)}$$
$`(13)`$
2) The extracted value of $`M_\mathrm{i}`$ is entered into the initial mass-final mass relation to get the mass $`M`$ of the post-AGB object.
3) To proceed, we need to evaluate $`M_\mathrm{e}^\mathrm{R}`$ (see $`\mathrm{\S }2.4`$), and correspondingly the wind time scale can be calculated (see Eq. 9).
4) A random time is then extracted, within a chosen range ($`0tt_{\mathrm{lim}}`$).
5) If the extracted time $`t`$ is larger than the wind timescale, we put $`t=tt_\mathrm{w}`$, and this value of $`t`$ is entered into the routines of $`\mathrm{\S }`$2.1 thus getting the corresponding location in the HR diagram. If $`t`$ is smaller than $`t_\mathrm{w}`$ the star is still in the wind phase, and we assign $`T_\mathrm{e}`$=3.7 and $`L=L_{\mathrm{PL}}`$. We call the resulting simulated star a wind object<sup>7</sup><sup>7</sup>7 Given our assumption on the MLR, we distinguish between wind objects and proto-PNN (in , below): wind objects are those transition objects whose evolution is wind dominated, proto-PNNi are all the other transition objects, located beyond the AGB. The assigned temperature for the wind object is rather arbitrary, but homogeneous to our mass loss choices. Even if we were to locate the wind objects within the transition area, the results would be not majorly affected, given the paucity of the samples. In a future paper, different mass loss choices will analyze the possible differences. Observationally, it may be hard to distinguish between these two types of objects..
If we do not consider a FF, we go to point (8).
6) Should we consider the effects of a final helium flash, another value of time $`t_{\mathrm{FF}}`$ is extracted, with $`0t_{\mathrm{FF}}t_{\mathrm{lim}}`$, so as to get a random value for the time at which the flash takes place <sup>8</sup><sup>8</sup>8Note that Iben (1984) has calculated slightly different probabilities for FF occurrence at different phases, but we use a flat probability for simplicity..
7) If $`t<t_{\mathrm{FF}}`$, then this $`t`$ is entered into the routines of $`\mathrm{\S }2.1`$ to get the luminosity and temperature of the star (in this case we have extracted a post-AGB star which will experience a final flash, but the flash has not taken place yet).
If $`t>t_{\mathrm{FF}}`$, then the time elapsed after the the flash ($`t^{}=tt_{\mathrm{FF}}`$) is entered into Equation (12) to get luminosity of the post-FF star. To calculate the temperature we use the routines of $`\mathrm{\S }2.1`$.
8) The transition time $`t_{\mathrm{tr}}`$ is obtained following the prescriptions in ยง2.4.
9) The object is finally classified into one of the following classes, and plotted with a different symbol for each class. The object is classified as a proto-Planetary Nebula Nucleus (proto-PNN) if $`tt_{\mathrm{tr}}`$. It is classified as a PNN if $`t+t_{\mathrm{tr}}t_{\mathrm{max}}`$, or $`t^{}+t_{\mathrm{FF}}+t_{\mathrm{tr}}t_{\mathrm{max}}`$, with the further distinction between objects that have experienced a final flash (PNNHe) and those which did not (PNNH). It is classified as a post-PNN object otherwise, again distinguishing post FF and no post FF objects.
#### 2.7.2 Luminosity limited sample
To build a luminosity-limited sample, we follow a similar procedure as in the time-limited sample, except in (4), the extraction of a random time, that is chosen in the range $`0tt_\mathrm{F}(M)+t_\mathrm{w}`$, where $`t_\mathrm{F}`$ is the fading time to the minimum luminosity, $`L=1.0L_{}`$ (see A3).
The we proceed as in points (5) through (8) of the procedure in $`\mathrm{\S }`$ 2.7.1, with the difference that $`t_{\mathrm{FF}}`$ is extracted within the new time limits.
## 3 Results
Our code is flexible, suitable to produce post-AGB population synthesis for many applications and many sets of parameters. Far from being exhaustive, this paper includes only a small part of the possible applications. We have explored the issues and questions that we though to be of wide interests, and allow to advance in out post-AGB evolution knowledge. Other applications will be developed in the future.
### 3.1 Synthetic diagrams and mass histograms: the basic model
The first synthetic population shown here is a luminosity-limited sample of post-AGB stars. We extract 1500 objects, and we separate the wind objects, the PNNi, the proto-PNNi, and the post-PNNi following the prescriptions of Chapter 2. In Figure 15 we show the synthetic HR diagram, where 313 objects of the 1500 extracted are in the PNN phase, 21 are in the wind phase, and 9 are proto-PNNi. The remaining objects are, following our prescription, already in the post-PNN phase.
In the simulation of Figure 15 the residual envelope mass, $`M_\mathrm{e}^\mathrm{R}`$, is chosen to be function of the core mass, as described in $`\mathrm{\S }`$2.4; we have used Weidemannโs IMFMR, Salpeterโs IMF, and no FF objects (i.e., all PNN are hydrogen burners). The maximum PN time has been set following Eq. A2, with K=4$`\times 10^5`$, a relatively long time for PNe to disappear. We use this time through $`\mathrm{\S }`$3.5, considering furtherly the effects of the maximum PN age in $`\mathrm{\S }`$3.6. We define the simulation plotted in Figure 15 as basic, and we will explore the parameter space by changing one parameter at the time with respect to this basic population.
As expected, the synthetic stars cluster toward the lower masses, only the very low luminosity part being populated by higher masses. This is exactly the effect we see in all complete galactic PNNi samples.
For more insight, let us plot the corresponding distribution of the visual magnitudes versus evolutionary timescale (Fig. 16). We translate temperatures and luminosities into visual magnitudes by following the bolometric correction from Code et al. (1976). As already evident in the previous plot, there are two mainly populated loci of the $`M_\mathrm{V}t`$ plane as well, corresponding one to plateau luminosity of the low-mass objects, and the other to the stellar crowding toward low luminosities of the intermediate- and high-mass objects. The large fraction of PNNi that have decayed to post-PNNi at late evolutionary times is very evident in this Figure. This effect is mass-dependent, and smaller mass PNN disappear earlier in their life.
In Figure 17 we show the (logarithmic) mass distribution for the basic synthetic population. This distribution has been obtained by running a sample of 1500 post-AGB stars with maximum HR diagram-life time of 30,000 years. The logarithmic scale has been chosen for a better view of the distribution in the whole mass range. The mass distribution of Figure 17 shows the clear clustering of model post-AGB stars around $``$0.6 $`M_{}`$, and a gradual spread to higher masses.
Due to the particular choice of the time interval explored ($`t<30,000yr`$), we do not obtain post-PNNi with this time-limited selection. In Table 2 we show the composition of a selection of the synthetic populations shown in this paper. The columns list, respectively, the number of PNNi, wind objects, post-PNNi, and proto-PNNi stars for each extraction. Composition of typical populations may vary from one random extraction to another. From Table 2 we can directly compare the time-limited to the luminosity-limited samples.
### 3.2 Effects of the initial mass function
To test the variations of post-AGB distribution for different Initial Mass Function we have run synthetic populations with the following IMFs (see $`\mathrm{\S }2.2`$ for meaning of variable x and for SF choice): (1) the standard, Salpeterโs (1955) IMF, with $`(1+x)=2.35`$ in the whole mass range considered; (2) Miller & Scaloโs (1979) IMF, with $`(1+x)=1.4`$ for $`M<1`$$`M_{}`$ and $`(1+x)=2.5`$ for $`M>1`$ $`M_{}`$; (3) the IMF by Kroupa et al. (1991), with $`(1+x)=1.85`$ for $`M<1`$ $`M_{}`$ (this distribution is not defined for larger masses), and (4) an updated empirical parametrization by Scalo (1998) with $`(1+x)=1.2`$ for $`M<1`$$`M_{}`$ and $`(1+x)=2.7`$ for $`M>1`$$`M_{}`$.
In Figure 18 we show HR distributions accounting for IMF variations, when the IMF is the only parameters changing across the panels. In each case, we have extracted a luminosity-limited sample of 1500 post-AGB stars, with Weidemannโs IMFMR, and mass-dependent residual envelope mass. The stars do not experience a final helium shell flash. In the case of Kroupa et al. โs low mass IMF, we use Scaloโs (1998) IMF for masses larger than solar.
By examining Figure 18 <sup>9</sup><sup>9</sup>9In Fig. 18, and after, we use the following conventional nomenclature for figure panels: a=lower left, b=upper left, c=lower right, d=upper right. Panel (a) always shows the basic sample. we see very little changes on the log $`L`$\- log $`T_{\mathrm{eff}}`$plane for different IMF in the mass range considered. The PNN populations in the case of variable IMF index (Fig. 18bcd) are larger for low masses respect to the basic model. The proto-PNN population changes slightly (we are dealing with low number statistics, anyway), but overall the distributions are similar in all four cases of Figure 18. Better to say, that the HR diagram is not the ideal way to pick up these differences, occurring in this case at low-luminosities, in a crowded part of the diagram. Observations of real situations such as those of the four panels of Figure 18 would not be distinguishable from one another.
In Figure 19 we illustrate the mass distribution for the synthetic populations corresponding to Figure 18, each panel showing a different IMF. Naturally, time-limited samples have been used. It is worth noting that peaks in the distribution of less than $`\sigma `$ can be produced by the randomness of the simulation.
The effects of the IMF are noticeable given the large sample of model stars. For example, Scaloโs IMF (Fig. 19d) would produce a broader distribution toward the higher masses. Note that the logarithmic scale partially hides the fact that all these mass distributions are extremely peaked around 0.55 $`M_{}`$.
### 3.3 Effects of the initial mass-final mass relation
As introduced in $`\mathrm{\S }2.3`$, we will use four different IMFMR and test their effects on the post-AGB populations. In Figures 20 and 21 we plot the results relative to (a) Weidemannโs, (b) Ciotti et al. , (c) Herwigโs, and (d) IRโs IMFMRs. Figure 20 shows the distributions on the HR diagram. On this plane, the sample with IRโs IMFMR stands aside, since the IRโs prescription produces more massive post-AGB.
The differences among the other distributions are more evident in Figure 21, where we plot the visual magnitude versus evolutionary time of the synthetic stars. By comparing the basic and Herwigโs IMFMR diagrams, above Mv=5 the two distributions are similar, but the differences appear and get more evident for fainter objects.
To compare these results with the actual observed stars is beyond the purpose of this paper. Nonetheless, we can state that the different IMFMRs could be inferred mostly at low brightness, with consequent higher uncertainty of the comparison with the data.
The mass distributions of Figures 22 set aside once again the IR IMFMR. This old parametrization allows a continuum of masses between 0.55 and 1.4 $`M_{}`$. The other 3 cases have similar low-mass distributions, while the high-mass distributions are similar in the Ciotti et al. โs and the Herwigโs IMFMRs, as the basic sample allows for few very high mass stars. The observational consequences of the four different IMFMRs are mainly constrained in the high-mass tail of the distributions.
### 3.4 The residual envelope mass
The envelope mass that is left on the star after the envelope ejection, $`M_\mathrm{e}^\mathrm{R}`$, plays a fundamental role in the following stellar fate. Our basic model uses $`M_\mathrm{e}^\mathrm{R}`$=f(M), with an ad hoc parametrization (see $`\mathrm{\S }2.4`$). Since the value of $`M_\mathrm{e}^\mathrm{R}`$ is undetermined, as well as its distribution with respect to the other physical parameters, we chose to explore scenarios of post-AGB evolution both with with variable and constant $`M_\mathrm{e}^\mathrm{R}`$. We also produce populations with random residual envelope masses, to represent the extreme (but not unrealistic) indetermination of $`M_\mathrm{e}^\mathrm{R}`$, and consequently of the transition time.
In Figure 23 we show the results with the different scenarios: the basic model (a) has mass-dependent $`M_\mathrm{e}^\mathrm{R}`$, model (b) has random $`M_\mathrm{e}^\mathrm{R}`$ (with $`M_\mathrm{e}^\mathrm{R}`$$`<`$0.1 $`M_{}`$), model (c) also has random $`M_\mathrm{e}^\mathrm{R}`$ (with $`M_\mathrm{e}^\mathrm{R}`$$`<1\times 10^4`$ $`M_{}`$), and model (d) has constant $`M_\mathrm{e}^\mathrm{R}`$=$`5\times 10^3`$ $`M_{}`$, independent on the progenitor mass (lazy-AGB stars have not been produced in these simulations, given the choice of K in equation A2).
The effects of the indetermination of the residual envelope mass are striking on the resulting post-AGB populations. A complete description on mass loss for low- and intermediate-mass stars does not exist to date. Our experiments (e.g., Fig. 23) represent a way to determine the effects of the mass dependence on the MLR on observable sets of PNNi and related objects. Once again, in this paper we do not compare these effects with the observed stars, we just set the stage for future comparisons. When using $`M_\mathrm{e}^\mathrm{R}`$=5$`\times 10^3`$$`M_{}`$, the distribution is very similar to the basic case.
We use the random $`M_\mathrm{e}^\mathrm{R}`$ to further analyze the effects of the mass loss on observed populations. The distributions that we obtain as a result of the simulations with random $`M_\mathrm{e}^\mathrm{R}`$ are very different depending on the upper limits that we set for $`M_\mathrm{e}^\mathrm{R}`$. For instance, for $`M_\mathrm{e}^\mathrm{R}`$$`<1\times 10^4`$ $`M_{}`$, no wind nor proto-PN are created. Furthermore, very few high-luminosity PNe are available with this low-limit random $`M_\mathrm{e}^\mathrm{R}`$ (see also Table 2). These effects could be observed in PNNi and proto-PNNi at known distances. At lower luminosities, the two distributions look similar especially if we take into account observing errors.
The mass distributions of post-AGB populations with different choices of $`M_\mathrm{e}^\mathrm{R}`$ are shown in Figure 24. The effects of changing $`M_\mathrm{e}^\mathrm{R}`$ are even more extreme in these time-limited samples. Let us examine Figures 24b and 24c (upper-left and lower-right panels). In both panels, $`M_\mathrm{e}^\mathrm{R}`$ has random values, the difference being the maximum allowed $`M_\mathrm{e}^\mathrm{R}`$. In the simulation of 24b, where the maximum residual envelope mass is 0.1 $`M_{}`$, the distributions is dominated by the objects in the wind phase, and the lack of PNNi (actually, the simulation produces two PNNi, hiddent in the log-representation). The most notable characteristics of panel 24c are the lack of wind objects or proto-PNNi. The very low $`M_\mathrm{e}^\mathrm{R}`$ make the transition time very short, and, as a consequence, the lack of proto-PNNi. Studies of a sample including OH/IR and other transition objects, and PNNi within the same environment, should reflect the mean value of $`M_\mathrm{e}^\mathrm{R}`$ in the sample.
In Figure 24d most post-AGB stars belong to the pre-Planetary Nebula (wind and proto-PNN) phases. Only a few percent of the stars are PNNi. The situation of Figure 24d has been produced only to show an extreme case, since it is highly unrealistic that all stars have the same residual envelope mass, independent of their initial mass on the main sequence, and on their core mass.
### 3.5 The final helium-shell flash
In this paper we did not use the helium-burning (nor a combination of hydrogen and helium burning) stellar models to determine the effect of a helium-burning PNN population. The parametrization that we use is described in $`\mathrm{\S }2.5`$, and it agrees with observations. Following, we show how the presence of helium-burning stars can effect the observable post-AGB populations. In Figure 25 we show the synthetic population on the HR diagram for a sample of 1500 post-AGB stars in which 20 $`\%`$ of the post-AGB stars have experienced a final flash. The fraction of stars in post-FF is chosen in agreement with the guesses based on the evolutionary models by Iben (1984). We compare this population with the basic model. We note, from Table 2, that the overall composition of the sample in terms of ratio of PNNi to other components is (statistically) very similar to the basic sample.
In Figure 25 we have separated the He- and H-burning stars. Panel (a) shows the basic model (1500 objects with Salpeterโs IMF, Weidemannโs IMFMR, $`M_\mathrm{e}^\mathrm{R}`$=f(M), and no post-FF). In the other panels, we show (b) the post-AGB stars having experienced a FF (20$`\%`$ of the sample); (c) stars that are in H-burning, post H-burning and/or will experience FF (the latter group of stars is observationally indistinguishable from H-burning stars); (d) the composite population of panels (b) and (c), thus the observable post-AGB population. Figure 26 shows the same populations of Figure 25 in the $`M_\mathrm{V}t`$ plane. The simulated He-burning populations stand out for lack of low-brightness PNNi.
It is not easy to observationally discern among H- and He-burning stars. The problem is that the stellar abundances are not easy to measure in PNNi. In general, H-depleted post-AGB stars are believed to be He-burners, but it is not so easy to single out the H-burning stars, as H-rich post-AGB could indeed be He-burners as well as H-burners. The simulation of Figures 25 and 26 can help us in this respect. In fact, the ratio of PNNi brighter than, say, 5 magnitudes in the basic sample is about a third of the total sample, while in the composite population (Fig. 26c) this ratio goes up to half the sample. A homogeneous, complete observed PNN sample should show this kind of discrepancy.
### 3.6 The planetary nebula life time
Through this paper we have assumed that the maximum age for a PN depends on the PN mass, as in Equation A2, with constant K=4$`\times 10^5`$. We have also produced other synthetic populations by keeping the nebular mass dependence, and changing the constant. If we keep all other parameters as in the basic model, we obtain that the number of PNNi goes to zero as K is lowered to about 5$`\times 10^3`$. This happens because the shorter lifetime of PNe makes most of the post-AGB population be in a post-PNNi phase. In this case, a large number of lazy AGB stars are also produced. In general, lazy AGB stars are produced by running the basic model with K$`<8\times 10^4`$.
If we were to eliminate the nebular mass dependence from $`t_{\mathrm{max}}`$ the situation changes noticeably. In Figure 27 we show a simulated population with all parameters identical to the basic model, except we have fixed the PN life time to 30,000 years for all nebulae. The result is remarkably different than the basic model (see Fig. 15 and Table 2). Most low-mass post-AGB stars are in post-PNN phase. If we were to lower the maximum PN life time to 10,000 years (not an unreasonable choice) we obtain that most PNNi in the diagram would have a high mass. This is in contrast with the observations, thus our simulations lead us to believe that the lifetime of PNe depends on the nebular mass. Detailed comparison with homogeneous dataset should be used to confirm this inclination.
## 4 Summary and future work
We have used up-to-date evolutionary tracks as templates to build a code for post-AGB population synthesis. Our models aim at understanding the fine tuning of post-AGB evolution, including the consequences of IMF and IMFMR, the transition time and its correlation with the residual envelope mass, the actual duration of the PN phase, the relative population of proto-PNNi, wind objects, and post-PNNi stars, and the occurrence and timing of the FF phase. Our synthetic tracks, available to obtain log $`T_{\mathrm{eff}}`$and log $`L`$for post-AGB stars of any stellar mass in the range $`0.85M_\mathrm{i}<9M_{}`$, are obtained with very high precision in reproducing the actual evolutionary tracks. The interpolation method is simple and fully explained in this paper, so that the readers can reproduce the synthetic populations included here.
In this paper we have shown a sample of the possible applications, without going into detailed comparisons with the observed data. Among the results shown here, we found that (1) the dependence of the synthetic populations on the assumed IMF and IMFMR is mild; the post-AGB populations are not ideal indicators of the IMF. (2) The residual envelope mass, after the envelope ejection, has a strong effect in determining the subsequent post-AGB evolution; its indetermination produces very high indetermination on the transition time, and ultimately on the resulting post-AGB populations. (3) The ratio of He- to H-burning PNNi can be reproduced with population synthesis.
The central importance of this paper consists in showing that the many fundamental variables of stellar evolution have a major role in determining post-AGB populations. The variation and indetermination of these parameters should not be overlooked in comparing data and theory.
The theoretical work contained in this paper was implemented with the goal of being versatile and useful for a full host of applications. Among other possible uses of these synthetic tracks are the studies of evolutionary effects on the Planetary Nebula Luminosity Function (PNLF), and how those translate in the variation of the extragalactic distance scale, as derived from the PNLF. Stanghellini (1995) has shown in a preparatory work that the PNLF is at variation with respect to the transition time. The updated models presented here allow the most detailed study of the PNLF on the evolutionary parameters.
Future applications also include the simulations of bulge, elliptical galaxy, and Magellanic Cloud populations, where a different treatment of the SF/IMF should be used, and possibly different stellar chemistry.
Thanks to Laura Greggio for her help in the interpolation techniques, to Antonella Nota and many others for their bibliographic indications, and to an anonymous referee for pointing out an error in the previous draft and for many suggestions.
## Appendix A Appendix
In order to parametrize the duration of the PN phase, we assume that the nebular mass is 0.02 $`M_{}`$ for M=0.55 $`M_{}`$ (core mass), and 1.5 $`M_{}`$ for M=1.2 $`M_{}`$. The planetary nebula mass can be described by the following correlation between nebular mass and plateau luminosity:
$$M_{\mathrm{PN}}=4.495\times 10^8L_{\mathrm{PL}}^{1.636},$$
$`(A1)`$
and
$$t_{\mathrm{max}}=KM_{\mathrm{PN}}{}_{}{}^{2/5}.$$
$`(A2)`$
To calculate the fading time to $`L=1.0`$$`L_{}`$, $`t_\mathrm{F}`$, we use the VW tracks and we interpolate versus the (post-AGB) mass, obtaining
$$t_\mathrm{F}(M)=7.09572.1441M.$$
$`(A3)`$
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# DISCOVERY OF PROTON DECAY : A MUST FOR THEORY, A CHALLENGE FOR EXPERIMENT11footnote 1Invited talk presented at the Intโl Workshop on Next Generation Nucleon Decay and Neutrino Detector, held at Stony Brook, September, 23-25, 1999,to appear in the Proceedings.
## 1 Introduction
It has been recognized since the early 1970โs that the price one must pay to achieve a unification of quarks and leptons and simultaneously a unity of the three gauge forces, commonly called โgrand unificationโ, is proton decay . This important process, which would provide the window for viewing physics at truly short distances ($`<10^{30}`$ cm), is yet to be seen. Nevertheless, as I will stress in this talk, there have appeared over the years an impressive set of facts, including the meeting of the gauge couplings and neutrino-oscillations, which not only favor the hypothesis of grand unification, but in fact select out a particular route to such unification, based on the ideas of supersymmetry and SU(4)-color . These facts together provide a clear signal that the discovery of proton decay cannot be far behind.
To be specific, working within the framework of a unified theory , that incorporates the ideas mentioned above, I would argue that an improvement in the current sensitivity for detecting proton decay by a modest factor of five to ten should either produce real events, or else the framework would be excluded. In this sense, and as I will elaborate further, the discovery of proton decay is now crucial to the survival of some elegant ideas on unification, which are otherwise so successful. By the same token, proving or disproving their prediction on proton decay poses a fresh challenge to experiment.
The pioneering efforts by several physicists in the mid 1950โs through the early 70โs had provided a lower limit on the proton lifetime of about $`10^{26}`$ yrs, independent of decay modes, and $`10^{29}`$-$`10^{30}`$ yrs for the $`e^+\pi ^0`$-mode. Subsequent searches at the Kolar Goldfield and the NUSEX detectors in the early 80โs pushed this limit to about $`10^{31}`$ years in the $`e^+\pi ^0`$-mode. Following the suggestion of proton decay in the context of grand unification, and thanks to the initiative of several experimenters, two relatively large-size detectors - IMB and Kamiokande - were built in the 80โs, where dedicated searches for proton decay were carried out with higher sensitivity. These detectors helped to push the lower limit in the $`e^+\pi ^0`$ channel to about $`10^{32}`$ yrs. This in turn clearly disfavored the minimal non-supersymmetric SU(5)-model of grand unification \- a conclusion that was strengthened subsequently by the measurements of the gauge couplings at LEP as well (see discussion later).
The searches for proton decay now continues with still greater sensitivity at the largest detector so far - at SuperKamiokande, completed in 1996. It is worth noting at this point that, although these detectors have not revealed proton decay yet, they did bring some major bonuses of monumental importance. These include : (a) the detection of the neutrinos from the supernova 1987a, (b) confirmation of the solar neutrino-deficit, and last but not least, (c) the discovery of atmospheric neutrino-oscillation. The SuperK water-Cerenkov detector with a fiducial volume of 22.5 kilotons currently provides (with three years of running) a lower limit on the inverse rate of proton decay of about $`1.6\times 10^{33}`$ yrs for the theoretically favored ($`\overline{\nu }K^+`$)-channel and of about $`3.8\times 10^{33}`$ yrs for the ($`e^+\pi ^0`$)-channel . It has the capability of improving these limits by a factor of two to three in each case within the next decade, unless of course it discovers real events for proton decay or strong candidate events in the meantime. I will return to this point and its relevance to theoretical expectations for proton decay in just a bit.
While proton decay is yet to be observed, it is worth stressing at this point, that the hypothesis of grand unification, especially that based on the ideas of SU(4)-color, left-right gauge symmetry, and supersymmetry, is now supported by several observations. As I will explain in sections 2-5, these include :
(a) The observed family structure : The five scattered multiplets of the standard model, belonging to a family, neatly become parts of a whole (a single multiplet), with their weak hypercharges precisely predicted by grand unification. Realization of this feature calls for an extension of the standard model symmetry G(213)$`=`$SU(2$`)_L\times `$U(1$`)_Y\times `$SU(3$`)^C`$ minimally to the symmetry group G(224)$`=`$SU(2$`)_L\times `$SU(2$`)_R\times `$SU(4$`)^C`$ , which can be extended further into the simple group SO(10) , but not SU(5) . The G(224) symmetry in turn introduces some additional attractive features (see Sec. 2), including especially the right-handed (RH) neutrinos ($`\nu _R`$โs) accompanying the left-handed ones ($`\nu _L`$โs), and $`B`$-$`L`$ as a local symmetry. As we will see, both of these features, which are special to G(224), now seem to be needed on empirical grounds.
(b) Meeting of the gauge couplings : Such a meeting is found to occur at a scale $`M_X2\times 10^{16}`$ GeV, when the three gauge couplings are extrapolated from their values measured at LEP to higher energies, in the context of supersymmetry . This dramatic phenomenon supports the ideas of both grand unification and supersymmetry. These in turn may well emerge from a string theory or M-theory (see discussion in Sec. 3).
(c) Mass of $`\nu _\tau \mathrm{๐}/\mathrm{๐๐}`$ eV : Subject to the well-motivated assumption of hierarchical neutrino masses, the recent discovery of atmospheric neutrino-oscillation at SuperKamiokande suggests a value for $`m(\nu _\tau )1/20eV`$. It has been argued (see e.g. Ref. ) that a mass for $`\nu _\tau `$ of this magnitude points to the need for RH neutrinos, and that it goes extremely well with the hypothesis of a supersymmetric unification, based on either a string-unified G(224) symmetry or SO(10). The SUSY unification-scale as well as SU(4)-color play crucial roles in making this argument.
(d) Some intriguing features of fermion masses and mixings : These include : (i) the โobservedโ near equality of the masses of the b-quark and the $`\tau `$-lepton at the unification-scale ($`m_b^0m_\tau ^0`$); (ii) the empirical Georgi-Jarsklog relations: $`m_s^0m_\mu ^0/3`$ and $`m_d^03m_e^0`$, and (iii) the observed largeness of the $`\nu _\mu `$-$`\nu _\tau `$ oscillation angle ($`\mathrm{sin}^22\theta _{\nu _\mu \nu _\tau }^{\text{osc}}0.83`$) , together with the smallness of the corresponding quark mixing parameter $`V_{bc}(0.04)`$ . As shown in recent work by Babu, Wilczek and me , it turns out that these features and more can be understood remarkably well (see discussion in Sec 5) within an economical and predictive SO(10)-framework based on a minimal Higgs system. The success of this framework is in large part due simply to the group-structure of SO(10). For most purposes, that of G(224) suffices.
(e) Baryogenesis : To implement baryogenesis successfully, in the presence of electroweak sphaleron effects , which wipe out out any baryon excess generated at high temperatures in the ($`B`$-$`L`$)-conserving mode, it has become apparent that one would need $`B`$-$`L`$ as a generator of the underlying symmetry, whose spontaneous violation at high temperatures would yield, for example, lepton asymmetry (leptogenesis). The latter in turn is converted to baryon-excess at lower temperatures by electroweak sphalerons. This mechanism, it turns out, yields even quantitatively the right magnitude for baryon excess . The need for $`B`$-$`L`$, which is a generator of SU(4)-color, again points to the need for G(224) or SO(10) as an effective symmetry near the unification-scale $`M_X`$.
The success of each of these five features (a)-(e) seems to be non-trivial. Together they make a strong case for both supersymmetric grand unification and simultaneously for the G(224)/SO(10)-route to such unification, as being relevant to nature. However, despite these successes, as long as proton decay remains undiscovered, the hallmark of grand unification - that is quark-lepton transformability \- would remain unrevealed.
The relevant questions in this regard then are : What is the predicted range for the lifetime of the proton - in particular an upper limit - within the emperically favored route to unification mentioned above? What are the expected dominant decay modes within this route? Are these predictions compatible with current lower limits on proton lifetime mentioned above, and if so, can they still be tested at the existing or possible near-future detectors for proton decay?
Fortunately, we are in a much better position to answer these questions now, compared to a few years ago, because meanwhile we have learnt more about the nature of grand unification. As noted above (see also Secs. 2 and 4), the neutrino masses and the meeting of the gauge coupleings together seem to select the supersymmetric G(224)/SO(10)-route to higher unification. The main purpose of my talk here will therefore be to address the questions raised above, in the context of this route. For the sake of comparison, however, I will state the corresponding results for the case of supersymmetric SU(5) as well.
My discussion will be based on a recent study of proton decay by Babu, Wilczek and me ,which, relative to previous ones, has three distinctive features :
(a) It systematically takes into account the link that exists between proton decay and the masses and mixings of all fermions, including the neutrinos.
(b) In particular, in adition to the contributions from the so-called โstandardโ $`d=5`$ operators (see Sec. 6), it includes those from a new set of $`d=5`$ operators, related to the Majorana masses of the RH neutrinos . These latter are found to be as important as the standard ones.
(c) The work also incorporates GUT-scale threshold effects, which arise because of mass-splittings between the components of the SO(10)-multiplets, and lead to differences between the three gauge couplings.
Each of these features turn out to be crucial to gaining a reliable insight into the nature of proton decay. Our study shows that the inverse decay rate for the $`\overline{\nu }K^+`$-mode, which is dominant, is less than about $`7\times 10^{33}`$ yrs. This upper bound is obtained by making generous allowance for uncertainties in the matrix elements and the SUSY-spectrum. Typically, the lifetime should of course be less than this bound. Furthermore, due to contributions from the new operators, the $`\mu ^+K^0`$-mode is found to be prominent, with a branching ratio typically in the range of 10-50%. By contrast, minimal SUSY SU(5), for which the new operators are absent, would lead to branching ratios $`10^3`$ for this mode. Thus our study of proton decay , correlated with fermion masses, strongly suggests that discovery of proton decay should be around the corner. In fact,one expects that at least candidate events should be observed in the near future already at SuperK. However, allowing for the possibility that the proton lifetime may well be closer to the upper bound stated above, a next-generation detector providing a net gain in sensitivity in proton decay-searches by a factor of 5-10, compared to SuperK, would certainly be needed not just to produce proton-decay events, but also to clearly distinguish them from the background. It would of course also be essential to study the branching ratios of certain sub-dominant but crucial decay modes, such as the $`\mu ^+K^0`$. The importance of such improved sensitivity, in the light of the successes of supersymmetric grand unification, is emphasized at the end.
## 2 Advantages of the Symmetry G(224) as a Step to Higher Unification
The standard model (SM) based on the gauge symmetry G(213)$`=`$SU(2$`)_L\times `$U(1$`)_Y\times `$SU(3$`)^C`$ has turned out to be extremely successful empirically. It has however been recognized since the early 1970โs that, judged on aesthetic merits, it has some major shortcomings. For example, it puts members of a family into five scattered multiplets, without providing a compelling reason for doing so. It also does not provide a fundamental reason for the quantization of electric charge. Nor does it explain the co-existence of quarks and leptons, and that of the three gauge forces, with their differing strengths. The idea of grand unification was postulated precisely to remove these shortcomings. That in turn calls for the existence of fundamentally new physics, far beyond that of the standard model. As mentioned before, recent experimental findings, including the meetings of the gauge couplings and neutrino-oscillations, seem to go extremely well with this line of thinking.
To illustrate the advantage of an early suggestion in this regard, consider the five standard model multiplets belonging to the electron-family as shown :
$`\left(\begin{array}{ccc}u_ru_yu_b& & \\ d_rd_yd_b& & \end{array}\right)_L^{\frac{1}{3}};\left(\begin{array}{ccc}u_ru_yu_b& & \end{array}\right)_R^{\frac{4}{3}};\left(\begin{array}{ccc}d_rd_yd_b& & \end{array}\right)_R^{\frac{2}{3}};\left(\begin{array}{c}\nu _e\\ e^{}\end{array}\right)_L^{\mathrm{\hspace{0.17em}1}};\left(e^{}\right)_R^{\mathrm{\hspace{0.17em}2}}.`$ (7)
Here the superscripts denote the respective weak hypercharges $`Y_W`$ (where $`Q_{em}=I_{3L}+Y_W/2`$) and the subscripts L and R denote the chiralities of the respective fields. If one asks : how one can put these five multiplets into just one multiplet, the answer turns out to be simple and unique. As mentioned in the introduction, the minimal extension of the SM symmetry G(213) needed, to achieve this goal, is given by the gauge symmetry :
$`\text{G(224)}=\text{SU(2})_L\times \text{SU(2})_R\times \text{SU(4})^C.`$ (8)
Subject to left-right discrete symmetry ($`LR`$), which is natural to G(224), all members of the electron family fall into the neat pattern :
$`F_{L,R}^e=\left[\begin{array}{cccc}u_ru_yu_b\nu _e& & & \\ d_rd_yd_be^{}& & & \end{array}\right]_{L,R}`$ (11)
The multiplets $`F_L^e`$ and $`F_R^e`$ are left-right conjugates of each other and transform respectively as (2,1,4) and (1,2,4) of G(224); likewise for the muon and the tau families. Note that the symmetries SU(2$`)_L`$ and SU(2$`)_R`$ are just like the familiar isospin symmetry, except that they operate on quarks and well as leptons, and distinguish between left and right chiralities. The left weak-isospin SU(2$`)_L`$ treats each column of $`F_L^e`$ as a doublet; likewise SU(2$`)_R`$ for $`F_R^e`$; the symmetry SU(4)-color treats each row of $`F_L^e`$ and $`F_R^e`$ as a quartet, interpreting lepton number as the fourth color. Note also that postulating either SU(4)-color or SU(2$`)_R`$ forces one to introduce a right-handed neutrino ($`\nu _R`$) for each family as a singlet of the SM symmetry. This requires that there be sixteen two-component fermions in each family, as opposed to fifteen for the SM. The symmetry G(224) introduces an elegant charge formula :
$`Q_{em}=I_{3L}+I_{3R}+{\displaystyle \frac{BL}{2}}`$ (12)
expressed in terms of familiar quantum numbers $`I_{3L}`$, $`I_{3R}`$ and $`B`$-$`L`$, which applies to all forms of matter (including quarks and leptons of all six flavors, gauge and Higgs bosons). Note that the weak hypercharge given by $`Y_W/2=I_{3R}+\frac{BL}{2}`$ is now completely determined for all members of the family. The values of $`Y_W`$ thus obtained precisely match the assignments shown in Eq. (7). Quite clearly, the charges $`I_{3L}`$, $`I_{3R}`$ and $`B`$-$`L`$, being generators respectively of SU(2$`)_L`$, SU(2$`)_R`$ and SU(4$`)^c`$, are quantized; so also then is the electric charge $`Q_{em}`$.
In brief, the symmetry G(224) brings some attaractive features to particle physics. These include :
(i) Organization of all 16 members of a family into one left-right self-conjugate multiplet;
(ii) Quantization of electric charge;
(iii) Quark-lepton unification (through SU(4) color);
(iv) Conservation of parity at a fundamental level ;
(v) Right-handed neutrinos ($`\nu _R^{}s`$) as a compelling feature; and
(vi) $`B`$-$`L`$ as a local symmetry.
As mentioned in the introduction, the two distinguishing features of G(224) - i.e. the existence of the RH neutrinos and $`B`$-$`L`$ as a local symmetry - now seem to be needed on empirical grounds.
Believing in a complete unification, one is led to view the G(224) symmetry as part of a bigger symmetry, which itself may have its origin in an underlying theory, such as string theory. In this context, one might ask : Could the effective symmetry below the string scale in four dimensions (see sec.3) be as small as just the SM symmetry G(213), even though the latter may have its origin in a bigger symmetry, which lives however only in higher dimensions? I will argue in Sec. 4 that the data on neutrino masses and the need for baryogenesis provide an answer to the contrary, suggesting clearly that it is the effective symmetry in four dimensions, below the string scale, which must minimally contain either G(224) or a close relative G(214)$`=`$SU(2$`)_L\times `$I$`{}_{3R}{}^{}\times `$SU(4$`)^C`$.
One may also ask : does the effective four dimensional symmetry have to be any bigger than G(224) near the string scale? In preparation for an answer to this question, let us recall that the smallest simple group that contains the SM symmetry G(213) is SU(5) . It has the virtue of demonstrating how the main ideas of grand unification, including unification of the gauge couplings, can be realized. However, SU(5) does not contain G(224) as a subgroup. As such, it does not possess some of the advantages listed above. In particular, it does not contain the RH neutrinos as a compelling feature, and $`B`$-$`L`$ as a local symmetry. Furthermore, it splits members of a family into two multiplets : $`\overline{5}+10`$.
By contrast, the symmetry SO(10) has the merit, relative to SU(5), that it contains G(224) as a subgroup, and thereby retains all the advantages of G(224) listed above. (As a historical note, it is worth mentioning that these advantages had been motivated on aesthetic grounds through the symmetry G(224) , and all the ideas of higher unification were in place , before it was noted that G(224)(isomorphic to SO(4)$`\times `$SO(6)) embeds nicely into SO(10) ). Now, SO(10) even preserves the 16-plet family-structure of G(224) without a need for any extension. By contrast, if one extends G(224) to the still higher symmetry E<sub>6</sub> , the advantages (i)-(vi) are retained, but in this case, one must extend the family-structure from a 16 to a 27-plet, by postulating additional fermions. In this sense, there seems to be some advantage in having the effective symmetry below the string scale to be minimally G(224) (or G(214)) and maximally no more than SO(10). I will compare the relative advantage of having either a string-derived G(224) or a string-SO(10), in the next section. First, I discuss the implications of the data on coupling unification.
## 3 The Need for Supersymmetry : MSSM versus String Unifications
It has been known for some time that the precision measurements of the standard model coupling constants (in particular $`\mathrm{sin}^2\theta _W`$) at LEP put severe constraints on the idea of grand unification. Owing to these constraints, the non-supersymmetric minimal SU(5), and for similar reasons, the one-step breaking minimal non-supersymmetric SO(10)-model as well, are now excluded . But the situation changes radically if one assumes that the standard model is replaced by the minimal supersymmetric standard model (MSSM), above a threshold of about 1 TeV. In this case, the three gauge couplings are found to meet , at least approximately, provided $`\alpha _3(m_Z)`$ is not too low (see Figs. in e.g. Refs. ). Their scale of meeting is given by
$`M_X2\times 10^{16}\text{GeV\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}(MSSM or SUSY\hspace{0.17em}\hspace{0.17em}SU(5))}`$ (13)
This dramatic meeting of the three gauge couplings, or equivalently the agreement of the MSSM-based prediction of $`\mathrm{sin}^2\theta _W(m_Z)_{\text{Th}}=0.2315\pm 0.003`$ with the observed value of $`\mathrm{sin}^2\theta _W(m_Z)=0.23124\pm 0.00017`$ , provides a strong support for the ideas of both grand unification and supersymmetry, as being relevant to physics at short distances.
The most straightforward interpretation of the observed meeting of the three couplings and of the scale $`M_X`$, is that a supersymmetric grand unification symmetry (often called GUT symmetry), like SU(5) or SO(10), breaks spontaneously at $`M_X`$ into the standard model symmetry G(213).
In the context of string or M theory, which seems to be needed to unify all the forces of nature including gravity and also to obtain a good quantum theory of gravity, an alternative interpretation is however possible. This is because, even if the effective symmetry in four dimensions emerging from a higher dimensional string theory is non-simple, like G(224) or G(213), string theory can still ensure familiar unification of the gauge couplings at the string scale. In this case, however, one needs to account for the small mismatch between the MSSM unification scale $`M_X`$ (given above), and the string unification scale, given by $`M_{st}g_{st}\times 5.2\times 10^{17}`$ GeV $`3.6\times 10^{17}`$ GeV (Here we have put $`\alpha _{st}=\alpha _{GUT}(\text{MSSM})0.04`$) . Possible resolutions of this mismatch have been proposed. These include : (i) utilizing the idea of string-duality which allows a lowering of $`M_{st}`$ compared to the value shown above, or alternatively (ii) the idea of a semi-perturbative unification that assumes the existence of two vector-like families, transforming as $`(16+\overline{16})`$, at the TeV-scale. The latter raises $`\alpha _{GUT}`$ to about 0.25-0.3 and simultaneously $`M_X`$, in two loop, to about $`(1/22)\times 10^{17}`$ GeV (Other mechanisms resolving the mismatch are reviewed in Refs. and ). In practice, a combination of the two mechanisms mentioned above may well be relevant<sup>2</sup><sup>2</sup>2I have in mind the possibility of string-duality lowering $`M_{st}`$ for the case of semi-perturbative unification (for which $`\alpha _{st}`$0.25, and thus, without the use of string-duality, $`M_{st}`$ would be about $`10^{18}`$ GeV) to a value of about (1-2)$`\times 10^{17}`$ GeV (say), and semi-perturbative unification raising the MSSM value of $`M_X`$ to about 5$`\times 10^{16}`$ GeV$``$ $`M_{st}`$(1/2 to 1/4) (say). In this case, an intermediate symmetry like G(224) emerging at $`M_{st}`$ would be effective only within the short gap between $`M_{st}`$ and $`M_X`$, where it would break into G(213). Despite this short gap, one would still have the benefits of SU(4)-color that are needed to understand neutrino masses (see sec.4). At the same time, Since the gap is so small, the couplings of G(224), unified at $`M_{st}`$ would remain essentially so at $`M_X`$, so as to match with the โobservedโ coupling unification, of the type suggested in Ref. ..
While the mismatch can thus quite plausibly be removed for a non-GUT string-derived symmetry like G(224) or G(213), a GUT symmetry like SU(5) or SO(10) would have an advantage in this regard because it would keep the gauge couplings together between $`M_{st}`$ and $`M_X`$ (even if $`M_XM_{st}/20`$), and thus not even encounter the problem of a mismatch between the two scales. A supersymmetric GUT-solution (like SU(5) or SO(10)), however, has a possible disadvantage as well, because it needs certain color triplets to become superheavy by the so-called double-triplet splitting mechanism (see Sec. 6 and Appendix), in order to avoid the problem of rapid proton decay. However, no such mechanism has emerged yet, in string theory, for the GUT-like solutions .
Non-GUT string solutions, based on symmetries like G(224) or G(2113) for example, have a distinct advantage in this regard, in that the dangerous color triplets, which would induce rapid proton decay, are often naturally projected out for such solutions . Furthermore, the non-GUT solutions invariably possess new โflavorโ gauge symmetries, which distinguish between families. These symmetries are immensely helpful in explaining qualitatively the observed fermion mass-hierarchy (see e.g. Ref. ) and resolving the so-called naturalness problems of supersymmetry such as those pertaining to the issues of squark-degeneracy , CP violation and quantum gravity-induced rapid proton decay .
Weighing the advantages and possible disadvantages of both, it seems hard at present to make a priori a clear choice between a GUT versus a non-GUT string-solution. As expressed elsewhere , it therefore seems prudent to keep both options open and pursue their phenomenological consequences. Given the advantages of G(224) or SO(10) in the light of the neutrino masses (see Secs. 2 and 4), I will thus proceed by assuming that either a suitable G(224)-solution with a mechanism of the sort mentioned above, or a realistic SO(10)-solution with the needed doublet-triplet mechanism, will emerge from string theory. We will see that with this broad assumption an economical and predictive framework emerges, which successfully accounts for a host of observed phenomena, and makes some crucial testable predictions. Fortunately, it will turn out that there are many similarities between the predictions of a string-unified G(224) and SO(10), not only for the neutrino and the charged fermion masses, but also for proton decay. I next discuss the implications of the mass of $`\nu _\tau `$ suggested by the SuperK data.
## 4 Mass of $`\nu _\tau `$: Evidence In Favor of the G(224) Route
One can obtain an estimate for the mass of $`\nu _L^\tau `$ in the context of G(224) or SO(10) by using the following three steps (see e.g.Ref.):
(i) Assume that B$``$L and $`I_{3R}`$, contained in a string-derived G(224) or SO(10), break near the unification-scale:
$`M_X2\times 10^{16}\text{GeV},`$ (14)
through VEVs of Higgs multiplets of the type suggested by string-solutions - i.e. $`(1,2,4)_H`$ for G(224) or $`\overline{16}_H`$ for SO(10), as opposed to $`126_H`$ . In the process, the RH neutrinos ($`\nu _R^i`$), which are singlets of the standard model, can and generically will acquire superheavy Majorana masses of the type $`M_R^{ij}\nu _R^{iT}C^1\nu _R^j`$, by utilizing the VEV of $`\overline{16}_H`$ and effective couplings of the form:
$`_M(SO(10))=f_{ij}\mathrm{\hspace{0.17em}\hspace{0.17em}16}_i16_j\overline{16}_H\overline{16}_H/M+h.c.`$ (15)
A similar expression holds for G(224). Here $`i,j=1,2,3`$, correspond respectively to $`e,\mu `$ and $`\tau `$ families. Such gauge-invariant non-renormalizable couplings might be expected to be induced by Planck-scale physics, involving quantum gravity or stringy effects and/or tree-level exchange of superheavy states, such as those in the string tower. With $`f_{ij}`$ (at least the largest among them) being of order unity, we would thus expect M to lie between $`M_{Planck}2\times 10^{18}`$ GeV and $`M_{string}4\times 10^{17}`$ GeV. Ignoring for the present off-diagonal mixings (for simplicity), one thus obtains <sup>3</sup><sup>3</sup>3The effects of neutrino-mixing and of possible choice of $`M=M_{string}4\times 10^{17}`$ GeV (instead of $`M=M_{Planck}`$) on $`M_{3R}`$ are considered in Ref. .:
$`M_{3R}{\displaystyle \frac{f_{33}\overline{16}_H^2}{M}}f_{33}(2\times 10^{14}\text{GeV})\eta ^2(M_{Planck}/M)`$ (16)
This is the Majorana mass of the RH tau neturino. Guided by the value of $`M_X`$, we have substituted $`\overline{16}_H=(2\times 10^{16}\text{GeV})\eta `$ ,with $`\eta 1/2`$ to 2(say).
(ii) Now using SU(4)-color and the Higgs multiplet $`(2,2,1)_H`$ of G(224) or equivalently $`10_H`$ of SO(10), one obtains the relation $`m_\tau (M_X)=m_b(M_X)`$, which is known to be successful. Thus, there is a good reason to believe that the third family gets its masses primarily from the $`10_H`$ or equivalently $`(2,2,1)_H`$ (see sec.5). In turn, this implies:
$`m(\nu _{Dirac}^\tau )m_{top}(M_X)(100\text{-}\mathrm{\hspace{0.17em}120})\text{GeV}`$ (17)
Note that this relationship between the Dirac mass of the tau-neutrino and the top-mass is special to SU(4)-color. It does not emerge in SU(5).
(iii) Given the superheavy Majorana masses of the RH neutrinos as well as the Dirac masses as above, the see-saw mechanism yields naturally light masses for the LH neutrinos. For $`\nu _L^\tau `$ (ignoring mixing), one thus obtains, using Eqs. (16) and (17),
$`m(\nu _L^\tau ){\displaystyle \frac{m(\nu _{Dirac}^\tau )^2}{M_{3R}}}[(1/20)\text{eV}(1\text{-}\mathrm{\hspace{0.17em}1.44})/f_{33}\eta ^2](M/M_{Planck})`$ (18)
Now,assuming the hierarchical pattern $`m(\nu _L^e)m(\nu _L^\mu )m(\nu _L^\tau )`$, which is suggested by the see-saw mechanism,and further that the SuperK observation represents $`\nu _L^\mu \nu _L^\tau `$ (rather than $`\nu _L^\mu \nu _X`$) oscillation, the observed $`\delta m^21/2(10^2\text{-}\mathrm{\hspace{0.17em}10}^3)`$eV<sup>2</sup> corresponds to $`m(\nu _L^\tau )`$ (1/15 - 1/40) eV. It seems truly remarkable that the expected magnitude of $`m(\nu _L^\tau )`$, given by Eq.(18), is just about what is suggested by the SuperK data, if $`f_{33}\eta ^2(M_{Planck}/M)`$ 1.3 to 1/2. Such a range for $`f_{33}\eta ^2(M_{Planck}/M)`$ seems most plausible and natural (see discussion in Ref. ). Note that the estimate (18) crucially depends upon the supersymmetric unification scale, which provides a value for $`M_{3R}`$, as well as on SU(4)-color that yields $`m(\nu _{Dirac}^\tau )`$.The agreement between the expected and the SuperK result thus clearly suggests that the effective symmetry below the string-scale should contain SU(4)-color. Thus, minimally it should be either G(214) or G(224), and maximally as big as SO(10), if not E<sub>6</sub>.
By contrast, if SU(5) is regarded as either a fundamental symmetry or as the effective symmetry below the string scale, there would be no compelling reason based on symmetry alone, to introduce a $`\nu _R`$, because it is a singlet of SU(5). Second, even if one did introduce $`\nu _R^i`$ by hand, their Dirac masses, arising from the coupling $`h^i\overline{5}_i5_H\nu _R^i`$, would be unrelated to the up-flavor masses and thus rather arbitrary (contrast with Eq. (17)). So also would be the Majorana masses of the $`\nu _R^i`$โs, which are SU(5)-invariant, and thus can be even of order string scale . This would give $`m(\nu _L^\tau )`$ in gross conflict with the observed value.
Before passing to the next section, it is worthnoting that the mass of $`\nu _\tau `$ suggested by SuperK, as well as the observed value of $`\mathrm{sin}^2\theta _W`$ (see Sec.3), provide valuable insight into the nature of GUT symmetry breaking. They both favor the case of a single-step breaking (SSB) of SO(10) or a string-unified G(224) symmetry at a scale of order $`M_X`$, into the standard model symmetry G(213), as opposed to that of a multi-step breaking (MSB). The latter would correspond, for example, to SO(10) (or G(224)) breaking at a scale $`M_1`$ into G(2213), which in turn breaks at a scale $`M_2<<M_1`$ into G(213). One reason why the case of single-step breaking is favored over that of multi-step breaking is that the latter can accomodate but not really predict $`\mathrm{sin}^2\theta _W`$, where as the former predicts the same successfully. Furthermore, since the Majorana mass of $`\nu _R^\tau `$ arises arises only after $`BL`$ and $`I_{3R}`$ break, it would be given, for the case of MSB, by $`M_{3R}f_{33}(M_2^2/M)`$, where $`MM_{st}`$ (say). If $`M_2M_X2\times 10^{16}`$ GeV, and M $`>M_X`$, one would obtain too low a value ($`<<10^{14}`$ GeV) for $`M_{3R}`$ (compare with Eq.(8)),and thereby too large a value for $`m(\nu _L^\tau )`$, compared to that suggested by SuperK. By contrast, the case of SSB yields the right magnitude for $`m(\nu _\tau )`$ (see Eq. (10)).
Thus the success of the result on $`m(\nu _\tau )`$ discussed above not only favors the symmetry G(224) or SO(10), but also clearly suggests that $`BL`$ and $`I_{3R}`$ break near the conventional GUT scale $`M_X2\times 10^{16}`$ GeV, rather than at an intermediate scale $`<<M_X`$. In other words, the observed values of both $`\mathrm{sin}^2\theta _W`$ and $`m(\nu _\tau )`$ favor only the simplest pattern of symmetry-breaking, for which SO(10) or a string-derived G(224) symmetry breaks in one step to the standard model symmetry, rather than in multiple steps. It is of course only this simple pattern of symmetry breaking that would be rather restrictive as regards its predictions for proton decay (to be dicussed in Sec.6). I next dicuss the problem of understanding the masses and mixings of all fermions.
## 5 Understanding Fermion Masses and Neutrino Oscillations in SO(10)
Understanding the masses and mixings of all quarks and charged leptons, in conjunction with those of the neutrinos, is a goal worth achieving by itself. It also turns out to be essential for the study of proton decay. I therefore present first a recent attempt in this direction, which seems most promising . A few guidelines would prove to be helpful in this regard. The first of these is motivated by the desire for economy and the rest by data.
1) Hierarchy Through Off-diagonal Mixings : Recall earlier attempts that attribute hierarchical masses of the first two families to matrices of the form :
$`M=\left(\begin{array}{cc}0ฯต& \\ ฯต\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}1}& \end{array}\right)m_s^{(0)},`$ (21)
for the $`(d,s)`$ quarks, and likewise for the $`(u,c)`$ quarks. Here $`ฯต1/10`$. The hierarchical patterns in Eq. (21) can be ensured by imposing a suitable flavor symmetry which distinguishes between the two families (that in turn may have its origin in string theory (see e.g. Ref ). Such a pattern has the virtues that (a) it yields a hierarchy that is much larger than the input parameter $`ฯต`$ : $`(m_d/m_s)ฯต^2ฯต`$, and (b) it leads to an expression for the cabibbo angle :
$`\theta _c\left|\sqrt{{\displaystyle \frac{m_d}{m_s}}}e^{i\varphi }\sqrt{{\displaystyle \frac{m_u}{m_c}}}\right|,`$ (22)
which is rather successful. Using $`\sqrt{m_d/m_s}0.22`$ and $`\sqrt{m_u/m_c}0.06`$, we see that Eq. (22) works to within about $`25\%`$ for any value of the phase $`\varphi `$. Note that the square root formula (like $`\sqrt{m_d/m_s}`$) for the relevant mixing angle arises because of the symmetric form of $`M`$ in Eq. (21), which in turn is ensured if the contributing Higgs is a 10 of SO(10). A generalization of the pattern in Eq. (21) would suggest that the first two families (i.e. the $`e`$ and the $`\mu `$) receive masses primarily through their mixing with the third family $`(\tau )`$, with $`(1,3)`$ and $`(1,2)`$ elements being smaller than the $`(2,3)`$; while $`(2,3)`$ is smaller than the $`(3,3)`$. We will follow this guideline, except for the modification noted below.
2) The Need for an Antisymmetric Component : Although the symmetric hierarchical matrix in Eq. (21) works well for the first two families, a matrix of the same form fails altogether to reproduce $`V_{cb}`$, for which it yields :
$`V_{cb}\left|\sqrt{{\displaystyle \frac{m_s}{m_b}}}e^{i\chi }\sqrt{{\displaystyle \frac{m_c}{m_t}}}\right|.`$ (23)
Given that $`\sqrt{m_s/m_b}0.17`$ and $`\sqrt{m_c/m_t}\mathrm{0.0.06}`$, we see that Eq. (23) would yield $`V_{cb}`$ varying between 0.11 and 0.23, depending upon the phase $`\chi `$. This is too big, compared to the observed value of $`V_{cb}0.04\pm 0.003`$, by at least a factor of 3. We interpret this failure as a clue to the presence of an antisymmetric component in $`M`$, together with symmetrical ones (thus $`m_{ij}m_{ji}`$), which would modify the relevant mixing angle to $`\sqrt{{\displaystyle \frac{m_i}{m_j}}}\sqrt{{\displaystyle \frac{m_{ij}}{m_{ji}}}}`$, where $`m_i`$ and $`m_j`$ denote the respective eigenvalues.
3) The Need for a Contribution Proportional to $`B`$-$`L`$ : The success of the relations $`m_b^0m_\tau ^0`$, and $`m_t^0m(\nu _\tau )_{Dirac}^0`$ (see Sec. 4), suggests that the members of the third family get their masses primarily from the VEV of a SU(4)-color singlet Higgs field that is independent of $`B`$-$`L`$. This is in fact ensured if the Higgs is a 10 of SO(10). However, the empirical observations of $`m_s^0m_\mu ^0/3`$ and $`m_d^03m_e^0`$ clearly call for a contribution proportional to $`B`$-$`L`$ as well. Further, one can in fact argue that the suppression of $`V_{bc}`$ (in the quark-sector) together with an enhancement of $`\theta _{\nu _\mu \nu _\tau }^{osc}`$ (in the lepton sector) calls for a contribution that is not only proportional to $`B`$-$`L`$ but is also antisymmetric in the family space (as suggested above in item (8)). We note below how both of these requirements can be met, rather easily, in SO(10), even for a minimal Higgs system.
4) Up-Down Asymmetry : Finally, the up and the down-sector mass matrices must not be proportional to each other, as otherwise the CKM angles would all vanish.
Following Ref. , I now present a simple and predictive mass-matrix, based on SO(10), that satisfies all three requirements, (8), (11) and (12). The interesting point is that one can obtain such a mass-matrix for the fermions by utilizing only the minimal Higgs system, that is needed anyway to break the gauge symmetry SO(10). It consists of the set :
$`H_{minimal}=\{45_H,\mathrm{\hspace{0.17em}16}_H,\overline{16}_H,\mathrm{\hspace{0.17em}10}_H\}.`$ (24)
Of these, the VEV of $`45_HM_X`$ breaks SO(10) into G(2213), and those of $`16_H=\overline{16}_HM_X`$ break G(2213) to G(213), at the unification-scale $`M_X`$. Now G(213) breaks at the electroweak scale by the VEV of $`10_H`$ to U$`(1)_{em}\times `$ SU$`(3)^c`$.
One might have introduced large-dimensional tensorial multiplets of SO(10) like 12$`6_H`$ and 12$`0_H`$, both of which possess cubic level Yukawa couplings with the fermions. In particular, the coupling $`16_i16_j(120_H)`$ would give the desired family-antisymmetric as well as ($`B`$-$`L`$)-dependent contribution. We do not however introduce these multiplets in part because they do not seem to arise in string solutions , and in part also because mass-splittings within such large-dimensional multiplets tend to give excessive threshold corrections to $`\alpha _3(m_z)`$ (typically exceeding 20%), rendering observed coupling unification fortuitous. By contrast, the multiplets in the minimal set (shown above) do arise in string solutions leading to SO(10). Furthermore, the threshold corrections for the minimal set are found to be naturally small, and even to have the right sign, to go with the observed coupling unification .
The question is : does the minimal set meet all the requirements listed above? Now $`10_H`$ (even several 10โs) can not meet the requirements of antisymmetry and $`(B`$-$`L)`$-dependence. Furthermore, a single $`10_H`$ cannot generate CKM-mixings. This impasse disappears,however, as soon as one allows for not only cubic, but also effective non-renormalizable quartic couplings of the minimal set of Higgs fields with the fermions. These latter couplings could of course well arise through exchanges of superheavy states (e.g. those in the string tower) involving renormalizable couplings, and/or through quantum gravity.
Allowing for such cubic and quartic couplings and adopting the guideline (7) of hierarchical Yukawa couplings, as well as that of economy, we are led to suggest the following effective lagrangian for generating Dirac masses and mixings of the three families (for a related but different pattern, involving a non-minimal Higgs system, see Ref ).
$`_{\mathrm{๐๐ฎ๐ค}}=h_{33}\mathbf{\hspace{0.17em}16}_\mathrm{๐}\mathbf{\hspace{0.17em}16}_\mathrm{๐}\mathbf{\hspace{0.17em}10}_๐+[h_{23}\mathbf{\hspace{0.17em}16}_\mathrm{๐}\mathbf{\hspace{0.17em}16}_\mathrm{๐}\mathbf{\hspace{0.17em}10}_๐+a_{23}\mathbf{\hspace{0.17em}16}_\mathrm{๐}\mathbf{\hspace{0.17em}16}_\mathrm{๐}\mathbf{\hspace{0.17em}10}_๐\mathbf{\hspace{0.17em}45}_๐/M`$ (25)
$`+g_{23}\mathbf{\hspace{0.17em}16}_\mathrm{๐}\mathbf{\hspace{0.17em}16}_\mathrm{๐}\mathbf{\hspace{0.17em}16}_๐\mathbf{\hspace{0.17em}16}_๐/M]+\{a_{12}\mathbf{\hspace{0.17em}16}_\mathrm{๐}\mathbf{\hspace{0.17em}16}_\mathrm{๐}\mathbf{\hspace{0.17em}10}_๐\mathbf{\hspace{0.17em}45}_๐/M`$
$`+g_{12}\mathbf{\hspace{0.17em}16}_\mathrm{๐}\mathbf{\hspace{0.17em}16}_\mathrm{๐}\mathbf{\hspace{0.17em}16}_๐\mathbf{\hspace{0.17em}16}_๐/M\}.`$
Here, $`M`$ could plausibly be of order string scale. Note that a mass matrix having essentially the form of Eq. (21) results if the first term $`h_{33}10_H`$ is dominant. This ensures $`m_b^0m_\tau ^0`$ and $`m_t^0m(\nu _{Dirac})^0`$. Following the assumption of progressive hierarchy (equivalently appropriate flavor symmetries<sup>4</sup><sup>4</sup>4Although no explicit string solution with the hierarchy in $`h_{ij}`$ mentioned above, together with the $`a_{ij}`$ and $`g_{ij}`$ couplings of Eq. (25), exists as yet, flavor symmetries of the type alluded to, as well as SM singlets carrying flavor-charges and acquiring VEVs of order $`M_X`$ that can lead to effective hierarchical couplings, do emerge in string solutions. And, there exist solutions with top Yukawa coupling being leading (see e.g. Refs. and ).), we presume that $`h_{23}h_{33}/10`$, while $`h_{22}`$ and $`h_{11}`$, which are set to be zeros, are progressively much smaller than $`h_{23}`$ (see discussion in Ref. ). Since $`45_H16_HM_X`$, while $`MM_{st}10M_X`$, the terms $`a_{23}45_H/M`$ and $`g_{23}16_H/M`$ can quite plausibly be of order $`h_{33}/10`$, if $`a_{23}g_{23}h_{33}`$. By the assumption of hierarchy, we presume that $`a_{12}a_{23}`$, and $`g_{12}g_{23}`$
It is interesting to observe the symmetry properties of the $`a_{23}`$ and $`g_{23}`$-terms. Although $`10_H\times 45_H=10+120+320`$, given that $`45_H`$ is along $`B`$-$`L`$, which is needed to implement doublet-triplet splitting (see Appendix), only 120 in the decomposition contributes to the mass-matrices. This contribution is, however, antisymmetric in the family-index and, at the same time, proportional to $`B`$-$`L`$. Thus the $`a_{23}`$ term fulfills the requirements of both antisymmetry and ($`B`$-$`L`$)-dependence, simultaneously<sup>5</sup><sup>5</sup>5The analog of $`10_H45_H`$ for the case of G(224) would be $`\chi _H(2,2,1)_H(1,1,15)_H`$. Although in general, the coupling of $`\chi _H`$ to the fermions need not be antisymmetric, for a string-derived G(224), the multiplet (1,1,15$`)_H`$ is most likely to arise from an underlying 45 of SO(10) (rather than 210); in this case, the couplings of $`\chi _H`$ must be antisymmetric like that of $`10_H45_H`$.. With only $`h_{ij}`$ and $`a_{ij}`$-terms, however, the up and down quark mass-matrices will be proportional to each other, which would yield $`V_{CKM}=1`$. This is remedied by the $`g_{ij}`$ coupling. Because, the $`16_H`$ can have a VEV not only along its SM singlet component (transforming as $`\stackrel{~}{\overline{\nu }}_R`$) which is of GUT-scale, but also along its electroweak doublet component โ call it $`16_d`$ โ of the electroweak scale. The latter can arise by the the mixing of $`16_d`$ with the corresponding doublet (call it $`10_d`$) in the $`10_H`$. The MSSM doublet $`H_d`$, which is light, is then a mixture of $`10_d`$ and $`16_d`$, while the orthogonal combination is superheavy (see Appendix). Since $`16_d`$ contributes only to the down-flavor mass matrices, but not to the up-flavor, the $`g_{23}`$ and $`g_{12}`$ couplings generate non-trivial CKM-mixings. We thus see that the minimal Higgs system satisfies apriori all the qualitative requirements (2)-(4), including the condition of $`V_{CKM}1`$. I now discuss that this system works well even quantitatively.
With these six effective Yukawa couplings, the Dirac mass matrices of quarks and leptons of the three families at the unification scale take the form :
$`U=\left(\begin{array}{ccc}0& ฯต^{}& 0\\ ฯต^{}& 0& ฯต+\sigma \\ 0& ฯต+\sigma & 1\end{array}\right)m_U,D=\left(\begin{array}{ccc}0& ฯต^{}+\eta ^{}& 0\\ ฯต^{}+\eta ^{}& 0& ฯต+\eta \\ 0& ฯต+\eta & 1\end{array}\right)m_D,`$ (32)
$`N=\left(\begin{array}{ccc}0& \mathrm{\hspace{0.17em}3}ฯต^{}& 0\\ 3ฯต^{}& 0& \mathrm{\hspace{0.17em}3}ฯต+\sigma \\ 0& 3ฯต+\sigma & 1\end{array}\right)m_U,L=\left(\begin{array}{ccc}0& \mathrm{\hspace{0.17em}3}ฯต^{}+\eta ^{}& 0\\ 3ฯต^{}+\eta ^{}& 0& \mathrm{\hspace{0.17em}3}ฯต+\eta \\ 0& 3ฯต+\eta & 1\end{array}\right)m_D.`$ (39)
Here the matrices are multiplied by left-handed fermion fields from the left and by antiโfermion fields from the right. $`(U,D)`$ stand for the mass matrices of up and down quarks, while $`(N,L)`$ are the Dirac mass matrices of the neutrinos and the charged leptons. The entries $`1,ฯต`$,and $`\sigma `$ arise respectively from the $`h_{33},a_{23}`$ and $`h_{23}`$ terms in Eq. (25), while $`\eta `$ entering into $`D`$ and $`L`$ receives contributions from both $`g_{23}`$ and $`h_{23}`$; thus $`\eta \sigma `$. Similarly $`\eta ^{}`$ and $`ฯต^{}`$ arise from $`g_{12}`$ and $`a_{12}`$ terms respectively. Note the quark-lepton correlations between $`U`$ and $`N`$ as well as $`D`$ and $`L`$, and the up-down correlations between $`U`$ and $`D`$ as well as $`N`$ and $`L`$. These correlations arise because of the symmetry property of G(224). The relative factor of $`3`$ between quarks and leptons involving the $`ฯต`$ entry reflects the fact that $`\mathrm{๐๐}_๐(BL)`$, while the antisymmetry in this entry arises from the group structure of SO(10), as explained above<sup>4</sup>. As we will see, this $`ฯต`$-entry helps to account for (a) the differences between $`m_s`$ and $`m_\mu `$, (b) that between $`m_d`$ and $`m_e`$, and also, (c) the suppression of $`V_{cb}`$ together with the enhancement of the $`\nu _\mu `$-$`\nu _\tau `$ oscillation angle.
The mass matrices in Eq. (39) contain 7 parameters<sup>6</sup><sup>6</sup>6Of these, $`m_U^0m_t^0`$ can in fact be estimated to within $`20\%`$ accuracy by either using the argument of radiative electroweak symmetry breaking, or some promising string solutions (see e.g. Ref. ). : $`ฯต`$, $`\sigma `$, $`\eta `$, $`m_D=h_{33}10_d`$, $`m_U=h_{33}10_U`$, $`\eta ^{}`$ and $`ฯต^{}`$. These may be determined by using, for example, the following input values: $`m_t^{phys}=174`$ GeV, $`m_c(m_c)=1.37`$ GeV, $`m_s(1`$ GeV$`)=110`$-$`116`$ MeV , $`m_u(1`$ GeV) $`6`$ MeV and the observed masses of $`e`$, $`\mu `$ and $`\tau `$, which lead to (see Ref. , for details) :
$`\sigma \mathrm{\hspace{0.17em}\hspace{0.17em}0.110},\eta \mathrm{\hspace{0.17em}0.151},ฯต\mathrm{\hspace{0.17em}0.095},|\eta ^{}|4.4\times 10^3\text{and}ฯต^{}2\times 10^4`$
$`m_Um_t(M_U)(100\text{-}120)\text{GeV},m_Dm_b(M_U)\mathrm{\hspace{0.17em}1.5}\text{GeV}.`$ (40)
We have assumed for simplicity that the parameters are real, because a good fitting suggests that the relative phases of at least $`\sigma `$, $`\eta `$ and $`ฯต`$ are small ($`<10^{}`$ say).Such fitting also fixes their relative signs. Note that in accord with our general expectations discussed above, each of the parameters $`\sigma `$, $`\eta `$ and $`ฯต`$ are found to be of order 1/10, as opposed to being <sup>7</sup><sup>7</sup>7This is one characteristic difference between our work and that of Ref. , where the (2,3)-element is even bigger than the (3,3). $`O(1)`$ or $`O(10^2)`$, compared to the leading (3,3)-element in Eq. (39). Having determined these parameters, we are led to a total of five predictions involving only the quarks (those for the leptons are listed separately) :
$`m_b^0m_\tau ^0(1\mathrm{\hspace{0.17em}8}ฯต^2);\text{thus}m_b(m_b)(4.6\text{-}4.9)\text{GeV}`$ (41)
$`|V_{cb}||\sigma \eta |\left|\sqrt{m_s/m_b}\left|{\displaystyle \frac{\eta +ฯต}{\eta ฯต}}\right|^{1/2}\sqrt{m_c/m_t}\left|{\displaystyle \frac{\sigma +ฯต}{\sigma ฯต}}\right|^{1/2}\right|\mathrm{\hspace{0.17em}0.045}`$ (42)
$`m_d(1\text{GeV})\mathrm{\hspace{0.17em}8}\text{MeV}`$ (43)
$`\theta _C\left|\sqrt{m_d/m_s}e^{i\varphi }\sqrt{m_u/m_c}\right|`$ (44)
$`|V_{ub}/V_{cb}|\sqrt{m_u/m_c}\mathrm{\hspace{0.17em}0.07}.`$ (45)
In making these predictions, we have extrapolated the GUT-scale values down to low energies using $`\alpha _3(m_Z)=0.118`$, a SUSY threshold of 500 GeV and $`\mathrm{tan}\beta =5`$. The results depend weakly on these choices, assuming $`\mathrm{tan}\beta 2`$-30. Further, the Dirac masses and mixings of the neutrinos and the mixings of the charged leptons also get determined. We obtain :
$`m_{\nu _\tau }^D(M_U)\mathrm{\hspace{0.17em}100}\text{-}120\text{GeV};m_{\nu _\mu }^D(M_U)\mathrm{\hspace{0.17em}8}\text{GeV},`$ (46)
$`\theta _{\mu \tau }^{\mathrm{}}\mathrm{\hspace{0.17em}3}ฯต+\eta \sqrt{m_\mu /m_\tau }\left|{\displaystyle \frac{\mathrm{\hspace{0.17em}3}ฯต+\eta }{3ฯต+\eta }}\right|^{1/2}\mathrm{\hspace{0.17em}0.437}`$ (47)
$`m_{\nu _e}^D[\mathrm{\hspace{0.17em}9}ฯต^{}_{}{}^{}2/(9ฯต^2\sigma ^2)]m_U\mathrm{\hspace{0.17em}0.4}\text{MeV}`$ (48)
$`\theta _{e\mu }^{\mathrm{}}\left|{\displaystyle \frac{\eta ^{}\mathrm{\hspace{0.17em}3}ฯต^{}}{\eta ^{}+\mathrm{\hspace{0.17em}3}ฯต^{}}}\right|^{1/2}\sqrt{m_e/m_\mu }\mathrm{\hspace{0.17em}0.85}\sqrt{m_e/m_\mu }\mathrm{\hspace{0.17em}0.06}`$ (49)
$`\theta _{e\tau }^{\mathrm{}}{\displaystyle \frac{1}{0.85}}\sqrt{m_e/m_\tau }(m_\mu /m_\tau )\mathrm{\hspace{0.17em}0.0012}.`$ (50)
In evaluating $`\theta _{e\mu }^{\mathrm{}}`$, we have assumed $`ฯต^{}`$ and $`\eta ^{}`$ to be relatively positive.
Given the bizarre pattern of quark and lepton masses and mixings, it seems remarkable that the simple pattern of fermion mass-matrices, motivated by the group theory of G(224)/SO(10), gives an overall fit to all of them which is good to within $`10\%`$. This includes the two successful predictions on $`m_b`$ and $`V_{cb}`$ (Eqs.(41 and (42)). Note that in supersymmetric unified theories, the โobservedโ value of $`m_b(m_b)`$ and renormalization-group studies suggest that, for a wide range of the parameter $`\mathrm{tan}\beta `$, $`m_b^0`$ should in fact be about 10-20$`\%`$ lower than $`m_\tau ^0`$ . This is neatly explained by the relation: $`m_b^0m_\tau ^0(18ฯต^2)`$ (Eq. (41)), where exact equality holds in the limit $`ฯต0`$ (due to SU(4)-color), while the decrease of $`m_b^0`$ compared to $`m_\tau ^0`$ by $`8ฯต^210\%`$ is precisely because the off-diagonal $`ฯต`$-entry is proportional to $`B`$-$`L`$ (see Eq. (39)).
Specially intriguing is the result on $`V_{cb}0.045`$ which compares well with the observed value of $`0.04`$. The suppression of $`V_{cb}`$, compared to the value of $`0.17\pm 0.06`$ obtained from Eq. (23), is now possible because the mass matrices (Eq. (39)) contain an antisymmetric component $`ฯต`$. That corrects the square-root formula $`\theta _{sb}=\sqrt{m_s/m_b}`$ (appropriate for symmetric matrices, see Eq. (21)) by the asymmetry factor $`|(\eta +ฯต)/(\eta ฯต)|^{1/2}`$ (see Eq. (19)), and similarly for the angle $`\theta _{ct}`$. This factor suppresses $`V_{cb}`$ if $`\eta `$ and $`ฯต`$ have opposite signs. The interesting point is that, the same feature necessarily enhances the corresponding mixing angle $`\theta _{\mu \tau }^{\mathrm{}}`$ in the leptonic sector, since the asymmetry factor in this case is given by $`[(3ฯต+\eta )/(3ฯต+\eta )]^{1/2}`$ (see Eq. (24)). This enhancement of $`\theta _{\mu \tau }^{\mathrm{}}`$ helps to account for the nearly maximal oscillation angle observed at SuperK (as discussed below). This intriguing correlation between the mixing angles in the quark versus leptonic sectors โ that is suppression of one implying enhancement of the other โ has become possible only because of the $`ฯต`$-contribution, which is simultaneously antisymmetric and is proportional to $`B`$-$`L`$. That in turn becomes possible because of the group-property of SO(10) or a string-derived G(224)<sup>4</sup>.
Taking stock, we see an overwhelming set of evidences in favor of $`B`$-$`L`$ and in fact for the full SU(4)-color-symmetry. These include: (i) the suppression of $`V_{cb}`$, together with the enhancement of $`\theta _{\mu \tau }^{\mathrm{}}`$, just mentioned above, (ii) the successful relation $`m_b^0m_\tau ^0(18ฯต^2)`$, (iii) the usefulness again of the SU(4)-color-relation $`m(\nu _{Dirac}^\tau )^0m_t^0`$ in accounting for $`m(\nu _L^\tau )`$( see Sec. 4 ), and (iv) the agreement of the relation $`|m_s^0/m_\mu ^0|=|(ฯต^2\eta ^2)/(9ฯต^2\eta ^2)|`$ with the data, in that the ratio is naturally less than 1, if $`\eta ฯต`$. The presence of $`9ฯต^2`$ in the denominator is because the off-diagonal entry is proportional to B-L. Finally, the need for ($`B`$-$`L`$)- as a local symmetry, to implement baryogenesis, has been noted in Sec.1.
Turning to neutrino masses, while all the entries in the Dirac mass matrix $`N`$ are now fixed, to obtain the parameters for the light neutrinos, one needs to specify those of the Majorana mass matrix of the RH neutrinos ($`\nu _R^{e,\mu ,\tau }`$). Guided by economy and the assumption of hierarchy, we consider the following pattern :
$`M_\nu ^R=\left(\begin{array}{ccc}x& 0& z\\ 0& 0& y\\ z& y& 1\end{array}\right)M_R.`$ (54)
As discussed in Sec. 4, the magnitude of $`M_R(5\text{-}15)\times 10^{14}`$ GeV can quite plausibly be justified in the context of supersymmetric unificaton<sup>8</sup><sup>8</sup>8This estimate for $`M_R`$ is retained even if one allows for $`\nu _\mu \text{-}\nu _\tau `$ mixing (see Ref. ). (e.g. by using $`MM_{st}4\times 10^{17}`$ GeV in Eq. (16)). To the same extent, the magnitude of $`m(\nu _\tau )(1/10\text{-}1/30)`$ eV, which is consistent with the SuperK value, can also be anticipated. Thus there are effectively three new parameters: $`x`$, $`y`$, and $`z`$. Since there are six observables for the three light neutrinos, one can expect three predictions. These may be taken to be $`\theta _{\nu _\mu \nu _\tau }^{osc}`$, $`m_{\nu _\tau }`$ (see Eq. (18)), and for example $`\theta _{\nu _e\nu _\mu }^{osc}`$.
Assuming successively hierarchical entries as for the Dirac mass matrices, we presume that $`|y|1/10,|z||y|/10`$ and $`|x|z^2`$. Now given that $`m(\nu _\tau )1/20`$ eV (as estimated in Eq. (18)), the MSW solution for the solar neutrino puzzle suggests that $`m(\nu _\mu )/m(\nu _\tau )1/10\text{-}1/30`$. The latter in turn yields : $`|y|(1/18\text{ to }1/23.6)`$, with $`y`$ having the same sign as $`ฯต`$ (see Eq. (40)). This solution for y obtains only by assuming that $`y`$ is $`O(1/10)`$ rather than $`O(1)`$. Combining now with the mixing in the $`\mu `$-$`\tau `$ sector determined above (see Eq. (47)), one can then determine the $`\nu _\mu \text{-}\nu _\tau `$ oscillation angle. The two predictions of the model for the neutrino-system are then :
$`m(\nu _\tau )(1/10\text{-}\mathrm{\hspace{0.17em}1}/30)\text{eV}`$ (55)
$`\theta _{\nu _\mu \nu _\tau }^{osc}\theta _{\mu \tau }^{\mathrm{}}\theta _{\mu \tau }^\nu \left(0.437+\sqrt{{\displaystyle \frac{m_{\nu _2}}{m_{\nu _3}}}}\right).`$ (56)
$`\text{Thus,}\mathrm{sin}^2\mathrm{\hspace{0.17em}2}\theta _{\nu _\mu \nu _\tau }^{osc}=(0.96,0.91,0.86,0.83,0.81)`$ (57)
$`\text{for}m_{\nu _2}/m_{\nu _3}=(1/10,1/15,1/20,1/25,1/30).`$ (58)
Both of these predictions are extremely successful.
Note the interesting point that the MSW solution, together with the requirement that $`|y|`$ should have a natural hierarchical value (as mentioned above), lead to $`y`$ having the same sign as $`ฯต`$; that (it turns out) implies that the two contributions in Eq.(56) must add rather than subtract, leading to an almost maximal oscillation angle . The other factor contributing to the enhancement of $`\theta _{\nu _\mu \nu _\tau }^{osc}`$ is, of course, also the asymmetry-ratio which increases $`|\theta _{\mu \tau }^{\mathrm{}}|`$ from 0.25 to 0.437 (see Eq. (47)). We see that one can derive rather plausibly a large $`\nu _\mu \text{-}\nu _\tau `$ oscillation angle $`\mathrm{sin}^2\mathrm{\hspace{0.17em}2}\theta _{\nu _\mu \nu _\tau }^{osc}0.8`$, together with an understanding of hierarchical masses and mixings of the quarks and the charged leptons, while maintaining a large hierarchy in the seesaw derived masses ($`m_{\nu _2}/m_{\nu _3}=1/10\text{-}1/30`$), all within a unified framework including both quarks and leptons. In the example exhibited here, the mixing angles for the mass eigenstates of neither the neutrinos nor the charged leptons are really large,in that $`\theta _{\mu \tau }^{\mathrm{}}0.43723^{}`$ and $`\theta _{\mu \tau }^\nu (0.18\text{-}0.31)(10\text{-}18)^{}`$, yet the oscillation angle obtained by combining the two is near-maximal. This contrasts with most previous work, in which a large oscillation angle is obtained either entirely from the neutrino sector (with nearly degenerate neutrinos) or almost entirely from the charged lepton sector.
While $`M_R(5\text{-}15)\times 10^{14}`$ GeV and $`y1/20`$ are better determined, the parameters $`x`$ and $`z`$ can not be obtained reliably at present because very little is known about observables involving $`\nu _e`$. Taking, for concreteness, $`m_{\nu _e}(10^5\text{-}10^4`$ (1 to few)) eV and $`\theta _{e\tau }^{osc}\theta _{e\tau }^{\mathrm{}}\theta _{e\tau }^\nu 10^3\pm 0.03`$ as inputs, we obtain : $`z(1`$-$`5)\times 10^3`$ and $`x(`$1 to few)$`(10^6\text{-}10^5)`$, in accord with the guidelines of $`|z||y|/10`$ and $`|x|z^2`$. This in turn yields : $`\theta _{e\mu }^{osc}\theta _{e\mu }^{\mathrm{}}\theta _{e\mu }^\nu 0.06\pm 0.015`$. Note that the mass of $`m_{\nu _\mu }3\times 10^3`$ eV, that follows from a natural hierarchical value for $`y(1/20)`$, and $`\theta _{e\mu }`$ as above, go well with the small angle MSW explanation<sup>9</sup><sup>9</sup>9Although the small angle MSW solution appears to be more generic within the approach outlined above, we have found that the large angle solution can still plausibly emerge in a limited region of parameter space, without affecting our results on fermion masses. of the solar neutrinos puzzle.
It is worthnoting that although the superheavy Majorana masses of the RH neutrinos cannot be observed directly, they can be of cosmological significance. The pattern given above and the arguments given in Sec. 3 and in this section suggests that $`M(\nu _R^\tau )(5\text{-}15)\times 10^{14}`$ GeV, $`M(\nu _R^\mu )(1\text{-}4)\times 10^{12}`$ GeV (for $`x1/20`$); and $`M(\nu _R^e)(1/2\text{-}10)\times 10^9`$ GeV (for $`x(1/2\text{-}10)10^6>z^2`$). A mass of $`\nu _R^e10^9`$ GeV is of the right magnitude for producing $`\nu _R^e`$ following reheating and inducing lepton asymmetry in $`\nu _R^e`$ decay into $`H^0+\nu _L^i`$, that is subsequently converted into baryon asymmetry by the electroweak sphalerons .
In summary, we have proposed an economical and predictive pattern for the Dirac mass matrices, within the SO(10)/G(224)-framework, which is remarakbly successful in describing the observed masses and mixings of all the quarks and charged leptons. It leads to five predictions for just the quark- system, all of which agree with observation to within 10%. The same pattern, supplemented with a similar structure for the Majorana mass matrix, accounts for both the large $`\nu _\mu `$-$`\nu _\tau `$ oscillation angle and a mass of $`\nu _\tau 1/20`$ eV, suggested by the SuperK data. It also accomodates a small $`\nu _e`$-$`\nu _\mu `$ oscillation angle relevant for theories of the solar neutrino deficit. Given this degree of success, it makes good sense to study proton decay concretely within this SO(10)/G(224)-framework. The results of this study are presented in the next section.
Before turning to proton decay, it is worth noting that much of our discussion of fermion masses and mixings, including those of the neutrinos, is essentially unaltered if we go to the limit $`ฯต^{}0`$ of Eq. (28). This limit clearly involves:
$`m_u=\mathrm{\hspace{0.17em}0},\theta _C\sqrt{m_d/m_s},m_{\nu _e}=\mathrm{\hspace{0.17em}0},\theta _{e\mu }^\nu =\theta _{e\tau }^\nu =\mathrm{\hspace{0.17em}0}.`$
$`|V_{ub}|\sqrt{{\displaystyle \frac{\eta ฯต}{\eta +ฯต}}}\sqrt{m_d/m_b}(m_s/m_b)(2.1)(0.039)(0.023)\mathrm{\hspace{0.17em}0.0019}`$ (59)
All other predictions remain unaltered. Now, among the observed quantities in the list above, $`\theta _C\sqrt{m_d/m_s}`$ is a good result. Considering that $`m_u/m_t10^5`$, $`m_u=0`$ is also a pretty good result. There are of course plausible small corrections which could arise through Planck scale physics; these could induce a small value for $`m_u`$ through the (1,1)-entry $`\delta 10^5`$. For considerations of proton decay, it is worth distinguishing between these two variants, which we will refer to as cases I and II respectively.
Case I : $`ฯต^{}\mathrm{\hspace{0.17em}2}\times \mathrm{\hspace{0.17em}10}^4,\delta =\mathrm{\hspace{0.17em}0}`$
Case II : $`\delta \mathrm{\hspace{0.17em}10}^5,ฯต^{}=0.`$ (60)
## 6 Expectations for Proton Decay in Supersymmetric Unified Theories
6.1 Turning to the main purpose of this talk, I present now the reason why the unification framework based on SUSY SO(10) or G(224), together with the understanding of fermion masses and mixings discussed above, strongly suggest that proton decay should be imminent.
Recall that supersymmetric unfied theories (GUTs) introduce two new features to proton decay : (i) First, by raising $`M_X`$ to a higher value of about $`2\times 10^{16}`$ GeV, they strongly suppress the gauge-boson-mediated $`d=6`$ proton decay operators, for which $`e^+\pi ^0`$ would have been the dominant mode (for this case, one typically obtains : $`\mathrm{\Gamma }^1(pe^+\pi ^0)|_{d=6}10^{36\pm 1.5}`$ yrs). (ii) Second, they generate $`d=5`$ proton decay operators of the form $`Q_iQ_jQ_kQ_l/M`$ in the superpotential, through the exchange of color triplet Higginos, which are the GUT partners of the standard Higgs(ino) doublets, such as those in the $`5+\overline{5}`$ of SU(5) or the 10 of SO(10). Assuming that a suitable doublet-triplet splitting mecahnism provides heavy GUT-scale masses to these color triplets and light masses to the doublets, these โstandardโ $`d=5`$ operators, suppressed by just one power of the heavy mass and the small Yukawa couplings, are found to provide the dominant mechanism for proton decay in supersymmetric GUT .
Now, owing to (a) Bose symmetry of the superfields in $`QQQL/M`$, (b) color antisymmetry, and especially (c) the hierarchical Yukawa couplings of the Higgs doublets, it turns out that these standard $`d=5`$ operators lead to dominant $`\overline{\nu }K^+`$ and comparable $`\overline{\nu }\pi ^+`$ modes, but in all cases to highly suppressed $`e^+\pi ^0`$, $`e^+K^0`$ and even $`\mu ^+K^0`$ modes. For instance, for minimal SUSY SU(5), one obtains (with $`\mathrm{tan}\beta 20`$, say) :
$`[\mathrm{\Gamma }(\mu ^+K^0)/\mathrm{\Gamma }(\overline{\nu }K^+)]_{std}^{SU(5)}[m_u/m_c\mathrm{sin}^2\theta ]R\mathrm{\hspace{0.17em}10}^3,`$ (61)
where $`R0.1`$ is the ratio of the relevant $`|`$matrix element$`|^2\times `$(phase space), for the two modes.
It was recently pointed out that in SUSY unified theories based on SO(10) or G(224), which assign heavy Majorana masses to the RH neutrinos, there exists a new set of color triplets and thereby very likely a new source of $`d=5`$ proton decay operators . For instance, in the context of the minimal set of Higgs multiplets<sup>10</sup><sup>10</sup>10The origin of the new $`d=5`$ operators in the context of other Higgs multiplets, in particular in the cases where $`126_H`$ and $`\overline{126}_H`$ are used to break $`B`$-$`L`$, has been discussed in Ref. . $`\{45_H,16_H,\overline{16}_H`$ and $`10_H\}`$ (see Sec. 5), these new $`d=5`$ operators arise by combining three effective couplings introduced before :โ i.e., (a) the couplings $`f_{ij}16i16j\overline{16}_H\overline{16}_H/M`$ (see Eq. (15)) that are required to assign Majorana masses to the RH neutrinos, (b) the couplings $`g_{ij}16i16j16_H16_H/M`$, which are needed to generate non-trivial CKM mixings (see Eq. (25)), and (c) the mass term $`M_{16}16_H\overline{16}_H`$. For the $`f_{ij}`$ couplings, there are two possible SO(10)-contractions, and we assume both to have comparable strength<sup>11</sup><sup>11</sup>11One would expect such a general contraction to hold, especially if the $`f_{ij}`$ couplings are induced by non-perturbative quantum gravity. Furthermore, the $`f_{ij}`$ couplings with the contraction of the pair ($`16_i\overline{16}_H`$), being effectively in 45 (rather than in 1) of SO(10), would be induced also by tree-level exchanges, if these pairs couple to the 45โs in the string tower. Such a contraction would lead to proton decay.. In this case, the color-triplet Higginos in $`\overline{16}_H`$ and $`16_H`$ of mass $`M_{16}`$ can be exchanged between $`\stackrel{~}{q}_iq_j`$ and $`\stackrel{~}{q}_kq_l`$-pairs. This exchange generates a new set of $`d=5`$ operators in the effective superpotential of the form
$`W_{new}f_{ij}g_{kl}(16_i\mathrm{\hspace{0.17em}16}_j)(16_k\mathrm{\hspace{0.17em}16}_l)\overline{16}_H16_H/M^2(1/M_{16}),`$ (62)
which induce proton decay. Note that these operators depend, through the couplings $`f_{ij}`$ and $`g_{kl}`$, both on the Majorana and on the Dirac masses of the respective fermions. This is why within SUSY SO(10) or G(224), proton decay gets intimately linked to the masses and mixings of all fermions, including neutrinos.
### 6.2 Framework for Calculating Proton Decay Rate
To establish notations, consider the case of minimal SUSY SU(5) and, as an example, the process $`\stackrel{~}{c}\stackrel{~}{d}\overline{s}\overline{\nu }_\mu `$, which induces $`p\overline{\nu }_\mu K^+`$. Let the strength of the corresponding $`d=5`$ operator, multiplied by the product of the CKM mixing elements entering into wino-exchange vertices, (which in this case is $`\mathrm{sin}\theta _C\mathrm{cos}\theta _C)`$ be denoted by $`\widehat{A}`$. Thus (putting $`\mathrm{cos}\theta _C=1`$), one obtains:
$`\widehat{A}_{\stackrel{~}{c}\stackrel{~}{d}}(SU(5))=(h_{22}^uh_{12}^d/M_{H_C})\mathrm{sin}\theta _c(m_cm_s\mathrm{sin}^2\theta _C/v_u^2)(\mathrm{tan}\beta /M_{H_C})`$ (63)
$`(1.9\times 10^8)(\mathrm{tan}\beta /M_{H_C})(2\times 10^{24}\text{GeV}\text{-1})(\mathrm{tan}\beta /2)(2\times 10^{16}\text{GeV}/M_{H_C}),`$
where $`\mathrm{tan}\beta v_u/v_d`$, and we have put $`v_u=174`$ GeV and the fermion masses extrapolated to the unification-scale โ i.e. $`m_c300`$ MeV and $`m_s40`$ MeV. The amplitude for the associated four-fermion process $`dus\overline{\nu }_\mu `$ is given by:
$`A_5(dus\overline{\nu }_\mu )=\widehat{A}_{\stackrel{~}{c}\stackrel{~}{d}}\times (2f)`$ (64)
where $`f`$ is the loop-factor associated with wino-dressing. Assuming $`m_{\stackrel{~}{w}}m_{\stackrel{~}{q}}m_{\stackrel{~}{l}}`$ one gets: $`f(m_{\stackrel{~}{w}}/m_{\stackrel{~}{q}}^2)(\alpha _2/4\pi )`$. Using the amplitude for $`(du)(s\nu _{\mathrm{}}`$), as in Eq. (64), ($`\mathrm{}=\mu `$ or $`\tau `$), one then obtains :
$`\mathrm{\Gamma }^1(p\overline{\nu }_\tau K^+)(2.2\times 10^{31})\text{yrs}\times `$ (65)
$`\left({\displaystyle \frac{0.67}{A_S}}\right)^2\left[{\displaystyle \frac{0.006\text{GeV}^3}{\beta _H}}\right]^2\left[{\displaystyle \frac{(1/6)}{(m_{\stackrel{~}{W}}/m_{\stackrel{~}{q}})}}\right]^2\left[{\displaystyle \frac{m_{\stackrel{~}{q}}}{1\text{TeV}}}\right]^2\left[{\displaystyle \frac{2\times 10^{24}\text{GeV}^1}{\widehat{A}(\overline{\nu })}}\right]^2.`$
Here $`\beta _H`$ denotes the hadronic matrix element defined by $`\beta _Hu_L(\stackrel{}{k})ฯต_{\alpha \beta \gamma }0|(d_L^\alpha u_L^\beta )u_L^\gamma |p,\stackrel{}{k}`$. While the range $`\beta _H=(0.003\text{-}0.03)`$ GeV<sup>3</sup> has been used in the past , given that one lattice calculation yields $`\beta _H=(5.6\pm 0.5)\times 10^3`$ GeV<sup>3</sup> , we will take as a plausible range : $`\beta _H=(0.006`$ GeV<sup>3</sup>)$`(1/22)`$.Here, $`A_S0.67`$ stands for the short distance renormalization factor of the $`d=5`$ operator. Note that the familiar factors that appear in the expression for proton lifetime โ i.e., $`M_{H_C}`$, ($`1+y_{tc}`$) representing the interference between the $`\stackrel{~}{t}`$ and $`\stackrel{~}{c}`$ contributions, and $`\mathrm{tan}\beta `$ (see e.g. Ref.) โ are all effectively contained in $`\widehat{A}(\overline{\nu })`$. Allowing for plausible and rather generous uncertainties in the matrix element and the spectrum we take:
$`\beta _H=(0.006\text{GeV}^3)(1/2\text{-}\mathrm{\hspace{0.17em}2})`$
$`m_{\stackrel{~}{w}}/m_{\stackrel{~}{q}}=\mathrm{\hspace{0.17em}1}/6(1/2\text{-}\mathrm{\hspace{0.17em}2}),\mathrm{and}m_{\stackrel{~}{q}}m_\stackrel{~}{\mathrm{}}\mathrm{\hspace{0.17em}1}\text{TeV}(1/\sqrt{2}\text{-}\sqrt{2}).`$ (66)
Using Eqs. (65-66), we get :
$`\mathrm{\Gamma }^1(p\overline{\nu }_\tau K^+)(2.2\times 10^{31}\text{yrs})[\mathrm{\hspace{0.17em}2.2}\times 10^{24}\text{GeV}^1/\widehat{A}(\overline{\nu }_{\mathrm{}})]^2[32\text{-}\mathrm{\hspace{0.17em}1}/32].`$ (67)
This relation, as well as Eq. (65) are general, depending only on $`\widehat{A}(\overline{\nu }_{\mathrm{}})`$ and on the range of parameters given in Eq. (66). They can thus be used for both SU(5) and SO(10).
The experimental lower limit on the inverse rate for the $`\overline{\nu }K^+`$ modes is given by ,
$`[{\displaystyle \underset{\mathrm{}}{}}\mathrm{\Gamma }(p\overline{\nu }_{\mathrm{}}K^+)]_{expt}^1>1.6\times 10^{33}\text{yrs}.`$ (68)
Allowing for all the uncertainties to stretch in the same direction (in this case, the square bracket = 32), and assuming that just one neutrino flavor (e.g. $`\nu _\mu `$ for SU(5)) dominates, the observed limit Eq. (68) provides an upper bound on the amplitude<sup>12</sup><sup>12</sup>12If there are sub-dominant $`\overline{\nu }_iK^+`$ modes with branching ratio $`R`$, the right side of Eq. (69) should be divided by $`\sqrt{1+R}`$.:
$`\widehat{A}(\overline{\nu }_{\mathrm{}})\sqrt{2}\times 10^{24}\text{GeV}^1`$ (69)
which holds for both SU(5) and SO(10). For minimal SU(5), using Eq. (63) and $`\mathrm{tan}\beta 2`$ (which is suggested on several grounds), one obtains a lower limit on $`M_{HC}`$ given by:
$`M_{HC}\mathrm{\hspace{0.17em}3}\times 10^{16}\text{GeV}(\text{SU}(5))`$ (70)
At the same time, higher values of $`M_{HC}>3\times 10^{16}`$ GeV do not go very well with gauge coupling unification. Thus keeping $`M_{HC}3\times 10^{16}`$ and $`\mathrm{tan}\beta 2`$, we obtain from Eq. (63): $`\widehat{A}(\text{SU}(5))(4/3)\times 10^{24}\text{GeV}^1`$. Using Eq. (67), this in turn implies that
$`\mathrm{\Gamma }^1(p\overline{\nu }K^+)\mathrm{\hspace{0.17em}1.5}\times 10^{33}\text{yrs}(\text{SU}(5))`$ (71)
This is a conservative upper limit. In practise, it is unlikely that all the uncertainties, including that in $`M_{HC}`$, would stretch in the same direction to nearly extreme values so as to prolong proton lifetime. A more reasonable upper limit, for minimal SU(5), thus seems to be: $`\mathrm{\Gamma }^1(p\overline{\nu }K^+)(SU(5))(0.7)\times 10^{33}`$ yrs. Given the experimental lower limit (Eq. (68)), we see that minimal SUSY SU(5) is already or almost on the verge of being excluded by proton decay-searches. We have of course noted in Sec. 4 that SUSY SU(5) does not go well with neutrino oscillations observed at SuperK.
Now, to discuss proton decay in the context of supersymmetric SO(10), it is necessary to discuss first the mechanism for doublet-triplet splitting. Details of this discussion may be found in Ref. . A synopsis is presented in the appendix.
### 6.3. Proton Decay in Supersymmetric SO(10)
The calculation of the amplitudes $`\widehat{A}_{std}`$ and $`\widehat{A}_{new}`$ for the standard and the new operators for the SO(10) model, are given in detail in Ref. . Here, I will present only the vresults. It is found that the four amplitudes $`\widehat{A}_{std}(\overline{\nu }_\tau K^+)`$, $`\widehat{A}_{std}(\overline{\nu }_\mu K^+)`$, $`\widehat{A}_{new}(\overline{\nu }_\tau K^+)`$ and $`\widehat{A}_{new}(\overline{\nu }_\mu K^+)`$ are in fact very comparable to each other, within about a factor of two, either way. Since there is no reason to expect a near cancellation between the standard and the new operators, especially for both $`\overline{\nu }_\tau K^+`$ and $`\overline{\nu }_\mu K^+`$ modes, we expect the net amplitude (standard + new) to be in the range exhibited by either one. Following Ref. , I therefore present the contributions from the standard and the new operators separately. Using the upper limit on $`M_{eff}3\times 10^{18}`$ GeV (see Appendix), we obtain a lower limit for the standard proton decay amplitude given by
$`\widehat{A}(\overline{\nu }_\tau K^+)_{std}\left[\begin{array}{c}(7\times 10^{24}\text{GeV}^1)(1/6\text{-\hspace{0.17em}1/4)}\text{case I}\\ (3\times 10^{24}\text{GeV}^1)(1/6\text{-\hspace{0.17em}1/2)}\text{case II}\end{array}\right]`$ (74)
Substituting into Eq. (67) and adding the contribution from the second competing mode $`\overline{\nu }_\mu K^+`$, with a typical branching ratio $`R0.3`$, we obtain
$`\mathrm{\Gamma }^1(\overline{\nu }K^+)_{std}\left[\begin{array}{c}(3\times 10^{31}\text{yrs.})(1.6\text{-}\mathrm{\hspace{0.17em}0.7})\\ (6.8\times 10^{31}\text{yrs.})(4\text{-}\mathrm{\hspace{0.17em}0.44})\end{array}\right](32\text{-}\mathrm{\hspace{0.17em}1}/32)`$ (77)
The upper and lower entries in Eqs. (74) and (77) correspond to the cases I and II of the fermion mass-matrix - i.e. $`ฯต^{}0`$ and $`ฯต^{}=0`$ \- respectively, (see Eq. (60)). The uncertainty shown inside the square brackets correspond to that in the relative phases of the different contributions. The uncertainty of (32 to 1/32) arises from that in $`\beta _H`$, $`(m_{\stackrel{~}{W}}/m_{\stackrel{~}{q}})`$ and $`m_{\stackrel{~}{q}}`$ (see Eq. (66)). Thus we find that for MSSM embedded in SO(10), the inverse partial proton decay rate should satisfy :
$`\mathrm{\Gamma }^1(p\overline{\nu }K^+)_{std}\left[\begin{array}{c}3\times 10^{31\pm 1.7}\text{yrs.}\\ 6.8\times 10^{31_{1.5}^{+2.1}}\text{y}rs.\end{array}\right]\left[\begin{array}{c}1.5\times 10^{33}\text{yrs.}\\ 7\times 10^{33}\text{yrs.}\end{array}\right](\text{SO(10)}).`$ (82)
The central value of the upper limit in Eq. (82) corresponds to taking the upper limit on $`M_{eff}`$. The uncertainties of matrix element and spectrum are reflected in the exponents.The uncertainity in the most sensitive entry of the fermion mass matrix - i.e. $`ฯต^{}`$ \- is fully incorporated (as regards obtaining an upper limit on the lifetime) by going from case I to case II . Note that this increases the lifetime by almost a factor of five. Any non-vanishing value of $`ฯต^{}`$ would only shorten the lifetime compared to case II. In this sense, the larger of the two upper limits quoted above is rather conservative.
Evaluating similarly the contributions from only the new operators, we obtain :
$`\mathrm{\Gamma }^1(\overline{\nu }K^+)_{new}(3\times 10^{31}\text{yrs})[16\text{-}\mathrm{\hspace{0.17em}1}/1.7]\{32\text{-}\mathrm{\hspace{0.17em}1}/32\}.`$ (83)
Note that this contribution is independent of $`M_{eff}`$. It turns out that it is also insensitive to $`ฯต^{}`$ ; thus it is nearly the same for cases I and II. Allowing for a net uncertainty at the upper end by as much as a factor of 20 to 200, arising jointly from the square and the curly brackets, i.e., without going to extreme ends of all parameters, the new operators related to neutrino masses, by themselves, lead to a proton decay lifetime bounded by:
$$\mathrm{\Gamma }^1(\overline{\nu }K^+)_{new}^{expected}(0.6\text{-}\mathrm{\hspace{0.17em}6})\times 10^{33}\text{yrs.}\text{(SO(10) or string G(224))}$$
(84)
It should be stressed that while the standard $`d=5`$ operators would be absent for a string-derived G(224)-model, the new $`d=5`$ operators, related to the Majorana masses of the RH neutrinos and the CKM mixings, would still be present for such a model. Thus our expectations for the proton decay lifetime (as shown in Eq. (84)) and the prominence of the $`\mu ^+K^0`$ mode (see below) hold for a string-derived G(224)-model, just as they do for SO(10).
### 6.4. The Charged Lepton Decay Mode $`(p\mu ^+K^0)`$
I now note a distinguishing feature of the SO(10) or the G(224) model presented here. Allowing for uncertainties in the way the standard and the new operators can combine with each other for the three leading modes i.e. $`\overline{\nu }_\tau K^+`$, $`\overline{\nu }_\mu K^+`$ and $`\mu ^+K^0`$, we obtain (see Ref. for details):
$`B(\mu ^+K^0)_{std+new}\left[1\%\text{to}\mathrm{\hspace{0.17em}\hspace{0.17em}50}\%\right]\rho \text{(SO(10) or string G(224))}`$ (85)
where $`\rho `$ denotes the ratio of the squares of relevant matrix elements for the $`\mu ^+K^0`$ and $`\overline{\nu }K^+`$ modes. In the absence of a reliable lattice calculation for the $`\overline{\nu }K^+`$ mode , one should remain open to the possibility of $`\rho 1/2`$ to 1 (say). We find that for a large range of parameters, the branching ratio $`B(\mu ^+K^0)`$ can lie in the range of 20 to 40% (if $`\rho 1`$). This prominence of the $`\mu ^+K^0`$ mode for the SO(10)/G(224) model is primarily due to contributions from the new operators. This contrasts sharply with the minimal SU(5) model, in which the $`\mu ^+K^0`$ mode is expected to have a branching ratio of only about $`10^3`$. In short, prominence of the $`\mu ^+K^0`$ mode, if seen, would clearly show the relevance of the new operators, and thereby reveal the proposed link between neutrino masses and proton decay .
### 6.5. Section Summary
In summary, our study of proton decay has been carried out within the SO(10) or the G(224)-framework<sup>13</sup><sup>13</sup>13As described in Secs. 3 and 5., with special attention paid to its dependence on fermion masses and threshold effects. The study strongly suggests an upperlimit on proton lifetime, given by
$`\tau _{proton}(1/2\text{-}\mathrm{\hspace{0.17em}1})\times 10^{34}\text{yrs},`$ (86)
with $`\overline{\nu }K^+`$ being the dominant decay mode. Although there are uncertainties in the matrix element, in the SUSY-spectrum, and in certain sensitive elements of the fermion mass matrix, especially $`ฯต^{}`$ (see Eq. (82) for predictions in cases I versus II), this upper limit is obtained by allowing for a generous range in these parameters and stretching all of them in the same direction so as to extend proton lifetime. In this sense, while the predicted lifetime spans a wide range, the upper limit quoted above is quite conservative. In turn, it provides a clear reason to expect that the discovery of proton decay should be imminent. The implication of this prediction for a next-generation detector is emphasized in the next section.
## 7 Concluding Remarks
The preceding sections show that one is now in possession of a set of facts, which may be viewed as the matching pieces of a puzzle ; in that all of them can be resolved by just one idea - that is grand unification. These include : (i) the observed family-structure, (ii) meeting of the three gauge coulings, (iii) neutrino oscillations; in particular the mass of $`\nu _\tau `$ (suggested by SuperK), (iv) the intricate pattern of the masses and mixings of all the fermions, including the smallness of $`V_{bc}`$ and the largeness of $`\theta _{\nu _\mu \nu _\tau }^{osc}`$, and (v) the need for $`B`$-$`L`$ to implement baryogenesis. All these pieces fit beautifully together within a single puzzle board framed by supersymmetric unification, based on SO(10) or a string-unified G(224)-symmetry.
The one and the most notable piece of the puzzle still missing, however, is proton decay. Based on a systematic study of this process within the supersymmetric SO(10)/G(224)-framework , which is clearly favored by the data, I have argued here that a conservative upper limit on the proton lifetime is about (1/2 - 1)$`\times 10^{34}`$ yrs. So, unless the fitting of all the pieces listed above is a mere coincidence, and I believe that that is highly unlikely, discovery of proton decay should be around the corner. In particular, as mentioned in the Introduction, we expect that candidate events should be observed in the near future already at SuperK. However, allowing for the possibility that proton lifetime may well be near the upper limit stated above, a next-generation detector providing a net gain in sensitivity by a factor five to ten, compared to SuperK, would be needed to produce real events and distinguish them unambiguously from the background. Such an improved detector would of course be essential to study the branching ratios of certain crucial though sub-dominant decay modes such as the $`\mu ^+K^0`$ .
The reason for pleading for such improved searches is that proton decay would provide us with a wealth of knowledge about physics at truly short distances ($`<10^{30}`$ cm), which cannot be gained by any other means. Specifically, the observation of proton decay, at a rate suggested above, with $`\overline{\nu }K^+`$ mode being dominant, would not only reveal the underlying unity of quarks and leptons but also the relevance of supersymmetry. It would also confirm a unification of the fundamental forces at a scale of order $`2\times 10^{16}`$ GeV. Furthermore, prominence of the $`\mu ^+K^0`$ mode, if seen, would have even deeper significance, in that in addition to supporting the three features mentioned above, it would also reveal the link between neutrino masses and proton decay, as discussed in Sec. 6. In this sense, the role of proton decay in probing into physics at the most fundamental level is unique . In view of how valuable such a probe would be and the fact that the predicted upper limit on the proton lifetime is only a factor of three to six higher than the empirical lower limit, the argument in favor of building an improved detector seems compelling.
To conclude, the discovery of proton decay would undoubtedly constitute a landmark in the history of physics. It would provide the last, missing piece of gauge unification and would shed light on how such a unification may be extended to include gravity.
Acknowledgements : I would like to thank Kaladi S. Babu and Frank Wilczek for a most enjoyable collaboration, and Joseph Sucher for valuable discussions. I would also like to thank the organizers of the NNN99 workshop, especially Chang Kee Jung, Millind Diwan and Hank Sobel, for arranging a stimulating meeting and also for the kind hospitality. The research presented here is supported in part by DOE grant no. DE-FG02-96ER-41015.
## Appendix <br>A Natural Doublet-Triplet Splitting Mechanism in SO(10)
In supersymmetric SO(10), a natural doubletโtriplet splitting can be achieved by coupling the adjoint Higgs $`\mathrm{๐๐}_๐`$ to a $`\mathrm{๐๐}_๐`$ and a $`\mathrm{๐๐}_๐^{}`$, with $`\mathrm{๐๐}_๐`$ acquiring a unificationโscale VEV in the $`B`$-$`L`$ direction : $`\mathrm{๐๐}_๐=(a,a,a,0,0)\times \tau _2`$ with $`aM_U`$. As discussed in Section 2, to generate CKM mixing for fermions we require $`(\mathrm{๐๐}_๐)_d`$ to acquire a VEV of the electroweak scale. To ensure accurate gauge coupling unification, the effective low energy theory should not contain split multiplets beyond those of MSSM. Thus the MSSM Higgs doublets must be linear combinations of the SU(2$`)_L`$ doublets in $`\mathrm{๐๐}_๐`$ and $`\mathrm{๐๐}_๐`$. A simple set of superpotential terms that ensures this and incorporates doublet-triplet splitting is :
$`W_H=\lambda \mathbf{\hspace{0.17em}10}_๐\mathbf{\hspace{0.17em}45}_๐\mathbf{\hspace{0.17em}10}_๐^{}+M_{10}\mathbf{\hspace{0.17em}10}_{๐}^{}{}_{}{}^{2}+\lambda ^{}\overline{\mathrm{๐๐}}_H\overline{\mathrm{๐๐}}_H\mathbf{\hspace{0.17em}10}_H+M_{16}\mathbf{\hspace{0.17em}16}_H\overline{\mathrm{๐๐}}_H.`$ (A1)
A complete superpotential for $`\mathrm{๐๐}_๐`$, $`\mathrm{๐๐}_๐`$, $`\overline{\mathrm{๐๐}}_๐`$, $`\mathrm{๐๐}_H`$, $`\mathrm{๐๐}_H^{}`$ and possibly other fields, which ensure that $`\mathrm{๐๐}_๐`$, $`\mathrm{๐๐}_๐`$ and $`\overline{\mathrm{๐๐}}_๐`$ acquire unification scale VEVs with $`\mathrm{๐๐}_๐`$ being along the $`(B`$-$`L)`$ direction, that exactly two Higgs doublets $`(H_u,H_d)`$ remain light, with $`H_d`$ being a linear combination of $`(\mathrm{๐๐}_๐)_d`$ and $`(\mathrm{๐๐}_๐)_d`$, and that there are no unwanted pseudoGoldstone bosons, can be constructed. With $`\mathrm{๐๐}_๐`$ in the $`B`$-$`L`$ direction, it does not contribute to the Higgs doublet mass matrix, so one pair of Higgs doublet remains light, while all triplets acquire unification scale masses. The light MSSM Higgs doublets are
$`H_u=\mathbf{\hspace{0.17em}10}_u,H_d=\mathrm{cos}\gamma \mathbf{\hspace{0.17em}10}_d+\mathrm{sin}\gamma \mathbf{\hspace{0.17em}16}_d,`$ (A2)
with $`\mathrm{tan}\gamma \lambda ^{}\overline{\mathrm{๐๐}}_๐/M_{16}`$. Consequently, $`\mathrm{๐๐}_d=(\mathrm{cos}\gamma )v_d`$, $`\mathrm{๐๐}_d=(\mathrm{sin}\gamma )v_d`$, with $`H_d=v_d`$ and $`\mathrm{๐๐}_d`$ and $`\mathrm{๐๐}_d`$ denoting the electroweak VEVs of those multiplets. Note that $`H_u`$ is purely in $`\mathrm{๐๐}_๐`$ and that $`\mathrm{๐๐}_d^2+\mathrm{๐๐}_d^2=v_d^2`$. This mechanism of doublet-triplet (DT) splitting is rather unique for the minimal Higgs systems in that it meets the requirements of both D-T splitting and CKM-mixing. In turn, it has three special consequences:
(i) It modifies the familiar SO(10)-relation $`\mathrm{tan}\beta v_u/v_d=m_t/m_b60`$ to:
$`\mathrm{tan}\beta /\mathrm{cos}\gamma m_t/m_b60`$ (A3)
As a result, even low to moderate values of $`\mathrm{tan}\beta 3`$ to 10 (say) are perfectly allowed in SO(10) (corresponding to $`\mathrm{cos}\gamma 1/20`$ to $`1/6`$).
(ii) The most important consequence of the DT-splitting mechanism outlined above is this: In contrast to SU(5), for which the strengths of the standard d=5 operators are proportional to $`(M_{H_c})^1`$ (where $`M_{H_C}few\times 10^{16}`$ GeV (see Eq. (70)), for the SO(10)-model, they become proportional to $`M_{eff}^1`$, where $`M_{eff}=(\lambda a)^2/M_{10^{}}M_U^2/M_{10^{}}`$. $`M_{10^{}}`$ can be naturally smaller (due to flavor symmetries) than $`M_U`$ and thus $`M_{eff}`$ correspondingly larger than $`M_U`$ by one to two orders of magnitude (see Ref. ). Now the proton decay amplitudes for SO(10) in fact possess an intrinsic enhancement compared to those for SU(5), owing primarily due to differences in their Yukawa couplings for the up sector (see Appendix C in Ref. ). As a result, these larger values of $`M_{eff}10^{18}`$ GeV are in fact needed for the SO(10)-model to be compatible with the observed limit on the proton lifetime. At the same time, being bounded above (see below), they allow optimism as regards future observation of proton decay.
(iii) $`M_{eff}`$ gets bounded above by considerations of coupling unification and GUT-scale threshold effects. Owing to mixing between $`10_d`$ and $`16_d`$ (see Eq. (A2)), the threshold correction to $`\alpha _3(m_z)`$ due to doublet-triplet splitting becomes proportional to $`\mathrm{ln}(M_{eff}\mathrm{cos}\gamma /M_U)`$. Inclusion of this correction and those due to splittings within the gauge and the Higgs multiplets (i.e. $`45_H`$, $`16_H`$, and $`\overline{16}_H`$)<sup>14</sup><sup>14</sup>14The correction to $`\alpha _3(m_z)`$ due to Planck scale physics through the effective operator $`F_{\mu \nu }F^{\nu \mu }45_H/M`$ vanishes due to antisymmetry in the SO(10)-contraction., together with the observed degree of coupling unification allows us to obtain a conservative upper limit on $`M_{eff}`$, given by :
$`M_{eff}\mathrm{\hspace{0.17em}3}\times 10^{18}\text{GeV}.`$ (A4)
This in turn helps provide an upper limit on the expected proton decay lifetime (see text).
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# Electromagnetic moments of relativistic higher spin baryons
## Abstract
We point out a source of ambiguities in the measurements of the electromagnetic moments of spin-3/2 baryons which rely on relativistic Lagrangian models. An anambiguous relation between the parameters of the spin-3/2 electromagnetic Lagrangian and the electromagnetic moments of the spin-3/2 particle exits in general only for on-shell situation, while the measurements are done on โvirtualโ baryons.
Accurate measurement of electromagnetic (EM) moments of $`N^{}`$ resonances could provide an important testing ground for many quark-models and lattice QCD predictions. Because of the short lifetime of the $`N^{}`$ resonances such measurements present experimental as well as theoretical challenge. At present one can hope only for indirect measurements. For example, the dipole magnetic moment of the $`\mathrm{\Delta }^{++,0}`$ isobars is extracted from the observables of radiative pion-nucleon scattering, $`\pi ^\pm +p\pi ^\pm +p+\gamma `$ (see, e.g., ), while the magnetic moment of the $`\mathrm{\Delta }^+`$ can possibly soon be measured at MAMI in the radiative pion photoproduction , $`\gamma +p\pi ^0+p+\gamma `$.
In these reactions the $`\gamma \mathrm{\Delta }\mathrm{\Delta }`$ current containing the EM moments is involved only through more complicated mechanisms, see Fig. 2. Moreover, other mechanisms (such as in Fig. 2) may interfere and must be estimated reliably. The necessary theoretical input is usually provided by some effective Lagrangian model based on hadronic degrees of freedom , where all the relevant mechanisms can be computed, while the parameters other than the EM moments can be fixed from different sources (e.g., by using the same effective Lagrangian for the description of other processes).
Here we shall assume a general form of the effective hadronic Lagrangian restricted only by a few fundamental principles such as Lorentz and gauge invariance, and possibly the conditions due to the approximate chiral symmetry.
We will consider the charge and dipole magnetic moment interactions of the spin-3/2 baryon and show how the parameters of the Lagrangian can unambiguously be related to the magnetic moment of the โrealโ spin-3/2 particle. For virtual particle this relation is in general ambiguous. This problem is also related to the so-called โoff-shell ambiguitiesโ of the covariant spin-3/2 description which generally does not exclude off-shell the unphysical spin-1/2 contributions. (Recently, Machavariani et al. have claimed to get rid of the โoff-shell ambiguitiesโ formulation by writing the spin-3/2 propagator in terms of on-shell Rarita-Schwinger vector-spinors. However, in doing so one looses relativistic covariance, see Ref. for details.) Furthermore, by considering the matrix elements of processes such as in Fig. 2, we can argue that certain consistency conditions on the form of the $`\pi N\mathrm{\Delta }`$ and $`\gamma N\mathrm{\Delta }`$ couplings may allow us to deal with the โoff-shell ambiguitiesโ.
To describe the spin-3/2 fields of (decuplet) baryons it is natural to use the covariant Rarita-Schwinger (RS) formalism where the field is a Lorentz vector-spinor $`\psi ^\mu (x)`$ with the following free Lagrangian<sup>1</sup><sup>1</sup>1Other forms of free spin-3/2 Lagrangian are frequently used in the literature, but they relate to Eq. (1) by a field redefinition $`\psi _\mu (g_{\mu \nu }+b\gamma _\mu \gamma _\nu )\psi ^\nu `$, with $`b\frac{1}{4}`$. Present discussion is fixed to representation (1), however is unchanged for other choices albeit the field redefinition is done in the whole Lagrangian. :
$$_{\mathrm{free}}=\overline{\psi }_\mu \gamma ^{\mu \nu \alpha }_\alpha \psi _\nu m\overline{\psi }_\mu \gamma ^{\mu \nu }\psi _\nu $$
(1)
where $`m`$ is the mass; $`\gamma ^{\mu \nu }`$ and $`\gamma ^{\mu \nu \alpha }`$ is the totally antisymmetrized product of, respectively, two and three gamma-matrices. This Lagrangian leads to the well-known RS propagator:
$$S^{\mu \nu }(p)=\frac{p/+m}{p^2m^2+i\epsilon }\left[g^{\mu \nu }+\frac{1}{3}\gamma ^\mu \gamma ^\nu +\frac{1}{3m}(\gamma ^\mu p^\nu \gamma ^\nu p^\mu )+\frac{2}{3m^2}p^\mu p^\nu \right],$$
(2)
and the following free-field equations:
$$(i/m)\psi _\mu =0,\gamma ^\mu \psi _\mu =0,^\mu \psi _\mu =0.$$
(3)
Consider now the EM properties of this field. The electric charge is included via minimal substitution, $`_\alpha _\alpha +ieA_\alpha `$, into Eq. (1). To explore the dipole magnetic moment we should include all the nonminimal, linear in $`F_{\mu \nu }=_{[\mu }A_{\nu ]}`$ terms. We then have:
$`_{\mathrm{int}}`$ $`=`$ $`e\{i\overline{\psi }_\mu \gamma ^{\mu \nu \alpha }\psi _\nu A_\alpha +l_1\overline{\psi }{}_{\mu }{}^{}F_{}^{\mu \nu }\psi _\nu +l_2\overline{\psi }{}_{\mu }{}^{}\gamma .F.\gamma \psi ^\mu +l_3F^{\mu \nu }[\overline{\psi }{}_{\mu }{}^{}\gamma _{\nu }^{}\gamma .\psi +\overline{\psi }.\gamma \gamma _\mu \psi _\nu ]`$ (4)
$`+l_4\overline{\psi }.\gamma (\gamma .F.\gamma )\gamma .\psi +il_5F^{\mu \nu }[\overline{\psi }{}_{\mu }{}^{}\gamma _{\nu }^{}\gamma .\psi \overline{\psi }.\gamma \gamma _\mu \psi _\nu ]\}\overline{\psi }_\mu \Lambda ^{\mu \nu }\psi _\nu .`$
It is further argued in that $`\Lambda ^{\mu \nu }`$ must be antisymmetric in order to maintain the correct number of spin degrees-of-freedom<sup>2</sup><sup>2</sup>2This condition does not remove the consistency problem of the minimal coupling , but only insures the proper degrees-of-freedom count.: $`2s+1=4`$. Hence, one finds $`l_2=l_4`$, $`l_3=2l_4`$, and the couplings reduce to
$$\Lambda ^{\mu \nu }=ie\gamma ^{\mu \nu \alpha }A_\alpha +em^1\left[L_1F^{\mu \nu }+iL_2\gamma ^5\stackrel{~}{F}^{\mu \nu }+iL_3(F^{\mu \alpha }\gamma _\alpha \gamma ^\nu \gamma ^\mu \gamma _\alpha F^{\alpha \nu })\right],$$
(5)
where $`\stackrel{~}{F}^{\mu \nu }=\epsilon ^{\mu \nu \varrho \sigma }_\varrho A_\sigma `$; $`L_i`$ are dimensionless parameters: $`L_1=m(l_1+2l_2)`$, $`L_2=2ml_2`$, $`L_3=ml_5`$. To relate these parameters to the physical magnetic moment we match this field theory to the soft-photon theorems.
Consider first the on-shell electromagnetic vertex \[obtained from (4) with (5)\]:
$`\epsilon _\lambda \Gamma ^\lambda (q)`$ $`=`$ $`em^1\epsilon _\lambda \overline{u}_\mu (p^{})\{m\gamma ^{\mu \nu \lambda }+L_1(q^\mu g^{\lambda \nu }q^\nu g^{\lambda \mu })+iL_2\gamma ^5\epsilon ^{\mu \nu \sigma \lambda }q_\sigma `$ (6)
$`+iL_3[(q^\mu \gamma ^\lambda q/g^{\lambda \mu })\gamma ^\nu +\gamma ^\mu (q^\nu \gamma ^\lambda q/g^{\lambda \nu })]\}u_\nu (p),`$
where $`\epsilon `$ and $`q=p^{}p`$ is the polarization and the 4-momentum of the photon; $`u`$ are the on-shell RS vector-spinors, solutions of Eq. (3) in momentum space.
Using the properties of the on-shell vector-spinor one can see that (i) $`L_3`$ term drops out, (ii) $`L_2`$ term contributes only to higher order in $`q`$ (basically because $`\gamma ^5`$ involves the โsmall componentsโ of the spinors). Also one can (iii) use various Gordon identities and (iv) verify the following nontrivial identity:
$$2\overline{u}_\varrho ^{}\gamma _{\mu \nu }u_\sigma =\overline{u}_\varrho ^{}(\gamma _{\mu \sigma }u_\nu +\gamma _{\sigma \nu }u_\mu )+(\overline{u}_\mu ^{}\gamma _{\varrho \nu }+\overline{u}_\nu ^{}\gamma _{\mu \varrho })u_\sigma +๐ช(q),$$
(7)
where $`u^{}u(p^{})`$, $`uu(p)`$, $`q=p^{}p`$. Using all that, the vertex to the 1st order in $`q`$ reads
$$\epsilon _\lambda \Gamma ^\lambda =\frac{e}{2m}\epsilon _\lambda \overline{u}_\mu ^{}\left[(p^{}+p)^\lambda g^{\mu \nu }gM_{3/2}^{\lambda \sigma ,\mu \nu }q_\sigma \right]u_\nu ,$$
(8)
with $`M_{3/2}`$ being the spin-3/2 Lorentz generator and
$$g=\frac{2}{3}(12L_1).$$
(9)
On the other hand, according to the soft-photon theorem $`g=(\mu /s)(e/2m)^1`$ is the gyromagnetic ratio; $`\mu `$ is the magnetic moment.<sup>3</sup><sup>3</sup>3Let as remark on the definition of the โanomalous magnetic momentโ, $`\mu _A`$. It is usually defined as the deviation from the value implied by the minimal coupling: $`\mu _A=\mu e/2m`$. However, taking the point of view that $`g=2`$ is the universal value for all elementary (pointlike) particles with spin , it is more correct to define it as the deviation from the fundamental value, i.e., as follows: $`\mu _A=\mu sem^1`$.
The same relation is obtained by analyzing the Compton scattering on spin-3/2 particle. For the forward Compton amplitude we obtain (to the 1st order in photon frequency $`\omega `$):
$$T_{fi}=\frac{e^2}{m}\epsilon ^{}.\epsilon \overline{u}^{}.u+\frac{ie^2\omega }{4m^2}(g2)^2\overline{u}_\mu ^{}M_{3/2}^{\lambda \sigma ,\mu \nu }u_\nu \epsilon _\lambda ^{}\epsilon _\sigma ,$$
(10)
again with the gyromagnetic ratio given by Eq. (9). The use of on-shell conditions such as Eq. (7) is made in obtaining this result.
Obviously, a potential problem arises when one tries to identify the magnetic moment of the virtual particle. One cannot use any of the statements (i)-(iv) to conclude an unambiguous relation between the parameters of the Lagrangian and the physical magnetic moment. In particular, because of the coupling to the spin-1/2 sector of the RS propagator, $`L_2`$ and $`L_3`$ terms may give contributions of the same order in photon energy as the minimal and $`L_1`$ terms. This problem does not arise if the spin-3/2 propagator could be replaced by a positive-energy spin-3/2 projection operator. However, such replacement can be done only at the expense of loosing the Lorentz and gauge invariance. (Gauge invariance is lost because the EM vertex satisfies the Ward-Takahashi identity for the full RS propagator and not for any truncated form).
A way to deal with the problem is to consider the complete matrix elements (for graphs like in Fig. 2) and use the specific โconsistentโ $`\pi N\mathrm{\Delta }`$ and $`\gamma N\mathrm{\Delta }`$ interactions . The latter interactions have the property of decoupling the spin-1/2 sector, and therefore the major source of the off-shell ambiguity should be eliminated. A study in this direction and development of models for the radiative $`\pi N`$ scattering and pion photoproduction is in progress .
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# Plasma Energy Loss into Kaluza-Klein Modes
## 1. Introduction
Recently, Barger et al. addressed astophysical constraints on extra dimensions by computing energy loss rates from the sun, red giant stars and (Type II) supernovae due to possible excitation of graviton modes, $`๐ข`$, in the case that the extra dimensions are compactified . The processes $`\gamma \gamma ๐ข,e^+e^{}๐ข,\gamma e๐ข,eN๐ขeN`$ (in the static nucleon approximation), and $`NN๐ขNN`$ where considered. They worked in the zero density approximation, varying only the temperature. Their calculation neglected plasma effects and they anticipated that this neglect should not be important because of the high power dependence on $`M_S`$, the inverse of the compactification dimension.
The purpose of this paper is to address the extent to which the process $`\gamma \gamma ๐ข`$ (and the crossed process $`\gamma \gamma ๐ข`$) is affected by a non-zero charged particle density and the presence of both longitudinal and transverse plasmons. Our aims are, in brief, to find numerical values for the energy loss rate per cm<sup>3</sup> for densities from $`1`$ to $`10^{15}`$ gm/cm<sup>3</sup> and temperatures from $`1`$ to $`10^9`$ eV, to confirm the Barger et al. expectation of density insensitivity, to show the extent to which the systems they considered are close to the border in density at which their expectation fails, to determine the relative contributions of longitudinal and transverse plasmon processes, to address the size of the contribution of the crossed process, and, importantly, to note the ambiguities in the form of a covariant interaction between plasmons and Kaluza-Klein modes.
It is clear that the expectation of Barger et al. must fail at sufficiently high density, for fixed temperature, since the energy loss rate goes as an integral over the Bose-Einstein distribution, $`(e^{\omega _i/kT}1)^1`$, in which the frequencies are given by a dispersion relation with an effective photon mass that grows with density. It is essentially a numerical question as to the point at which suppression sets in and the rate (in density) at which it proceeds. Intuition is hampered by the fact that the natural parameter is the electron chemical potential which is not simply related to the particle density. We give the numerical results over the temperature and density ranges cited above.
An important problem in carrying out this work is the question of the appropriate Lagrangian. The free space coupling between the electromagnetic field and gravitons can be found in textbooks, for example , and has been generalized to the case of higher dimensional Kaluza-Klein excitations . However, we have not found a parallel literature for the case in which the free space photon is replaced by a plasmon satisfying a non-trivial dispersion relation. This difficulty is addressed in Section 2. We adopt a diagrammatic approach and also make an approximation that we test numerically. Also given in Section 2 is the formalism used for the numerical calculations of Section 3. We conclude in Section 4 with a brief summary.
## 2. Formalism
In a medium with nonzero temperature and density, radiation satisfies the dispersion relation
$$\omega ^2|\stackrel{}{k}|^2=\mathrm{\Pi }_a(\omega ,|\stackrel{}{k}|).$$
(1)
$`\mathrm{\Pi }_a(\omega ,|\stackrel{}{k}|)`$ is the transverse or longitudinal component of the polarization tensor,
$$\mathrm{\Pi }_a(\omega ,|\stackrel{}{k}|)=\epsilon _a^\mu \mathrm{\Pi }_{\mu \nu }\epsilon _a^\nu ,$$
(2)
where $`\epsilon _a^\mu `$ ($`a=L,T`$) are the polarization vectors and the polarization tensor $`\mathrm{\Pi }_{\mu \nu }`$ is the photon self energy in the medium. The contribution to this self energy from fermions in the medium is calculated in the medium rest frame by adding a term
$$\frac{p/+m}{e^{(E\pm \mu )/T}+1}2\pi i\delta (p^2m^2)$$
(3)
to the usual (vacuum) progagator. Here, $`T`$ is the temperature, $`E=p^0`$ and $`\mu `$ is the electron chemical potential, which is related to the electron number density $`n_e`$ by
$$2\frac{d^{\mathrm{\hspace{0.17em}3}}p}{(2\pi )^3}\left[\frac{1}{e^{(E\mu )/T}+1}\frac{1}{e^{(E+\mu )/T}+1}\right]=n_e$$
(4)
To lowest order in the fine structure constant $`\alpha `$, the polarization tensor is given by
$`\mathrm{\Pi }^{\mu \nu }`$ $`=`$ $`16\pi \alpha {\displaystyle \frac{d^{\mathrm{\hspace{0.17em}3}}p}{(2\pi )^32E}\left[\frac{1}{e^{(E\mu )/T}+1}+\frac{1}{e^{(E+\mu )/T}+1}\right]}`$ (5)
$`\times {\displaystyle \frac{(pk)^2g^{\mu \nu }+k^2p^\mu p^\nu pk(k^\mu p^\nu +k^\nu p^\mu )}{(pk)^2}}.`$
It turns out that $`\mathrm{\Pi }_L`$ and $`\mathrm{\Pi }_T`$ can be approximated to within 1% for all temperatures and densities by
$`\mathrm{\Pi }_L`$ $`=`$ $`\omega _P^2\left[1G(v_{}^2|\stackrel{}{k}|^2/\omega ^2)\right]+v_{}^2|\stackrel{}{k}|^2|\stackrel{}{k}|^2,`$ (6)
$`\mathrm{\Pi }_T`$ $`=`$ $`\omega _P^2\left[1+{\displaystyle \frac{1}{2}}G(v_{}^2|\stackrel{}{k}|^2/\omega ^2)\right],`$ (7)
where $`v_{}`$ is an average value of $`v=|\stackrel{}{p}|/E`$ for the electron (the only fermion which contributes for stellar temperatures and densities). Explicitly,
$$v_{}=\omega _1/\omega _P,$$
(8)
with $`\omega _1`$ given by
$$\omega _1^2=\frac{4\alpha }{\pi }_0^{\mathrm{}}d|\stackrel{}{p}||\stackrel{}{p}|\left(\frac{5}{3}v^3v^5\right)f_E,$$
(9)
where $`f_E`$ is the sum of the electron and positron distributions (the square bracket in Eq. (5) above). The plasma frequency $`\omega _P`$ is given by
$$\omega _P^2=\frac{4\alpha }{\pi }_0^{\mathrm{}}d|\stackrel{}{p}||\stackrel{}{p}|\left(v\frac{1}{3}v^3\right)f_E,$$
(10)
and the function $`G(x)`$ is
$$G(x)=\frac{3}{x}\left[1\frac{2x}{3}\frac{1x}{2\sqrt{x}}\mathrm{log}\left(\frac{1+\sqrt{x}}{1\sqrt{x}}\right)\right].$$
(11)
It will be important below to note that, for transverse photons, $`k^2=\omega ^2|\stackrel{}{k}|^2`$ is $`\omega _P^2`$ for $`|\stackrel{}{k}|=0`$ and increases as $`|\stackrel{}{k}|`$ increases, while, for longitudinal photons, $`k^2=\omega _P^2`$ at $`|\stackrel{}{k}|=0`$ and decreases as $`|\stackrel{}{k}|`$ increases. Integration over $`|\stackrel{}{k}|`$ for longitudinal photons must be cut off at the point where $`k^2`$ becomes negative,
$$|\stackrel{}{k}|_{\mathrm{max}}^2=\frac{3\omega _P^2}{v_{}^2}\left[\frac{1}{2v_{}^2}\mathrm{log}\left(\frac{1+v_{}}{1v_{}}\right)\right].$$
(12)
We include in our numerical evaluations the renormalization constants
$$Z_a^1=1\frac{\mathrm{\Pi }_a}{\omega ^2},$$
(13)
although this is inconsistent with calculating only to the lowest order in $`\alpha `$. It is a check on our results that they do not change significantly when the $`Z_a`$ are set to unity. The $`Z_a`$ are given by
$`Z_T`$ $`=`$ $`{\displaystyle \frac{2\omega ^2(\omega ^2v_{}^2|\stackrel{}{k}|^2)}{\omega ^2[2\omega _P^22(\omega ^2|\stackrel{}{k}|^2)]+(\omega ^2+|\stackrel{}{k}|^2)(\omega ^2v_{}^2|\stackrel{}{k}|^2)}}`$ (14)
$`Z_L`$ $`=`$ $`{\displaystyle \frac{2\omega ^2(\omega ^2v_{}^2|\stackrel{}{k}|^2)}{3\omega _P^2\omega ^2+v_{}^2|\stackrel{}{k}|^2}}{\displaystyle \frac{\omega ^2}{\omega ^2|\stackrel{}{k}|^2}}.`$ (15)
The rate of graviton emission can be calculated using the Lagrangian for the coupling of Kaluza-Klein field $`๐ข_\stackrel{}{n}^{\mu \nu }`$, corresponding to the mass excitation $`m_\stackrel{}{n}^2=(2\pi )^2\stackrel{}{n}^2/R^2`$, to the photon energy-momentum tensor $`T_{\mu \nu }`$. Neglecting gauge terms, this coupling is
$``$ $`=`$ $`{\displaystyle \frac{\kappa }{2}}๐ข_\stackrel{}{n}^{\mu \nu }T_{\mu \nu }`$ (16)
$`=`$ $`{\displaystyle \frac{\kappa }{2}}\left(๐ข_\stackrel{}{n}^{\mu \nu }F_\mu ^\lambda F_{\nu \lambda }{\displaystyle \frac{1}{4}}๐ข_{\stackrel{}{n},\mu }^\mu F^{\lambda \rho }F_{\lambda \rho }\right),`$
where $`F_{\mu \nu }`$ is the electromagnetic field tensor. We consider only the coupling of the spin-2 component of the Kaluza-Klein field; the spin-0 component does not couple to photons.
The matrix element for $`\gamma (k_1)\gamma (k_2)๐ข`$ obtained from Eq. (16) is
$$=\frac{\kappa }{2}\epsilon ^\lambda (k_1)\epsilon ^\rho (k_2)\epsilon ^{\mu \nu }(k_1+k_2)T_{\mu \nu ,\lambda \rho }^{(0)},$$
(17)
where
$$T_{\mu \nu ,\lambda \rho }^{(0)}=k_1k_2C_{\mu \nu ,\lambda \rho }+D_{\mu \nu ,\lambda \rho }(k_1,k_2),$$
(18)
with
$`C_{\mu \nu ,\lambda \rho }`$ $`=`$ $`\eta _{\mu \lambda }\eta _{\nu \rho }+\eta _{\mu \rho }\eta _{\nu \lambda }\eta _{\mu \nu }\eta _{\lambda \rho }`$ (19)
$`D_{\mu \nu ,\lambda \rho }(k_1,k_2)`$ $`=`$ $`\eta _{\mu \nu }k_{1\rho }k_{2\lambda }[\eta _{\mu \rho }k_{1\nu }k_{2\lambda }+\eta _{\mu \lambda }k_{1\lambda }k_{2\nu }+(\mu \nu )].`$ (20)
The sum over polarizations of the Kaluza-Klein state is
$$\underset{s=1}{\overset{5}{}}\epsilon _{\mu \nu }^s(k)\epsilon _{\lambda \rho }^s(k)=\frac{1}{2}B_{\mu \nu ,\lambda \rho }(k),$$
(21)
with
$`B_{\mu \nu ,\lambda \rho }(k)`$ $`=`$ $`E_{\mu \lambda }E_{\nu \rho }+E_{\mu \rho }E_{\nu \lambda }{\displaystyle \frac{2}{3}}E_{\mu \nu }E_{\lambda \rho }`$ (22)
$`E_{\mu \nu }`$ $`=`$ $`\eta _{\mu \nu }{\displaystyle \frac{k_\mu k_\nu }{m_\stackrel{}{n}^2}}.`$ (23)
The coupling Eq. (18) is gauge invariant even if $`k_1^2`$ and $`k_2^2`$ are not zero, e.g. $`k_1^\mu T_{\mu \nu ,\lambda \rho }^{(0)}=0`$. However, it is not conserved, $`(k_1+k_2)^\mu T_{\mu \nu ,\lambda \rho }^{(0)}0`$, if $`k_1^2`$ and/or $`k_2^2`$ differs from zero. We cannot write a conserved coupling by using the energy-momentum tensor for a massive vector field because $`k_1^2`$ is not necessarily equal to $`k_2^2`$. This means that if we square $``$ of Eq. (17), and use Eq. (21), we get extra terms of the form $`k_1^2/m_\stackrel{}{n}^2`$ or $`k_2^2/m_\stackrel{}{n}^2`$ from the second term in Eq. (23).
To have a conserved amplitude with $`k_1^2k_2^20`$, we must include all the diagrams of Fig. (1). The Feynman rules for the
fermion-fermion-$`๐ข`$ coupling and the fermion-fermion-photon-$`๐ข`$ are given in Refs. , and the loops are calculated by using Eq. (3) for one of the legs. We have shown the the sum of these diagrams is gauge invariant and conserved for arbitrary $`k_1^2`$ and $`k_2^2`$ at finite temperature and density. However, this was done without actually evaluating the diagrams. In particular, diagram (d) is very tedious and we have not computed it. Instead, we have used only diagram (a) (Eq. (17) above) but have evaluated every energy loss twice - once including the $`k_1^2/m_\stackrel{}{n}^2`$ and $`k_2^2/m_\stackrel{}{n}^2`$ terms and once omitting them. In every case, the results were almost identical. While this proves nothing, it does seem to indicate that performing the complete one-loop calculation would not give a substantially different answer.
The reaction rate must be summed over the Kaluza-Klein states, which is done by integrating over
$$dm_n^2\frac{4\pi m_n^{(n2)}}{\kappa ^2M_S^{n+2}},$$
(24)
where $`n`$ is the number of extra dimensions. $`M_S`$ is the string scale which is related to the compactification scale $`R`$ and Newtonโs constant $`G_N`$. Specifically, we use
$$M_S^{n+2}R^n=\frac{(4\pi )^{n/2}\mathrm{\Gamma }(n/2)}{4G_N}.$$
(25)
Our definition of $`M_S^{n+2}`$ differs from that of by a factor of $`2`$, i.e. their $`M_S`$ is larger by a factor $`2^{1/(n+2)}`$. As a consequence, values of the energy loss per unit volume obtained from our tables must be multiplied by 2 when comparing with Barger et al..
For 2 particles $``$ 1 particle reactions there remains a delta function from phase space which identifies $`m_n^2`$ with the center of mass squared energy $`s`$. Thus, the integral over $`m_n^2`$, Eq. (24), replaces $`m_n^2`$ by $`s`$ and our results depend on $`n`$ through the factor $`s^{(n2)/2}/M_S^{(n+2)}`$.
The rate of energy loss per unit volume is given by the standard expression
$$๐ฌ_{a,b}=A\frac{d^3k_1}{(2\pi )^3}\frac{1}{e^{\omega _1/T}1}\frac{d^3k_2}{(2\pi )^3}\frac{1}{e^{\omega _2/T}1}(\omega _1+\omega _2)Z_a(|\stackrel{}{k}_1|)Z_b(|\stackrel{}{k}_2|)v\sigma ,$$
(26)
where $`v\sigma `$ denotes the cross section times the relative velocity . The initial photons have $`k_i^2=\omega _i^2|\stackrel{}{k}_i|^2`$, $`i=1,2`$ and can be transverse, $`a=b=T`$, longitudinal, $`a=b=L`$ or mixed, e.g. $`a=L,b=T`$. The factor $`A`$ gives the number of spin states: $`A=4,2,1`$ for $`TT`$, $`TL`$ or $`LL`$. For longitudinal photons, the $`|\stackrel{}{k}|`$ integrals are cut off at $`|\stackrel{}{k}|_{\mathrm{max}}`$ given by Eq. (12). The corresponding expression for the energy loss in the decay $`TL๐ข`$ is
$`๐ฌ_{TL}`$ $`=`$ $`2{\displaystyle \frac{d^3k_T}{(2\pi )^32\omega _T}\frac{1}{e^{\omega _T/T}1}\frac{d^3k_L}{(2\pi )^32\omega _L}\frac{(\omega _T\omega _L)}{1e^{\omega _L/T}}Z_T(|\stackrel{}{k}_T|)Z_L(|\stackrel{}{k}_L|)}`$ (27)
$`\times {\displaystyle \frac{(2\pi )^4\left((k_Tk_L)^2\right)^{(n2)/2}}{\kappa ^2M_S^{n+2}}}||^2,`$
where $``$ is given by Eq. (17) with $`k_2k_L`$.
## 3. Calculations
The first step in the calculations is to obtain $`\mu (T,\rho )`$ from Eq. (4). In doing this we assume that the electron number density, $`n_e`$ is related to the mass density $`\rho `$ by $`n_e=\rho /m_p`$ where $`m_p`$ is the proton mass. This is useful for comparison purposes and is a reasonable order of magnitude approximation but needs correction (by less than an order of magnitude) for a supernova or a neutron star. The results of the calculation of $`\mu `$ are given in Tables 10 and 11 for the matrix of $`\rho `$ and $`T`$ values: $`\rho =1.0,\mathrm{\hspace{0.17em}10.0},\mathrm{}\mathrm{\hspace{0.17em}10}^{15}`$ gm/cm<sup>3</sup> and $`T=1,\mathrm{\hspace{0.17em}10}^2,\mathrm{}\mathrm{\hspace{0.17em}10}^9\mathrm{eV}`$. Two tables are given ($`\mu (T,\rho )`$ and $`\stackrel{~}{\mu }(T,\rho )=\mu (T,\rho )m_e`$) in order to make clear both the deviation of $`\mu `$ from $`m_e`$ (taken as 0.51 MeV) at low temperature and its deviation from zero at high. Note the rapid variation of $`\mu `$ (for the lowest densities) at the temperature $`T`$ around $`0.1\mathrm{MeV}`$ where pair production first begins to be copious, and the slower but similar variation at higher densities $`\rho `$. The variation is slower for higher densities because the electron-positron density difference needs to have a large value. These variations are illustrated in Fig. 2, where, for display purposes, the lowest value of $`\mu `$ shown, $`10^2`$, is an upper limit on the exact numbers in Table 10.
In calculating $`\mu `$ we used an iteration procedure and required that the output density value equal the input to better than a percent.
We now pass on to the results of calculating $`๐ฌ=d^4E/dtdV`$ for the case of $`n=2`$ extra dimensions. It was possible to evaluate the integral in Eq. (26) over the cosine of the angle between the two plasmons analytically so that, for all the processes under consideration, only two integrals remain in finding the energy loss rate - the integrals over the two plasmon momenta.
It should be noted that there are only four processes to consider: (1) $`T+T๐ข`$, (2) $`T+L๐ข`$, (3) $`TL+๐ข`$, and (4) $`L+L๐ข`$. This is because, as the plasmon momentum $`|\stackrel{}{k}|`$ increases, the effective mass of a transverse plasmon increases while the effective mass of a longitudinal plasmon decreases . Thus the missing processes, $`LT+๐ข`$ and $`LL+๐ข`$ are forbidden by energy-momentum conservation. The assertion is clear for the first process since $`m_T>m_L`$. For the second, we note that, in the rest frame of the decaying longitudinal plasmon, $`m_L=\omega _P`$ and conservation of energy and momentum implies that the graviton mass, $`m_๐ข`$, satisfies
$$m_๐ข^2=\left(\omega _P\omega _L(|\stackrel{}{k}|)\right)^2|\stackrel{}{k}|^2,$$
(28)
where $`\omega _L(|\stackrel{}{k}|)`$ is the energy of the final plasmon. Using the dispersion relation for longitudinal plasmons, it can be shown that the right side of Eq. (28) is less than zero for $`|\stackrel{}{k}|>0`$.
The results of the calculations are given in Tables 1-5. These have the energy loss rates for the four processes, and for the sum, for a matrix of density and temperature values - $`1.0`$ to $`10^{15}`$ gm/cm<sup>3</sup> for density $`\rho `$ and $`1`$ to $`10^9`$ eV for temperature $`T`$ \- in both cases in factors of $`10`$ increments. In these tables, $`T`$ increases from left to right while $`\rho `$ increases from top to bottom. The entries are logs to the base 10 of the energy loss rate in ergs per cm<sup>3</sup>-s. Note that Barger et al. give results per unit mass, but results per unit volume are better for our purposes since they show more clearly the way in which the zero-density approximation breaks down as the density increases. We give the results for $`M_S=1`$ TeV. Two additional tables, 6 and 7, give respectively the number of the process that dominates for the reaction (zero if the rate is zero, i.e. below $`10^{320}`$) and the fraction of the total represented by the dominant contribution.
In Fig. 3 ($`T+T๐ข`$), we see at a glance the effect cited in the
Introduction: for fixed $`T`$ the energy loss rate is independent of the density until $`\rho `$ increases to a point where the effective photon mass and plasmon density are sufficiently high that the rate drops exponentially. The numerical values are given in Table 1.
The analogous plot for $`L+L๐ข`$ is shown in Fig. 4. Here, the the rate grows slightly for a fixed $`T`$ before dropping exponentially in the mass of the longitudinal plasmon with increasing $`|\stackrel{}{k}|`$. Again, numerical values are found in Table 4.
The Sun, a red giant, and a Type II supernova are, in the ($`\rho ,T`$) plane, given by Barger et al. to be at (156 gm/cm<sup>3</sup>, 1.3 keV), (10<sup>6</sup> gm/cm<sup>3</sup>, 8.6 keV), and (10<sup>15</sup> gm/cm<sup>3</sup>, 30 MeV) respectively. We see from Table 1 that the Sun is in a low density region where the zero-density approximation holds while the supernova is on the edge of a constant density region. We also see from Table 6 that the other processes contribute little in these two cases. The red giant (RG) case is more interesting. In Table 1, we see from the $`T=10^4`$ and $`T=10^3`$ columns that the RG is in the gentle fall off region for the former, but the steep fall off region for the latter. Examining the region between 1 and 10 keV more closely gives, with T varying in 1.0 keV increments, for the log of the TT energy loss rate: -16.7, -5.1, -0.91, 1.36, 2.85, 3.94, 4.78, 5.48, 6.07, 6.57. The total energy loss rate varies as: -15.96, -4.14, 0.37, 2.65, 3.77, 4.59, 5.25, 5.82, 6.31, 6.77. In short, at the RG, the TT energy loss rate, 1.4 erg/g$``$s, is down by a factor 4.5 from its zero-density approximation , but the other processes bring the total rate up to 2.6 erg/g$``$s, or within about a factor 2 of the zero density result.
Finally we turn to the dependence of the emission rates on the number n of (large) extra dimensions. We give, in Tables 8 and 9, the results for the total rates for n=3 and n=4. One sees, in the low-density limit, the $`T^{7+n}`$ behavior pointed out by Barger et al. for the $`TT๐ข`$ rate which dominates for $`n=3,4`$ in the same places as in the $`n=2`$ case.
## 4. Summary
The approximations of zero density and purely transverse ($`T+T๐ข`$) photon annihilation into gravitons of must fail, for fixed temperature, at high densities. We have computed a first estimate of finite density corrections to two-plasmon production, for both transverse and longitudinal plasmons of Kaluza-Klein excitations, as well as the decay process $`TL+๐ข`$ as a function of plasma density and temperature over a wide range of interest in both parameters. Our conclusion is that the zero-density, pure transverse approximation is satisfactory for the sun, marginal for supernovae, and fails by about a factor of 5 for red giants. It is interesting to note that, while very little of the $`\rho T`$ plane is occupied, astrophysical systems appear to be preferentially located relatively close to the boundaries (in $`\rho `$) at which the transverse photon approximation begins to fail. Our calculation is approximate in that we omit most of the diagrams of Fig. (1). However, we believe our results indicate that the full calculation of all the diagrams is unlikely to modify our conclusions.
## Acknowledgements
One of us (V. T.) wishes to thank R. Mohapatra for helpful conversations. This work was supported in part by the National Science Foundation under grant PHY-9802439 and by the Department of Energy under Contract Nos. DE-FG03-93ER40757 and DE-FG03-95ER40908.
## References
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# 1 Introduction
## 1 Introduction
Various models of random surfaces built out of triangles embedded into continuous space $`R^d`$ and surfaces built out of plaquettes embedded into Euclidean lattice $`Z^d`$ have been considered in the literature . These models are based on area action and suffer the problem of non-scaling behaviour of the string tension and the dominance of branched polymers . Several studies have analyzed the physical effects produced by rigidity of the surface introduced by adding dimension-less extrinsic curvature term to the area action
$$TS(area)+\frac{1}{\alpha }F(extrinsiccurvature),$$
where T is a string tension, $`\alpha `$ is dimension-less coupling constant and $`F(extcurv)1`$ is dimension-less functional. Comprehensive review of work in this area up to 1997 can be found in .
In authors suggested so-called gonihedric model of random surfaces which is also based on the concept of extrinsic curvature, but it differs in two essential points from the models considered in the previous studies. First it claims that only extrinsic curvature term should be considered as the fundamental action of the theory
$$S=mA(extrinsiccurvature),$$
here $`m`$ has dimension of mass and there is no area term in the action. Secondly it is required that the dependence on the extrinsic curvature should be such that the action will have dimension of length $`A(extrinsiccurvature)length.`$ This means that the action should measure the surfaces in terms of length and it should be proportional to the linear size of the surface, as it was for the path integral. The last property will guarantee that in the limit when the surface degenerates into a single world line the functional integral over surfaces will naturally transform into the Feynman path integral for point-like relativistic particle (see Figure 1,2) <sup>1</sup><sup>1</sup>1This is in contrast with the previous proposals when the extrinsic curvature term is a dimension-less functional $`F(extrinsiccurvature)1`$ and can not provide this property.
$$mA_{xy}(extrinsiccurvature)m_x^y๐l.$$
As it was demonstrated in quantum fluctuations generate the area term in the effective action
$$S_{eff}=mA(extrinsiccurvature)+T_{eff}S(area)+\mathrm{},$$
that is nonzero string tension $`T_{eff}`$. Therefore at the tree level the theory describes free quarks with string tension equal to zero, but now quantum fluctuations generate nonzero string tension and, as a result, quark confinement . The theory may consistently describe asymptotic freedom and confinement as it is expected to be the case in QCD.
Our aim is to find out what the continuum limit of this regularized theory of random surfaces is, and if indeed the continuum limit corresponds to nontrivial string theory, does it describe QCD phenomena. This is a complicated problem and one of the possible ways to handle the problem is to go further and to formulate the theory also on a lattice. An equivalent representation of the model on Euclidean lattice $`Z^d`$ has been formulated in . This representation depends on dimension $`d`$ of the embedding space. In three dimensions it is a spin model in which the interaction between spins is organized in a very specific way (see equation (8) ).In high dimensions it is a gauge spin system which contains one, two and three plaquette interaction terms.
In this article I shall consider the above model of random surfaces embedded into 3d Euclidean lattice $`Z^3`$. The reason to focus on that particular case is motivated by the fact that one can geometrically construct the corresponding transfer matrix and find its exact spectrum . The spectrum of the transfer matrix which depends only on symmetric difference of initial and final loops has been evaluated exactly in terms of correlation functions of the 2d Ising model in . This transfer matrix has the form
$$K_{\kappa =0}(Q_1,Q_2)=exp\{\beta [k(Q_1)+2l(Q_1Q_2)+k(Q_2)]\},$$
where $`Q_1`$ and $`Q_2`$ are closed polygon-loops on a two-dimensional lattice, $`k(Q)`$ is the curvature and $`l(Q)`$ is the length of the polygon-loop $`Q`$ <sup>2</sup><sup>2</sup>2We shall use the word โloopโ for the โpolygon-loopโ.. This transfer matrix describes the propagation of the initial loop $`Q_1`$ to the final loop $`Q_2`$.
Our aim in this article is to extend these results in several directions. First to construct the transfer matrix for more general case of nonzero self-intersection coupling constant $`\kappa `$
$$K(Q_1,Q_2)=K_{\kappa =0}(Q_1,Q_2)exp\{4\kappa \beta [k_{int}(Q_1)+2l(Q_1\stackrel{}{}Q_2)+k_{int}(Q_2)]\},$$
where $`k_{int}(Q)`$ is number of self-intersection vertices of the loop $`Q`$. For $`\kappa =0`$ this expression coincides with the previous one. In the limit $`\kappa \mathrm{},`$ one can see that on every time slice there will be no loops with self-intersections and we have propagation of nonsingular self-avoiding oriented loops.
Secondly, I will demonstrate that diagonalization of any transfer matrices $`K(Q_1Q_2)`$, which depend only on symmetric difference of loops $`Q_1Q_2`$ can be performed by using generalization of Fourier transformation in loop space. This transformation will be defined by using the loop eigenfunctions
$$\mathrm{\Psi }_P(Q)=e^{i\pi s(PQ)}$$
which are numbered by the loop momentum $`P`$ and $`s(Q)`$ is the area of the region with boundary loop $`Q`$. They are in a pure analogy with plane waves in quantum mechanics $`\psi _p(x)=e^{ipx}`$. We shall prove that any transfer matrix $`K(Q_1Q_2)`$ can be diagonalized by using this loop Fourier transformation
$$K(P_1,P_2)=\underset{\{Q_1,Q_2\}}{}K(Q_1Q_2)e^{i\pi s(P_1Q_1)i\pi s(P_2Q_2)}=\mathrm{\Lambda }_{P_1}\delta (P_1,P_2),$$
here $`\mathrm{\Lambda }_P`$ are eigenvalues depending on loop-momentum $`P`$. Because the transfer matrix plays the role of the evolution operator, one can define a corresponding Hamiltonian operator as $`H(P)=ln\mathrm{\Lambda }(P).`$ At this point one can talk about loop quantum mechanics which should be defined by using conjugate loop operators $`\widehat{Q}`$ and $`\widehat{P}`$ and the Hamiltonian $`H(\widehat{P},\widehat{Q})`$.
I shall also consider the model of random three-dimensional manifolds which are build by gluing together tetrahedra in continuous space $`R^4`$ and out of 3d cubes in 4d Euclidean lattice $`Z^4`$. This system has very close nature with the gonihedric model because in both cases the action is proportional to the linear size of the manifold . In the authors where able to construct a transfer matrix for this theory as well. This transfer matrix describes the propagation of two-dimensional membrane $`M`$ and our main result is that both theories can be solved exactly by using generalization of the Fourier transformation in loop and in membrane spaces. They represent two nontrivial theories in three and four dimensions respectively which can be solved exactly.
In the next sections I shall present background material, the definition of the system of random surfaces in continuous space as well as on Euclidean lattice and describe its basic properties. In the third section the construction of the transfer matrix for the loops will be extended to the case of nonzero self-intersection coupling constant $`\kappa `$. In the fourth main section we shall introduce the generalization of Fourier transformation in the loop space and demonstrate that loop Fourier transformation allows to diagonalize all transfer matrices which depend only on symmetric difference of loops. The net result is that all eigenvalues $`\mathrm{\Lambda }_P`$ of 3d transfer matrix are exactly equal to the loop correlation functions of the corresponding 2d system
$$(\frac{\mathrm{\Lambda }_{P_{\stackrel{}{r_1},\mathrm{}\stackrel{}{r_n}}}}{\mathrm{\Lambda }_{\mathrm{}}})^{3d}=<\sigma _{\stackrel{}{r_1}}\mathrm{}\sigma _{\stackrel{}{r_n}}>_{2d}.$$
(1)
In this formula the loop momentum $`P`$ is defined by the vectors $`\stackrel{}{r_1},\mathrm{}\stackrel{}{r_n}`$ (see Figure 8). In its form the last relation is very close with the ones discussed in .
In the last section I shall consider the transfer matrix for membranes. This matrix also can be diagonalized by Fourier transformation and all its egenvalues are equal to the correlation functions of the corresponding 3d system. In particular we shall see that free energy of the membrane system is equal to the free energy of gonihedric system of loops and is equal to the free energy of 2d Ising model.
## 2 Random Surfaces and Path Integral
Surfaces in Continuous Space. Gonihedric string has been defined as a model of random surfaces with the action which is proportional to the linear size of the surface
$$m\underset{<ij>}{}\lambda _{ij}|\pi \alpha _{ij}|^\zeta \zeta (d2)/d,$$
(2)
where $`\lambda _{ij}`$ is the length of the edge $`<ij>`$ of the triangulated surface $`M`$ and $`\alpha _{ij}`$ is the dihedral angle between two neighboring triangles of $`M`$ sharing a common edge $`<ij>`$ (see Figure 1,2) <sup>3</sup><sup>3</sup>3The angular factor defines the rigidity of the random surfaces and for $`\zeta (d2)/d`$ it increases sufficiently fast near angles $`\alpha =\pi `$ to suppress transverse fluctuations .. The action was defined for the self-intersecting surfaces as well , it accounts self-intersections of different orders by ascribing weights to self-intersections. These weights are proportional to the length of the intersection multiplied by the angular factor which is equal to the sum of all angular factors corresponding to dihedral angles in the intersection (see Figure 3)
$$\kappa m\underset{<ij>}{}\lambda _{ij}(|\pi \alpha _{ij}^1|^\zeta +\mathrm{}+|\pi \alpha _{ij}^r|^\zeta ).$$
(3)
The number of terms inside parentheses depends on the order of the intersection $`r`$, where $`r`$ is the number of triangles sharing the edge $`<ij>`$. Coupling constant $`\kappa `$ is called self-intersection coupling constant . Thus the total action is equal to
$$A(M)=m\underset{<ij>}{}\lambda _{ij}|\pi \alpha _{ij}|^\zeta +m\kappa \underset{<ij>}{}\lambda _{ij}(|\pi \alpha _{ij}^1|^\zeta +\mathrm{}+|\pi \alpha _{ij}^r|^\zeta ).$$
(4)
Both terms in the action have the same dimension and the same geometrical nature, the action (2) measures two-dimensional surfaces in terms of length, self-intersections are also measured in terms of length (3). If $`\kappa `$ is small then self-intersections are permitted, if $`\kappa `$ increases then there is strong repulsion along the curves of self-intersections, and the surfaces tend to be self-avoiding .
The model is a natural extension of the Feynman path integral for point-like relativistic particle (see Figure 1,2). This can be seen from (2) in the limit when the surface degenerates into a single world line <sup>4</sup><sup>4</sup>4This property of the gonihedric action guarantees that spike instability does not appear here because the action is proportional to the total length of the spikes and suppresses the corresponding fluctuations .
$$\underset{<ij>}{}\lambda _{ij}|\pi \alpha |^\varsigma \underset{<ij>}{}\lambda _{ij}.$$
At the classical level the string tension is equal to zero and quarks viewed as open ends of the surface are propagating freely without interaction $`\sigma _{classical}=0`$. This is because the action (2) is equal to the perimeter of the flat Wilson loop $`A(M)m(R+T).`$ As it was demonstrated in , quantum fluctuations generate nonzero string tension $`\sigma _{quantum}=\frac{d}{a^2}(1ln\frac{d}{\beta }),`$ where $`d`$ is the dimension of the spacetime, $`\beta `$ is the coupling constant, $`a`$ is a scaling parameter and $`\varsigma =\frac{d2}{d}`$ in (2). In the scaling limit $`\beta \beta _c=d/e`$ the string tension has a finite limit.
We shall also consider similar model of random membranes with the action which is proportional to the linear size of the corresponding worldvolume. This action for three-dimensional manifold has been defined as
$$A(N)=m\underset{<i,j>}{}\lambda _{ij}|2\pi \omega _{ij}|^\zeta $$
(5)
where $`\lambda _{ij}`$ is the length of the edge $`<ij>`$ of three-dimensional manifold $`N`$ which is constructed by gluing together three-dimensional tetrahedra through their triangular faces and $`2\pi \omega _{ij}`$ is the deficit angle on the edge $`<ij>`$, here $`\omega _{ij}`$ is the sum of the dihedral angles between triangular faces of tetrahedra which have the common edge $`<ij>`$.
Equivalent spin system on a lattice . In addition to the formulation of the theory in the continuum space the system allows an equivalent representation on Euclidean lattices where a surface is associated with a collection of plaquettes . The corresponding statistical weights are proportional to the total number of non-flat edges $`n_2(M)`$ of the surface $`M`$ as it is defined in (2) . The edges are non-flat when the dihedral angle between plaquettes is equal to $`\pi /2`$. The surfaces may have self-intersections in the form of four plaquettes intersecting on a link. The weight associated with self-intersections is proportional to $`4\kappa n_4`$ where $`n_4(M)`$ is the number of edges with four intersecting plaquettes, and $`\kappa `$ is the self-intersection coupling constant (3) . The edges of the surface with self-intersections comprise the singular part of the surface. The edges of the surface where only two plaquettes are intersecting comprise the regular part of the surface.
The partition function can be represented as a sum over two-dimensional surfaces $`M`$ of the type described above, embedded into a three-dimensional lattice:
$$Z(\beta )=\underset{\{surfacesM\}}{}e^{2\beta A(M)},$$
(6)
where $`A(M)`$ is the energy of the surface $`M`$ constructed from plaquettes
$$A_{gonihedric}(M)=n_2(M)+4\kappa n_4(M).$$
(7)
To study statistical and scaling properties of the system one can directly simulate surfaces by gluing together plaquettes following the rules described above, or to use the duality between random surfaces and spin systems on Euclidean lattice. This duality allows to study equivalent spin system with specially adjusted interaction between spins . In three dimensions the corresponding Hamiltonian is equal to
$$H_{gonihedric}^{3d}=2\kappa \underset{\stackrel{}{r},\stackrel{}{\alpha }}{}\sigma _\stackrel{}{r}\sigma _{\stackrel{}{r}+\stackrel{}{\alpha }}+\frac{\kappa }{2}\underset{\stackrel{}{r},\stackrel{}{\alpha },\stackrel{}{\beta }}{}\sigma _\stackrel{}{r}\sigma _{\stackrel{}{r}+\stackrel{}{\alpha }+\stackrel{}{\beta }}\frac{1\kappa }{2}\underset{\stackrel{}{r},\stackrel{}{\alpha },\stackrel{}{\beta }}{}\sigma _\stackrel{}{r}\sigma _{\stackrel{}{r}+\stackrel{}{\alpha }}\sigma _{\stackrel{}{r}+\stackrel{}{\alpha }+\stackrel{}{\beta }}\sigma _{\stackrel{}{r}+\stackrel{}{\beta }},$$
(8)
where $`\stackrel{}{r}`$ is a three-dimensional vector on the lattice $`Z^3`$, the components of which are integer and $`\stackrel{}{\alpha }`$, $`\stackrel{}{\beta }`$ are unit vectors parallel to the axes. The interface energy of this lattice spin system coincides with the linear action (4), (7).
The form of the Hamiltonian $`H(\kappa )`$ and the symmetry of the system essentially depend on $`\kappa `$: when $`\kappa 0`$ one can flip the spins on arbitrary parallel layers and thus the degeneracy of the ground state is equal to $`32^N`$, where $`N^3`$ is the size of the lattice. When $`\kappa =0`$ the system has even higher symmetry, all states, including the ground state are exponentially degenerate. This degeneracy of the ground state is equal to $`2^{3N}`$ . This is because now one can flip the spins on arbitrary layers, even on intersecting ones. The corresponding Hamiltonian contains only exotic four-spin interaction term.
The understanding of critical behaviour of this system comes from the analysis of the transfer matrix which describes the propagation of the closed loops $`Q`$ in time direction with an amplitude which is proportional to the sum of the length $`l(Q)`$ of the loop and of its total curvature $`k(Q)`$. Because the amplitude essentially depends on length, as it takes place in 2d Ising model, this result directs to the idea that the system in 3d will show the phase transition which should be of the same nature as it is in 2d Ising system and that the transfer matrix describes the propagation of an almost free string of Ising model . Direct numerical Monte-Carlo simulations and analytical solution of different approximations to the full transfer matrix confirm strong connection with the 2d Ising model .
Below I shall remind the derivation of the corresponding transfer matrix for the case $`\kappa =0`$, which is based on geometrical theorem proven in and then I shall derive the transfer matrix for more general case $`\kappa 0`$. This will allow to study physical picture of string propagation which follows from the transfer matrix approach.
## 3 Loop Transfer Matrix
Geometrical Theorem . In this section I shall review the construction of the transfer matrix for the system (8) with $`\kappa =0`$
$$H_{gonihedric}^{3d}(\kappa =0)=\frac{1}{2}\underset{\stackrel{}{r},\stackrel{}{\alpha },\stackrel{}{\beta }}{}\sigma _\stackrel{}{r}\sigma _{\stackrel{}{r}+\stackrel{}{\alpha }}\sigma _{\stackrel{}{r}+\stackrel{}{\alpha }+\stackrel{}{\beta }}\sigma _{\stackrel{}{r}+\stackrel{}{\beta }},$$
(9)
and then extend the result to nonzero case $`\kappa 0`$.
In order to find transfer matrix for this system we have to use geometrical theorem proven in . The geometrical theorem provides an equivalent representation of the action $`A(M)`$ in terms of total curvature $`k(E)`$ of the polygons which appear in the intersection of the two-dimensional plane $`E`$ with the given two-dimensional surface $`M`$ (see Figure 4,5). This curvature $`k(E)`$ should then be integrated over all planes $`E`$ intersecting the surface $`M`$.
The intersection of the two-dimensional plane $`E`$ with the polyhedral surface $`M`$ is the union of the polygons (see Figure 5)
$$Q_1(E),\mathrm{},Q_k(E)$$
and the absolute total curvature $`k(E)`$ of these polygons is equal to
$$k(E)=\underset{i=1}{\overset{k}{}}k(Q_i)=\underset{<i,j>}{}|\pi \alpha _{ij}^E|,$$
(10)
where $`\alpha _{ij}^E`$ are the angles of these polygons. They are defined as the angles in the intersection of the two-dimensional plane $`E`$ with the edge $`<ij>`$ (see Figure 5). The meaning of (10) is that it measures the total revolution of the tangent vectors to polygons $`Q_1(E),\mathrm{},Q_k(E)`$. Integrating the total curvature $`k(E)`$ (10) over all intersecting planes $`E`$ we shall get the action $`A(M)`$
$$A(M)=\frac{1}{2\pi }_{\{E\}}k(E)๐E.$$
(11)
Geometrical Theorem on a lattice . One can find the same representation (11) for the action $`A(M)`$ on a cubic lattice $`Z^3`$ introducing a set of planes $`\{E_x\},\{E_y\},\{E_z\}`$ perpendicular to $`x,y,z`$ axis on the dual lattice $`Z_d^3`$. These planes will intersect a given surface $`M`$ and on each of these planes we shall have an image of the surface $`M`$. Every such image is represented as a collection of closed polygons $`Q(E)`$ appearing in the intersection of the plane with surface $`M`$ (see Figure 6). The energy of the surface $`M`$ is equal now to the sum of the total curvature $`k(E)`$ of all these polygons on different planes
$$A(M)=\underset{\{E_x,E_y,E_z\}}{}k(E).$$
(12)
The total curvature $`k(E)`$ is simply the total number of polygon right angles. In the case when the self-intersection coupling constant $`\kappa =0`$ is equal to zero we must compute the total curvature $`k(E)`$ of these polygons ignoring angles at the self-intersection points!
With (12) the partition function of the system can be written in the form
$$Z(\beta )=\underset{\{M\}}{}exp\{2\beta \underset{\{E\}}{}k(E)\}$$
(13)
where the sum in the exponent can be represented as a product
$$exp\{2\beta \underset{\{E\}}{}k(E)\}=\underset{\{E\}}{}e^{2\beta k(E)}=\underset{\{E_z\}}{}e^{2\beta k(E_z)}\underset{\{E_y\}}{}e^{2\beta k(E_y)}\underset{\{E_x\}}{}e^{2\beta k(E_x)}.$$
(14)
The goal is to express the energy functional (12) and the product (14) in terms of quantities defined only on two-dimensional planes in one fixed direction, letโs say through $`\{E_z\}`$.
Transfer Matrix for $`\kappa =0`$ . The question is: what kind of information do we need to know on these planes $`\{E_z\}`$ to recover the values of the total curvature $`k(E_x)`$ and $`k(E_y)`$ on the planes $`\{E_y\}`$ and $`\{E_x\}`$ for the given surface $`M`$?
Let us consider the sequence of two planes $`E_z^i`$ and $`E_z^{i+1}`$ and denote by $`Q_i`$ the polygon-image of the surface $`M`$ on the plane $`E_z^i`$ and by $`Q_{i+1}`$ the polygon-image of $`M`$ on the plane $`E_z^{i+1}`$. The contribution to the curvature $`k(E_x)+k(E_y)`$ of the polygons which are on the perpendicular planes between $`E_z^i`$ and $`E_z^{i+1}`$ is equal to the length of the polygons $`Q_i`$ and $`Q_{i+1}`$ without length of the common bonds
$$l(Q_i)+l(Q_{i+1})2l(Q_iQ_{i+1})=l(Q_iQ_{i+1}),$$
(15)
where the polygon-loop $`Q_1Q_2Q_1Q_2\backslash Q_1Q_2`$ is a union of links $`Q_1Q_2`$ without common links $`Q_1Q_2`$. Therefore (see formula (12))
$$A(M)=\underset{\{E\}}{}k(E)=\underset{\{E_z\}}{}k(Q_i)+l(Q_iQ_{i+1}).$$
(16)
The partition function (13) can be now represented in the form
$$Z(\beta )=\underset{\{Q_1,Q_2,\mathrm{},Q_N\}}{}K_\beta (Q_1,Q_2)\mathrm{}K_\beta (Q_N,Q_1)=trK_\beta ^N,$$
(17)
where $`K_\beta (Q_1,Q_2)`$ is the transfer matrix of size $`\gamma \times \gamma `$, defined as
$$K_{\kappa =0}^{gonihedric}(Q_1,Q_2)=exp\{\beta [k(Q_1)+2l(Q_1Q_2)+k(Q_2)]\},$$
(18)
where $`Q_1`$ and $`Q_2`$ are closed polygon-loops on a two-dimensional lattice $`T^2`$ of size $`N\times N`$ and $`\gamma =2^{N^2}`$ is the total number of polygon-loops on a lattice $`T^2`$. The transfer matrix (18) can be viewed as describing the propagation of the polygon-loop $`Q_1`$ at time $`\tau `$ to another polygon-loop $`Q_2`$ at the time $`\tau +1`$ (see Figure 7)<sup>5</sup><sup>5</sup>5Layer-to-layer transfer matrices for three-dimensional statistical systems have been considered in the literature . Using Yang-Baxter and Tetrahedron equations one can compute the spectrum of the transfer matrix in a number of cases . In the given case the transfer matrix has geometrical interpretation which helps to compute the spectrum..
The functional $`k(Q)`$ is the total curvature of the polygon-loop $`Q`$ which is equal to the number of corners of the polygon ( the vertices with self-intersection are not counted!) and $`l(Q)`$ is the length of $`Q`$ which is equal to the number of its links. The length functional $`l(Q_1Q_2)`$ was defined above (15) . The operation $``$ maps two polygon-loops $`Q_1`$ and $`Q_2`$ into a polygon-loop $`Q=Q_1Q_2`$ <sup>6</sup><sup>6</sup>6Note that the operations $``$ and $``$ do not have this property. These operations acting on a polygon-loops can produce link configurations which do not belong to $`\mathrm{\Pi }`$. The symmetric difference of sets $`Q_1Q_2`$ is an important concept in functional analysis .. The length functional $`l(Q_1Q_2)`$ defines a distance between two polygon-loops $`Q_1`$ and $`Q_2`$.
The eigenvalues of the transfer matrix $`K_\beta (Q_1,Q_2)`$ define all statistical properties of the system and can be found as a solution of the following integral equation in the loop space $`\mathrm{\Pi }`$
$$\underset{\{Q_2\}}{}K_\beta (Q_1,Q_2)\mathrm{\Psi }(Q_2)=\mathrm{\Lambda }(\beta )\mathrm{\Psi }(Q_1),$$
(19)
where $`\mathrm{\Psi }(Q)`$ is a function on loop space. The Hilbert space of complex functions $`\mathrm{\Psi }(Q)`$ on $`\mathrm{\Pi }`$ will be denoted as $`H=L^2(\mathrm{\Pi })`$. The eigenvalues define the partition function (13)
$$Z^{3d}(\beta )=\mathrm{\Lambda }_0^N+\mathrm{}+\mathrm{\Lambda }_{\gamma 1}^N,$$
(20)
and in the thermodynamical limit the free energy is equal to
$$\beta f_{3d}(\beta )=\underset{N\mathrm{}}{lim}\frac{1}{N^3}lnZ^{3d}(\beta )=\underset{N\mathrm{}}{lim}\frac{1}{N^2}ln\mathrm{\Lambda }_0.$$
(21)
Finite time propagation amplitude of an initial loop $`Q_i`$ to a final loop $`Q_f`$ for the time interval $`t=M/\beta `$ can be defined as
$$K(Q_i,Q_f)=\mathrm{\Lambda }_0^M\underset{\{Q_1,Q_2,\mathrm{},Q_{M1}\}}{}K_\beta (Q_i,Q_1)\mathrm{}K_\beta (Q_{M1},Q_f),$$
(22)
where we have introduced natural normalization to the biggest eigenvalue $`\mathrm{\Lambda }_0`$ and $`MN`$.
Loop Transfer Matrix for $`\kappa 0`$. We would like to generalize the construction of the transfer matrix to the case of nonzero self-intersection coupling constant $`\kappa `$. As we already mentioned above surfaces on a lattice may have self-intersections and we associate the coupling constant $`\kappa `$ with these self-intersection edges of the surface. When we consider the sections of the lattice surface $`M`$ by the planes $`E`$ the self-intersection edges will appear on a plane as polygon-loops with self-intersections. When $`\kappa =0`$ we have been instructed to compute the total curvature $`k(E)`$ of these polygon-loops simply ignoring angles at the self-intersection points. Now when $`\kappa 0`$ we have to take into account self-intersection vertices of polygon-loops.
The contribution to the action which comes from simple polygon right angles is the same as before (16)
$$\underset{\{E_z\}}{}k(Q_i)+l(Q_iQ_{i+1}).$$
(23)
The contribution coming from self-intersections should be divided into contribution coming from vertical and horizontal edges of the surface. The vertical and horizontal edges of the surface are defined relative to the โtimeโ planes $`E_z`$. The contribution from vertical edges is simply the number of self-intersection vertices of loops on $`E_z`$ planes
$$4\kappa \underset{\{E_z\}}{}k_{int}(Q_i),$$
(24)
where $`k_{int}(Q)`$ is the number of self-intersection vertices of loop $`Q`$.
To count contribution from horizontal edges of the surface one should introduce the orientation of the polygon-loops on every $`E_z`$ plane. The same orientation is ascribed to all polygon-loops on a given plane $`E_z`$: clock-wise or anti-clock-wise. Orientation is automatic if $`Q`$ is thought as interface between spins. Conjugation $`\stackrel{}{Q}`$ $``$ $`\stackrel{}{Q}`$ is equivalent to $`Z_2`$ spin flip. So there is $`\stackrel{}{Q}`$ and $`\stackrel{}{Q}`$ with opposite orientation.
Now let us consider oriented polygon-loops on two planes $`E_z^i`$ and $`E_z^{i+1}`$. Then one should select that bonds on these polygon-loops which are common. Only bonds with the opposite orientation should be counted. The number of common bonds with opposite orientation we shall denote as $`l(Q_i\stackrel{}{}Q_j)`$. The notation $``$ is used to denote the common bonds, the symbol $``$ on the top of it to denote that they should have opposite orientations. Therefore the contribution from horizontal edges can be expressed as
$$4\kappa \underset{\{E_z\}}{}l(Q_i\stackrel{}{}Q_j)$$
(25)
and the contribution from all self-intersections as
$$4\kappa \underset{\{E_z\}}{}k_{int}(Q_i)+l(Q_i\stackrel{}{}Q_{i+1}).$$
(26)
The total action (7) now takes the form:
$$A(M)=\underset{\{E_z\}}{}k(Q_i)+l(Q_iQ_{i+1})+4\kappa [k_{int}(Q_i)+l(Q_i\stackrel{}{}Q_{i+1})].$$
(27)
The partition function (13) can be now represented in the same form (17) where $`K_\beta (Q_1,Q_2)`$ is again the transfer matrix of size $`\gamma \times \gamma `$, defined as
$`K_\beta (Q_1,Q_2)=exp\{\beta [k(Q_1)+2l(Q_1Q_2)+k(Q_2)]\}`$
$`exp\{4\kappa \beta [k_{int}(Q_1)+2l(Q_1\stackrel{}{}Q_2)+k_{int}(Q_2)]\},`$ (28)
where $`Q_1`$ and $`Q_2`$ are closed oriented polygon loops on a two-dimensional lattice $`T^2`$ of size $`N\times N`$ and $`\gamma `$ is the total number of polygon-loops on a lattice $`T^2`$. Let us now consider the limit
$$\kappa \mathrm{},$$
then from our expression for the transfer matrix (28) one can conclude that on every time slice there will be no polygon-loops with self-intersections. This means that they are just nonsingular self-avoiding oriented loops. In addition there is strong correlation between loops on consecutive planes: common bonds with opposite orientation are also strongly suppressed. Which means that propagation of oriented loop in time has small probability to invert its orientation.
## 4 Loop Fourier transformation
In the approximation when we drop the curvature term in the expression (18) for the loop transfer matrix $`K(Q_1,Q_2)`$ the spectrum of the matrix has been evaluated exactly in the articles . In this section I shall present an alternative method of derivation which is based on generalization of Fourier transformation and allows to carry out analogy with quantum mechanics of point particles and to emphasize generality of the approach.
I shall consider two different approximations to the initial transfer matrix (18) suggested in and . In the first one it was suggested to substruct the term $`2k(Q_1Q_2)`$ and to consider the transfer matrix of the form
$$\stackrel{~}{K}(Q_1Q_2)=exp\{\beta [k(Q_1Q_2)+2l(Q_1Q_2)]\},$$
(29)
and in the second one to consider the transfer matrix
$$K(Q_1Q_2)=exp\{2\beta l(Q_1Q_2)\}.$$
(30)
In both cases the transfer matrix depends only on the symmetric difference $`Q_1Q_2`$ and the integral equation (19) take the form
$$\underset{\{Q_2\}}{}K_\beta (Q_1Q_2)\mathrm{\Psi }(Q_2)=\mathrm{\Lambda }(\beta )\mathrm{\Psi }(Q_1).$$
(31)
The fact that transfer matrix now depends only on the symmetric difference is of paramount importance in solving these models. The reason is that we can make a change of variable, from $`Q_2`$ to $`Q`$ at fixed $`Q_1`$
$$Q=Q_1Q_2$$
(32)
and see that equation transforms to the form
$$\underset{\{Q\}}{}K_\beta (Q)\mathrm{\Psi }(QQ_1)=\mathrm{\Lambda }(\beta )\mathrm{\Psi }(Q_1).$$
(33)
The integral equation will be solved if we can find the set of loop functions $`\mathrm{\Psi }(Q)`$ which satisfy the following functional equation on loop space $`\mathrm{\Pi }`$
$$\mathrm{\Psi }(QQ_1)=\mathrm{\Psi }(Q)\mathrm{\Psi }(Q_1).$$
(34)
For these functions the equation (33) reduces to
$$\underset{\{Q\}}{}K_\beta (Q)\mathrm{\Psi }(Q)=\mathrm{\Lambda }(\beta )$$
(35)
and we can find all eigenvalues. The constant function $`\mathrm{\Psi }=1`$ is the solution of the functional equation (34) and due to Frobenius-Perron theorem corresponds to the largest and non-degenerate eigenvalue $`\mathrm{\Lambda }_{\mathrm{}}`$ of the symmetric transfer matrix $`K(Q_1Q_2)`$.
The functional equation (34) is analogous to the equation $`\psi (x+y)=\psi (x)\psi (y)`$ appearing in ordinary quantum mechanics, which allows to find all eigenvectors and eigenvalues of the integral equation
$$K_\beta (xy)\psi (y)๐x=\lambda (\beta )\psi (y),K_\beta (xy)=exp(\beta (xy)^2)$$
(36)
with the well known solutions
$$\psi (x)=e^{ipx},\lambda _p(x)=e^{p^2/4\beta }.$$
Our idea is to find all solutions of the functional equation (34) and the analog of the plane wave solution in loop space. From functional equation (34) it follows that if $`Q=Q_1`$ then
$$\mathrm{\Psi }(\mathrm{})=\mathrm{\Psi }^2(Q)$$
(37)
and if $`Q=\mathrm{}`$ then $`\mathrm{\Psi }(\mathrm{})`$ is equal to zero or to one, therefore nontrivial solution is
$$\mathrm{\Psi }^2(Q)=1\mathrm{\Psi }(Q)=\pm 1.$$
(38)
This suggests the following form of the eigenfunctions
$$\mathrm{\Psi }_P(Q)=e^{i\pi s(PQ)}$$
(39)
which are numbered by the loop momentum $`P`$ and take values $`\pm 1`$. The $`s(Q)`$ is the area of the region with the boundary loop $`Q`$. These solutions are in pure analogy with plane waves in quantum mechanics . To check that they are indeed solutions of the functional equation (34) we have to use the set identity
$$(Q_1Q_2)P=(Q_1P)(Q_2P)$$
(40)
which is easy to verify and the relation $`s(Q_i)+s(Q_{i+1})2s(Q_iQ_{i+1})=s(Q_iQ_{i+1})`$. To carry out the analogy with plane wave solution of quantum mechanics deeper one should interpret the loop $`P`$ as a loop momentum conjugate to a loop coordinate variable $`Q`$. To prove orthogonality of these loop wave functions let us compute the product
$$\mathrm{\Psi }_{P_1}\mathrm{\Psi }_{P_2}=\underset{\{Q\}}{}\mathrm{\Psi }_{P_1}(Q)\mathrm{\Psi }_{P_2}(Q)=\underset{\{Q\}}{}e^{i\pi s(P_1Q)+i\pi s(P_2Q)}=\underset{\{Q\}}{}e^{i\pi s((P_1P_2)Q)}$$
(41)
where we again used the above set identity and the relation (15) for $`s(Q)`$. The last sum is equal to zero if $`P_1P_2\mathrm{}`$ and is equal to the number of loops $`\gamma `$ on the lattice if $`P_1P_2=\mathrm{}`$, thus
$$\underset{\{Q\}}{}e^{i\pi s((P_1P_2)Q)}=2^{N^2}\delta (P_1,P_2).$$
(42)
The last statement follows from the fact that the product between the constant function $`\mathrm{\Psi }_{\mathrm{}}=1`$, corresponding to the largest and non-degenerate eigenvalue $`\mathrm{\Lambda }_{\mathrm{}}`$, and all others $`exp(i\pi s(PQ))`$ is equal to zero if $`P`$ is nonempty
$$\underset{\{Q\}}{}1e^{i\pi s((PQ)}=0,ifP\mathrm{}.$$
(43)
The number of functions in this orthogonal set is equal therefore to the number of closed loops $`P`$, that is $`\gamma =2^{N^2}`$.
Having in hand all eigenfuctions we can find all eigenvalues from (35) for transfer matrices (29) and (30)
$$\stackrel{~}{\mathrm{\Lambda }}_P=\underset{\{Q\}}{}e^{i\pi s(PQ)\beta [k(Q)+2l(Q)]}\mathrm{\Lambda }_P=\underset{\{Q\}}{}e^{i\pi s(PQ)2\beta l(Q)}.$$
(44)
Because the expression $`_{\{Q\}}e^{2\beta l(Q)}`$ is the partition function of the 2d Ising model and $`_{\{Q\}}e^{\beta [2l(Q)+k(Q)]}`$ is the partition function $`\stackrel{~}{Z}^{2d}`$ of the two-dimensional model solved in we see that the largest eigenvalue $`\mathrm{\Lambda }_{\mathrm{}}`$ is exactly equal to a corresponding partition function
$$\stackrel{~}{\mathrm{\Lambda }}_{\mathrm{}}=\stackrel{~}{Z}^{2d}\mathrm{\Lambda }_{\mathrm{}}=Z^{2dIsing}.$$
(45)
For the ratio of eigenvalues we shall have <sup>7</sup><sup>7</sup>7In what follows we shall consider the transfer matrix (30). Analogous formulas are valid for the transfer matrix (29) with curvature term $`k(Q)`$.
$$\frac{\mathrm{\Lambda }_P}{\mathrm{\Lambda }_{\mathrm{}}}=\underset{\{Q\}}{}e^{i\pi s(PQ)2\beta l(Q)}/Z^{2dIsing}<e^{i\pi s(PQ)}>_Q$$
(46)
which is identically equal to the spin correlation functions of the 2d Ising model
$$\frac{\mathrm{\Lambda }_{P_{\stackrel{}{r_1},\mathrm{}\stackrel{}{r_n}}}}{\mathrm{\Lambda }_{\mathrm{}}}=<\sigma _{\stackrel{}{r_1}}\mathrm{}\sigma _{\stackrel{}{r_n}}>_{2dIsing}.$$
(47)
In this formula the loop momentum $`P`$ is defined by the vectors $`\stackrel{}{r_1},\mathrm{}\stackrel{}{r_n}`$. In particular for the smallest (one-box) loop momentum $`P=\mathrm{}`$ we shall have
$$\frac{\mathrm{\Lambda }_{\mathrm{}}}{\mathrm{\Lambda }_{\mathrm{}}}=<\sigma _\stackrel{}{r}>=\mu (\beta )$$
(48)
which coincides with the magnetization of 2d Ising model. Because it does not depend on position of the loop on the lattice, this eigenvalue is $`N^2`$ degenerate. For the two-box loop momentum $`P=`$ consisting of two neighboring boxes we shall have
$$\frac{\mathrm{\Lambda }_{}}{\mathrm{\Lambda }_{\mathrm{}}}=<\sigma _\stackrel{}{r}\sigma _{\stackrel{}{r+1}}>=u(\beta )$$
(49)
which is $`2N^2`$ degenerate and after summation over all $`2N^2`$ bonds of the lattice we shall get the internal energy of the 2d Ising model.
We can also use loop wave functions to define generalized loop-Fourier transformation as
$$\mathrm{\Psi }(P)=\frac{1}{2^{N^2}}\underset{\{Q\}}{}e^{i\pi s(PQ)}\mathrm{\Psi }(Q)$$
(50)
and then compute the loop Fourier transformation of the transfer matrices (29) and (30) as:
$`K(P_1,P_2)={\displaystyle \underset{\{Q_1,Q_2\}}{}}K(Q_1Q_2)e^{i\pi s(P_1Q_1)i\pi s(P_2Q_2)}=\mathrm{\Lambda }_{P_1}{\displaystyle \underset{\{Q_2\}}{}}e^{i\pi s(P_1Q_2)i\pi s(P_2Q_2)}`$
$`=\mathrm{\Lambda }_{P_1}\delta (P_1,P_2).`$ (51)
The finite time propagation amplitude of an initial loop $`P_i`$ to the final loop $`P_f`$ for the time interval $`t=M/\beta `$ (22) we have
$`K(P_i,P_f)=({\displaystyle \frac{\mathrm{\Lambda }_{P_i}}{\mathrm{\Lambda }_{\mathrm{}}}})^M\delta (P_i,P_f)=e^{Mln(\frac{\mathrm{\Lambda }_{P_i}}{\mathrm{\Lambda }_{\mathrm{}}})}\delta (P_i,P_f)=e^{MH(P_i)}\delta (P_i,P_f),`$ (52)
therefore for the loop Hamiltonian we shall get
$$H(P)=ln(\frac{\mathrm{\Lambda }_P}{\mathrm{\Lambda }_{\mathrm{}}})$$
(53)
where $`H(P)`$ is the energy of the coherent state with loop momentum $`P`$.
Spin system which corresponds to the matrix $`K(Q_iQ_{i+1})`$ (30). The natural question which arises here is how nontrivial is the 3d system which has been solved in this way? To make the answer as clear as possible we should go back and express the system which has been solved again in terms of spins on a 3d lattice and to see what is left out of original interactions.
We ignore the curvature terms $`k(Q_i)`$ in the transfer matrix, this means that we ignore the contribution coming from the edges of the surface which are along the time direction, or $`z`$ direction. The contribution coming from the edges which are parallel to the $`x`$ and $`y`$ direction are left and their contribution is equal to (15)
$$\underset{\{E_z\}}{}l(Q_iQ_{i+1}).$$
(54)
In terms of spin variables the original action can be expressed as a sum of energies coming from all edges of the surface (see formulas (5a) and (7) in )
$$H=\underset{overalledges}{}H_{edge},$$
(55)
where
$$H_{edge}=U_1U_1+U_1U_1=\sigma _1\sigma _2\sigma _1\sigma _2.$$
(56)
is a four-spin interaction term with spins distributed around the edge. Because we ignore the edges in the $`z`$ direction, we have to sum in (55) only over $`A_{edges}`$ in $`x`$ and $`y`$ direction. This corresponds to the summation over four-spin interaction terms (56) which lie only on the planes $`E_x`$ and $`E_y`$. These planes are normal to the $`x`$ and $`y`$ direction. Therefore our approximation corresponds to the 3d lattice spin system in which four-spin interactions take place only on vertical planes $`E_x`$ and $`E_y`$
$$H_{Q_1Q_2}=\underset{E_x,E_y}{}\sigma \sigma \sigma \sigma .$$
(57)
It is obvious that this spin system can not be factored into non interacting two-dimensional subsystems. We can ask what is the effective interaction between spins which lie on a given two-dimensional $`E_z`$ plane? Let us label spins which are on that plane as $`\sigma `$, spins which are on the previous plane as $`\lambda `$, and spins which are on the next plane as $`\mu `$. The part of the Hamiltonian which contains only $`\sigma `$ spins can be written in the form:
$$\underset{i}{}\sigma _i\sigma _{i+1}\lambda _i\lambda _{i+1}+\sigma _i\sigma _{i+1}\mu _i\mu _{i+1}=\underset{i}{}(\lambda _i\lambda _{i+1}+\mu _i\mu _{i+1})\sigma _i\sigma _{i+1}=\underset{i}{}J_{eff}\sigma _i\sigma _{i+1},$$
(58)
where the effective coupling $`J_{eff}`$ depends on spins on the neighboring planes. It takes the values $`2,0,2`$ with probabilities $`1/6,4/6,1/6`$ simply because
$$J_{eff}=\lambda _i\lambda _{i+1}+\mu _i\mu _{i+1}.$$
For the observer who lives on the plane $`E_z`$ spin interactions look like 2d Ising model with randomly distributed coupling constant $`J_{eff}`$.
3d Ising Transfer Matrix. Similar expression for the transfer matrix has been derived in for the 3d Ising model
$$K_\beta ^{3dIsing}(P_1,P_2)=exp\{\beta [l(P_1)+2s(P_1P_2)+l(P_2)]\},$$
(59)
where $`s(P)`$ is the area functional. It is instructive to compare 3d-gonihedric and 3d Ising systems. One can see simple hierarchical structure of the corresponding transfer matrices (18) , (59). If one orders the functionals in the increasing geometrical complexity
$$curvaturek(P),lengthl(P),areas(P),$$
(60)
then it becomes clear that gonihedric transfer matrix depends on the first two functionals: curvature and length, while the 3d Ising model depends on the last two functionals: length and area. Thus if one puts the models in order of increasing complexity one should consider gonihedric model and then Ising model .
Let us also consider similar to (29), (30) approximations to the 3d Ising transfer matrix, they are:
$$\stackrel{~}{K}(Q_1Q_2)=exp\{\beta [l(Q_1Q_2)+2s(Q_1Q_2)]\},K(Q_1Q_2)=exp\{\beta 2s(Q_1Q_2)\},$$
(61)
Both approximations can be solved by using the set of eigenfunctions (39)
$$\stackrel{~}{\mathrm{\Lambda }}_P=\underset{\{Q\}}{}e^{i\pi s(PQ)\beta [l(Q)+2s(Q)]}\mathrm{\Lambda }_P=\underset{\{Q\}}{}e^{i\pi s(PQ)2\beta s(Q)}.$$
(62)
The first relation expresses the eigenvalues through the 2d Ising model in background magnetic field, the model which has not been solved yet exactly , and in the second case through free spins in background magnetic field.
## 5 Matrices depending only on symmetric difference
It is obvious that the main reason why we have been able to reduce solution of three-dimensional statistical system to a solution of two-dimensional system is special dependence of the transfer matrix argument from symmetric difference $`Q_1Q_2`$.
We are aimed to formulate the following general result: all transfer matrices of three-dimensional statistical systems which depend only on symmetric difference $`Q_1Q_2`$
$$Z^{3d}(\beta )=TrK^N=e^{\beta f_{3d}(\beta )N^3},K(Q_1Q_2)=e^{\beta \mathrm{\Omega }(Q_1Q_2)}$$
(63)
where $`\mathrm{\Omega }`$ is arbitrary energy functional, can be diagonalized by using orthogonal set of functions $`\mathrm{\Psi }_P(Q)=e^{i\pi s(PQ)}.`$ The eigenvalues can be found exactly in the same way as we described above and are equal to the following averages
$$\mathrm{\Lambda }_P=\underset{\{Q\}}{}e^{i\pi s(PQ)}K(Q)=\underset{\{Q\}}{}e^{i\pi s(PQ)\beta \mathrm{\Omega }(Q)}.$$
(64)
The eigenvalues $`\mathrm{\Lambda }_P`$ of three-dimensional system are exactly equal to the correlation functions of the two-dimensional system with the partition function
$$Z^{2d}(\beta )=\underset{\{Q\}}{}e^{\beta \mathrm{\Omega }(Q)}=e^{\beta f_{2d}(\beta )N^2}=\lambda _0^N+\mathrm{}+\lambda _{2^N1}^N,$$
(65)
where $`f_{2d}(\beta )`$ is the free energy of the 2d model and $`\lambda _i`$ are the eigenvalues of the corresponding transfer matrix. In particular from (64) we see that the largest eigenvalue $`\mathrm{\Lambda }_{\mathrm{}}`$ is equal to the partition function of $`2d`$ system
$$\mathrm{\Lambda }_{\mathrm{}}=\underset{\{Q\}}{}e^{\beta \mathrm{\Omega }(Q)}Z^{2d}(\beta )=\lambda _0^N+\mathrm{}+\lambda _{2^N1}^N,$$
(66)
therefore we can express free energy of 3d system in terms of largest eigenvalue $`\mathrm{\Lambda }_{\mathrm{}}`$ and then through the partition function of the corresponding two-dimensional system
$`\beta f_{3d}(\beta )=\underset{N\mathrm{}}{lim}{\displaystyle \frac{1}{N^3}}lnZ^{3d}(\beta )=\underset{N\mathrm{}}{lim}{\displaystyle \frac{1}{N^3}}ln\{\mathrm{\Lambda }_{\mathrm{}}^N+\mathrm{}\}`$
$`=\underset{N\mathrm{}}{lim}{\displaystyle \frac{1}{N^3}}ln\{[\lambda _0^N+\mathrm{}+\lambda _{2^N1}^N]^N+\mathrm{}\}=\beta f_{2d}(\beta ).`$ (67)
Thus we have amazing equality
$$f_{3d}(\beta )=f_{2d}(\beta )$$
which proves identical critical behaviour of 2d and 3d systems! Because
$$\underset{\{Q\}}{}e^{i\pi s(PQ)\beta \mathrm{\Omega }(Q)}/Z^{2d}(\beta )=<e^{i\pi s(PQ)}>_Q==<\sigma _{\stackrel{}{r_1}}\mathrm{}\sigma _{\stackrel{}{r_n}}>_{2d},$$
from (64) it follows that the ration of eigenvalues is equal to:
$$\frac{\mathrm{\Lambda }_{P_{\stackrel{}{r_1},\mathrm{}\stackrel{}{r_n}}}}{\mathrm{\Lambda }_{\mathrm{}}}=<\sigma _{\stackrel{}{r_1}}\mathrm{}\sigma _{\stackrel{}{r_n}}>_{2d}.$$
(68)
In this formula the loop momentum $`P`$ is defined by the vectors $`\stackrel{}{r_1},\mathrm{}\stackrel{}{r_n}`$. It is clear that the correlation functions of the three-dimensional system are more complicated conglomerates.
It is easy to understand that the result is even more general: any $`ddimensional`$ system which has transfer matrix depending only on symmetric difference $`Q_1Q_2`$ can be reduced to a corresponding $`(d1)dimensional`$ system and its eigenvalues are simply the correlation functions of the lower dimensional system. Below we shall consider an important example of such dimensional reduction in four-dimensions.
Transfer Matrix for Membranes In the authors have constructed the transfer matrix in four-dimensions which describes a propagation of two-dimensional membrane. Below we shall see that this system also can be solved in terms of three-dimensional system considered in previous sections.
Transfer matrix which describes the propagation of two-dimensional membrane $`M`$ has the following form
$$K_\beta (M_1,M_2)=exp\{\beta [\chi (M_1)+2A(M_1M_2)+\chi (M_2)]\},$$
(69)
where $`\chi (M)`$ defines the module of the Euler character of the surface $`M`$
$$\chi (M)=\underset{<i>}{}|2\pi \omega _i|,$$
(70)
summation is over all vertices of triangulated surface $`M`$ and $`A(M)`$ is the gonihedric functional which we defined earlier. If we introduce the set of orthogonal functions
$$\mathrm{\Psi }_P(M)=e^{i\pi v(PM)},$$
(71)
where $`v(M)`$ is the volume of the region whose boundary is the surface $`M`$, then transfer matrices which depend only on symmetric difference of surfaces $`M_1`$ and $`M_2`$
$$\stackrel{~}{K}(M_1,M_2)=exp\{\beta [\chi (M_1M_2)+2A(M_1M_2)]\},K(M_1,M_2)=exp\{2\beta A(M_1M_2)\},$$
(72)
can be diagonalized by using the above set of functions (71)
$$\stackrel{~}{\mathrm{\Lambda }}_P=\underset{\{M\}}{}e^{i\pi v(PM)\beta [\chi (M)+2A(M)]},\mathrm{\Lambda }_P=\underset{\{M\}}{}e^{i\pi v(PM)2\beta A(M)}.$$
(73)
For $`P=\mathrm{}`$ we have
$$\mathrm{\Lambda }_{\mathrm{}}=\underset{\{M\}}{}e^{2\beta A(M)}$$
(74)
and it coincides with the partition function of the gonihedric system (6). The free energy of this four-dimensional model is equal therefore to
$$f_{4d}(\beta )=f_{3d}(\beta )=f_{2dIsing}(\beta )$$
(75)
and the eigenvalues (73) are the correlation function of the 3d gonihedric system
$$(\frac{\mathrm{\Lambda }_P}{\mathrm{\Lambda }_{\mathrm{}}})^{4d}=(Correlationfunctions)^{3d}.$$
(76)
Thus the eigenvalues of the transfer matrix in four-dimensions can be described by the correlation functions of the system of one less dimension.
## 6 Discussion
We have seen that transfer matrix approach allows to solve large class of statistical systems when transfer matrix depends on symmetric difference of propagating loops, membranes or p-branes. The eigenvalues are equal to the correlation functions of the corresponding statistical system, which has smaller dimension than the original one. This dimensional reduction expresses the hierarchical structure in which more complicated high-dimensional systems are a superposition of lower-dimensional ones . This structure is especially visible when one uses generalization of Fourier transformation in loop space or in p-brane spaces.
In conclusion I would like to acknowledge R.Kirschner, B.Geyer and W.Janke for discussions and kind hospitality at Leipzig University. This work was supported in part by the EEC Grant no. HPMF-CT-1999-00162.
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# HIGH ENERGY PHYSICS AND QUINTESSENCEaafootnote aCERN-TH/2000-078
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# Bohmian description of a decaying quantum system
## 1 Introduction
Consider a nonrelativistic particle in a given potential field. In Bohmโs ontological interpretation of quantum mechanics, the position and velocity of the particle are both well-defined at any instant of time. The particle moves along a Bohm trajectory, irrespective of whether or not the particle is observed . The Schrรถdinger wave function accompanies and guides the particle . We do not talk about the collapse of the wave function caused by observation. A certain โsurrealisticโ aspect of the Bohm trajectory in the presence of a Welcher Weg (which way) detector has been discussed recently but we do not consider such a situation in this Letter .
The particle obeys Newtonโs equation of motion with a potential which consists of the usual given potential and the โquantum potentialโ. The latter is related to the amplitude of the wave function. The initial condition of the motion, however, can be determined only statistically. With this statistical uncertainty implemented, Bohmโs theory agrees with the traditional quantum mechanics for all observable quantities. But the picture that Bohmโs theory offers is strikingly different from the traditional one. A variety of problems have been examined on the basis of Bohmโs theory. Among them let us mention the recent work by Leavens et al. and Oriols et al. on the one-dimensional tunneling problem, which has direct relevance to what we propose to examine. The purpose of this Letter is to extend the application of Bohmโs theory to the decay problem.
We consider a model with a particle in a central potential $`V(r)`$. The potential has a repulsive barrier that supports one or more unstable bound states or resonances. The particle is initially confined inside the potential barrier and at a certain time, $`t=0`$, it begins to leak out. The model and its variants can be used to simulate a decay process through tunneling such as the nuclear $`\alpha `$ decay or the emission of an electron from an artificial atom (quantum dot) . In order to apply Bohmโs theory to such a system, one has to know the wave function of the system explicitly. Bohmโs approach has not been attempted for the decay problem so far. This is because very little has been known about the wave function of such a decaying system, particularly outside the potential barrier. A recently developed technique, however, has made it possible to solve the time-dependent Schrรถdinger equation accurately, at least for the type of model that we consider below, no matter how large $`r`$ and $`t`$ are .
In Section 2 we set up the model and we present a solution of the time-dependent Schrรถdinger equation of the model with an appropriate initial condition. In Section 3 we work out the Bohmian description of the decay process and discuss various features such as the exponential decay law and deviation from it at very small time. A summary is given in Section 4.
## 2 Model
We assume a simple model with the repulsive potential barrier,
$$V(r)=(\lambda /a)\delta (ra),$$
(1)
where $`\lambda >0`$ and $`a>0`$ are constants. A particle is initially confined inside the potential barrier and it begins to leak out at $`t=0`$. This model has been used by a number of authors to examine various features of the decaying quantum system such as the exponential decay law and deviations from it at very small time as well as at very large time . We use units such that $`\mathrm{}=1`$ and $`2m=1`$ where $`m`$ is the mass of the particle of the model. In numerical illustrations we set $`a=1`$. For the strength of the potential we take $`\lambda =6`$, which is one of the cases assumed in earlier work . The $`\lambda `$ corresponds to the $`G`$ of .
We consider only the $`S`$ state. The wave function $`\psi (r,t)`$ (actually the wave function times $`r`$) is determined by the time-dependent Schrรถdinger equation for $`t>0`$,
$$\mathrm{i}\frac{\psi (r,t)}{t}=\left[\frac{^2}{r^2}+V(r)\right]\psi (r,t),\psi (0,t)=0.$$
(2)
For the initial condition for $`\psi (r,t)`$, let us assume the normalized function
$$\psi (r,0)=\sqrt{\frac{2}{a}}\mathrm{sin}\left(\frac{\pi r}{a}\right)\theta (ar),$$
(3)
where $`\theta (x)=1`$ (0) if $`x>0`$ ($`x<0)`$.
Although the model is very simple, its time-dependent Schrรถdinger equation is nontrivial to solve . The wave function of the model in the entire space was found in an analytical form only recently . It reads as
$$\psi (r,t)=\underset{\nu }{}c_\nu [M(k_\nu ,ra,t)+N_{}(k_\nu ,ra,t)],$$
(4)
$$N_\pm (k,x,t)=\frac{\mathrm{i}\lambda }{2ka}[M(k,x,t)\pm M(k,x,t)]\theta (x),$$
(5)
where the summation is over $`\nu =\pm 1,\pm 2,\mathrm{}`$. The $`k_\nu `$โs, which are determined by solving the equation $`ka\mathrm{cot}ka+\lambda \mathrm{i}ka=0`$ for $`k`$, are the positions of the poles of the S matrix. They are all in the lower half of the complex $`k`$ plane. We designate the poles in the fourth (third) quadrant with $`\nu =1,2,\mathrm{}`$ ($`\nu =1,2,\mathrm{}`$). For $`r>a`$, $`c_\nu =2\pi \mathrm{i}a_\nu (r)\mathrm{e}^{\mathrm{i}k_\nu a}`$ for the $`a_\nu (r)`$ defined in . It is given by
$$c_\nu =\frac{2\pi \sqrt{2a}k_\nu }{(k_\nu ^2a^2\pi ^2)[(1+\lambda \mathrm{i}k_\nu a)\mathrm{cot}k_\nu a\mathrm{i}k_\nu a]}.$$
(6)
The $`c_\nu `$โs satisfy $`_\nu c_\nu /k_\nu =0`$, but $`_\nu c_\nu 0`$. The $`M(k,x,t)`$ is the Moshinsky function
$$M(k,x,t)=\frac{1}{2}\mathrm{e}^{\mathrm{i}k^2t}\mathrm{e}^{\mathrm{i}kx}\mathrm{erfc}(y),y=\mathrm{e}^{\mathrm{i}\pi /4}\left(\frac{x2kt}{2\sqrt{t}}\right),$$
(7)
where $`\mathrm{erfc}(y)=(2/\sqrt{\pi })_y^{\mathrm{}}\mathrm{e}^{u^2}du`$. In the limit of $`t0`$, $`M(k,x,t)`$ becomes discontinuous at $`x=0`$. Otherwise $`M(k,x,t)`$ is a smooth function of $`x`$. The $`\psi (r,t)`$ for $`t>0`$ is continuous at $`r=a`$ as can be seen from $`N_{}(k,0,t)=0`$.
The $`r`$-derivatives of $`\psi (r,t)`$ are given by
$`\psi ^{}(r,t)`$ $`=`$ $`\mathrm{i}{\displaystyle \underset{\nu }{}}c_\nu \{k_\nu [M(k_\nu ,ra,t)`$ (8)
$`+`$ $`N_+(k_\nu ,ra,t)]+\chi (ra,t)\},`$
$`\psi ^{\prime \prime }(r,t)={\displaystyle \underset{\nu }{}}c_\nu \{k_\nu ^2[M(k_\nu ,ra,t)+N_{}(k_\nu ,ra,t)]`$
$`+{\displaystyle \frac{ra+2k_\nu t}{2t}}\chi (ra,t)\}+{\displaystyle \frac{\lambda }{a}}\delta (ra)\psi (a,t),`$ (9)
$$\chi (x,t)=\frac{\mathrm{e}^{\mathrm{i}\pi /4}}{2\sqrt{\pi t}}\mathrm{exp}\left(\frac{\mathrm{i}x^2}{4t}\right).$$
(10)
It is not difficult to confirm that $`\psi `$ of Eq. (4) does satisfy Eq. (2). The $`\delta `$-function part of $`\psi ^{\prime \prime }`$ exactly cancels $`V\psi `$ in the Schrรถdinger equation. In deriving $`\psi `$, $`\psi ^{}`$ and $`\psi ^{\prime \prime }`$ we used $`_\nu c_\nu /k_\nu =0`$. The convergence of the $`\nu `$ summation for $`\psi `$ can be dramatically improved by adding $`\chi (ra,t)_\nu c_\nu /k_\nu `$, which is formally zero, to the right hand side of Eq. (4).
## 3 Bohmian description
We write the wave function as
$$\psi (r,t)=R(r,t)\mathrm{exp}[\mathrm{i}S(r,t)],$$
(11)
where $`R(r,t)`$ and $`S(r,t)`$ are both real. Then Eq. (2) leads to
$$\frac{S}{t}+(S^{})^2+U(r,t)=0,$$
(12)
$$U(r,t)=V(r)+Q(r,t),Q(r,t)=R^{\prime \prime }/R,$$
(13)
$$\frac{R^2}{t}+2(R^2S^{})^{}=0,$$
(14)
where $`S^{}=S/r`$ and $`R^{\prime \prime }=^2R/r^2`$. The $`Q(r,t)`$ is the quantum potential.
In Bohmโs interpretation, the particle obeys Newtonโs equation of motion with the potential $`U(r,t)=V(r)+Q(r,t)`$. Equation (12) is the classical Hamilton-Jacobi equation with potential $`U(r,t)`$. The momentum $`p=mv`$ ($`m=1/2`$) of the particle at $`(r,t)`$ is given by
$$p=mv=S^{}(r,t).$$
(15)
This $`p`$ should not be confused with the usual quantum momentum operator. Equation (15) together with an appropriate initial condition determines the particle trajectory. The motion is causal and deterministic. The underlying potential $`U(r,t)`$ is of a nonlocal and holistic nature. It depends on the wave function which in turn is related to aspects of the system at points different from $`r`$.
The motion is subject to uncertainty in the sense that the initial condition is known only statistically. The particle may start at any point $`r=r_0`$ within the potential barrier; $`0<r_0<a`$. Since $`S(r,0)=0`$ for the $`\psi (r,0)`$ of Eq. (3), the initial velocity is zero, i.e., $`v(r_0,0)=0`$. Two trajectories starting at different points at $`t=0`$ do not cross each other <sup>1</sup><sup>1</sup>1 Since we are considering only the $`S`$ state the wave function is independent of angles. The velocity associated with a Bohm trajectory has only a radial component. The particle that starts at $`r_0`$ moves always in the radial direction. In this Letter by a trajectory we mean a plot of $`r`$ versus $`t`$, not a path in the $`xyz`$ configuration space.. It is understood that the probability with which the particle starts at $`r_0`$ is proportional to $`|\psi (r_0,0)|^2`$. An observable quantity of the system is calculated by taking an average over all possible trajectories weighted according to $`|\psi (r_0,0)|^2`$. Equation (14) guarantees the conservation of the probability of the particle in the statistical ensemble of the trajectories.
With the wave function given in Section 2 we can determine $`R`$, $`S`$ and their derivatives by using $`\mathrm{Im}(\psi ^{}/\psi )=S^{}`$, $`\mathrm{Re}(\psi ^{\prime \prime }/\psi )=(R^{\prime \prime }/R)S_{}^{}{}_{}{}^{2}`$, etc. At $`r=a`$, it can be shown that $`R`$, $`S`$ and $`S^{}`$ are all continuous, while $`R^{}(a+0,t)R^{}(a0,t)=(\lambda /a)R(a,t)`$. The quantum potential $`Q`$ contains a $`\delta `$-function term, which cancels the $`\delta `$-function of $`V`$ in potential $`U`$. When $`t>0`$, $`U`$ is continuous at $`r=a`$. At $`t=0`$, we have
$$Q(r,0)=\left(\frac{\pi }{a}\right)^2,r<a.$$
(16)
Unlike for $`t>0`$, there is no exact cancellation of the $`\delta `$-function terms of $`Q`$ and $`V`$ at $`t=0`$ and hence $`U(r,0)`$ is singular at $`r=a`$. As soon as $`t`$ becomes positive, this singularity disappears. Let us define the energy $`E(r,t)`$ of the particle by
$$E(r,t)=p^2+U(r,t).$$
(17)
The particle can start at any point $`0<r_0<a`$. Since $`v(r_0,0)=0`$, the initial value of the energy is $`E(r_0,0)=(\pi /a)^2`$ which is independent of $`r_0`$. Note, however, $`E(r,t)`$ is not conserved during the motion because $`U(r,t)`$ depends on $`t`$ explicitly.
The trajectories can be labeled in terms of the starting point $`r_0`$. Instead of $`r_0`$, it is often more convenient to use
$$s(r_0)=_0^{r_0}\psi ^2(r,0)dr=\frac{r_0}{a}\frac{1}{2\pi }\mathrm{sin}\left(\frac{2\pi r_0}{a}\right),$$
(18)
which is the probability for the particle being in the region $`(0,r_0)`$ at $`t=0`$. Note that $`s(a)=1`$, $`s(a/2)=1/2`$ and $`s(0)=0`$. For the trajectory that starts at $`r_0`$ at $`t=0`$, let us denote the position function with $`r(r_0,t)`$. We determine $`r(r_0,t)`$ by solving Eq. (15). The consistency between the trajectory and the wave function requires that
$$s(r_0)=_0^{r(r_0,t)}|\psi (r,t)|^2dr,$$
(19)
at any time. This is based on the noncrossing nature of the trajectories. In fact, Eq. (19) can be used as an alternative method of determining $`r(r_0,t)`$ . We have numerically solved Eq. (15) starting with a small but finite value of $`t`$, in order to avoid the highly oscillatory behavior of $`\psi (r,t)`$ as $`t0`$. We have verified our solutions by varying the initial value of $`t`$ and also by using Eq. (19).
So far we have assumed that the trajectories start at $`t=0`$. We will assume the same in the rest of this Letter unless we state otherwise. Let us note, however, that the starting time can be chosen at will. For a trajectory that starts at position $`r`$ and time $`t>0`$ with velocity $`S^{}(r,t)`$, one can unambiguously trace its history back to $`t=0`$. We will do this later for the trajectories shown in Fig. 4.
The quantum decay process generally goes through three stages, I, II and III, which are characterized by different $`t`$ dependence of the nonescape probability $`P(t)`$ ,
$$P(t)=_0^a|\psi (r,t)|^2dr.$$
(20)
In the initial stage I we have, approximately, $`1P(t)t^2`$. We have the exponential law $`P(t)=\mathrm{e}^{\mathrm{\Gamma }t}`$ in stage II and the power law $`P(t)1/t^3`$ in final stage III <sup>2</sup><sup>2</sup>2There is a claim that $`P(t)1/t`$ in stage III . On the other hand it was argued that it is $`1/t^3`$ . For our model we can calculate $`P(t)`$ directly by using the explicit wave function of Eq. (4). We find that $`P(t)1/t^3`$ in stage III, in agreement with ..
Between the three stages are transition periods in which $`P(t)`$ exhibits an irregular, oscillatory behaviour . Stage II usually spans most of the life time of the decaying system. In our model with the chosen parameters the transition between stages I and II is around $`t=0.2`$ and that between II and III is around $`t=12`$. Toward the end of stage II, $`P(t)`$ becomes as small as $`10^8`$. By then the system has almost completely decayed. We present results in five figures.
In Fig. 1 we show trajectories with equally spaced $`s`$, with interval of $`\mathrm{\Delta }s=1/N`$, $`N=30`$. These trajectories have equal statistical weights. In other words, each of the trajectories occurs with the same probability. The trajectory density at $`(r,t)`$ is proportional to the probability density $`|\psi (r,t)|^2`$. The trajectories do not cross each other. A trajectory that starts near the barrier ($`a=1`$) escapes earlier. Outside the barrier the trajectories become nearly straight. The slopes of the four or five trajectories that leave the barrier the earliest are somewhat steeper than the others. If we assume that $`E(r,t)`$ of Eq. (17) is conserved and that $`U(r,t)`$ is negligible for $`r>a`$, we obtain $`E(r,t)=Q(r,0)=(\pi /a)^2`$ for $`r>a`$. This leads to $`p=\pi /a`$ and $`v=2\pi `$ ($`2m=1`$, $`a=1`$) outside the barrier. This is approximately the case as can be seen from the slopes of the trajectories of Fig. 1.
Let us label the trajectories with $`n=1,2,\mathrm{},29`$, starting with the one that escapes first. In terms of the starting point $`r_0`$, trajectory 1 is the closest to the barrier. For each of the trajectories we define the escape time $`t_n`$ as the time at which the particle crosses $`a=1`$ outward. (Such escape times or โexit timesโ have been formally discussed by Daumer et al. in the context of the scattering problem in three dimensions.) At $`t=t_n`$ out of the $`N`$ trajectories, $`(Nn)`$ trajectories are still within the boundary. This means that the nonescape probability $`P(t)`$ is given by
$$P(t)=\frac{Nn}{N},t_n<t<t_{n+1}.$$
(21)
Ideally $`N`$ should be taken as an infinitely large number. If the exponential decay law holds exactly, that is, if $`P(t)=\mathrm{e}^{\mathrm{\Gamma }t}`$, we obtain
$$t_n=\frac{1}{\mathrm{\Gamma }}\mathrm{ln}\left(\frac{N}{Nn}\right),$$
(22)
where $`\mathrm{\Gamma }`$ is related to the half-life $`\tau _{1/2}`$ through $`\tau _{1/2}=\mathrm{ln}2/\mathrm{\Gamma }`$. Figure 2 shows $`t_n`$ versus $`\mathrm{ln}[N/(Nn)]`$. The dots correspond to the trajectories shown in Fig. 1.
By fitting the numerically calculated nonescape probability of the same model as ours with $`P(t)=\mathrm{e}^{\mathrm{\Gamma }t}`$, Winter obtained $`1/\mathrm{\Gamma }=0.644`$ which leads to $`\tau _{1/2}=0.446`$ . In Fig. 2 the dashed line shows the $`n`$ dependence of $`t_n`$ given by Eq. (22) with Winterโs $`1/\mathrm{\Gamma }`$. Except for the first several ones, the dots follow the exponential curve very well. The 15-th trajectory starts at $`r_0=0.5`$ and $`s=0.5`$. It crosses the barrier at $`t_{15}=0.468`$, which is the half-life $`\tau _{1/2}`$. This is slightly larger than the value 0.446 which is based on Winterโs estimate. (In Section 3 of below Eq. (17) the half-life for $`G=6`$ was misquoted as $`\tau _{1/2}=1.08`$. The correct value is 0.446.) Note that the decay process begins at a rate slower than predicted by the exponential law. This explains why the $`\tau _{1/2}`$ estimated by the trajectory of $`s=1/2`$ is greater than the one based on the exponential law.
For $`t<0.2`$ the escape time does not follow the exponential law very well. It is in fact better fitted with $`t_n\sqrt{n}`$. In Fig. 3 we plot the escape time against $`\sqrt{n}`$. In order to see the details we have increased the number of trajectories, which are again equally spaced with respect to $`s`$ but with a smaller interval $`\mathrm{\Delta }s=1/N`$ with $`N=100`$. They are labelled with $`n=1,2,\mathrm{},99`$ but the escape time is shown only for the first 25 trajectories. The first three dots are almost exactly on a straight line which passes through the origin. This means that $`1P(t)t^2`$ and hence $`\mathrm{d}P(t)/\mathrm{d}t=0`$ at $`t=0`$. In the beginning of the decay process, $`P(t)`$ decreases more slowly than the exponential law predicts. This is analogous to the โstandby mechanismโ which Elberfeld and Kleber discussed in their analysis of time-dependent tunneling of a semi-infinite wave train through a thin barrier. This deviation from the exponential law for very small $`t`$ is a general feature of the quantum decay process which is related to the possibility of the quantum Zeno effect as discussed in . For experimental evidence, see .
Let us take the model as a simulation of an $`\alpha `$-decaying nucleus and examine how the nuclear charge $`Z(t)`$ (in units of $`e>0`$) varies as a function of $`t`$ when the $`\alpha `$ particle of charge 2 is emitted. In the traditional theory the charge number $`Z(t)`$ of the nucleus is given by
$$Z(t)=Z(0)2[1P(t)].$$
(23)
The nonescape probability is well approximated by $`P(t)=e^{\mathrm{\Gamma }t}`$ for most of the time. The $`Z(t)`$ changes from $`Z(0)`$ to $`Z(0)2`$ gradually. This is because the wave function $`\psi (r,t)`$ of the $`\alpha `$ particle leaks out gradually.
In the Bohmian description, if one follows the $`n`$-th trajectory, the nuclear charge changes from $`Z(0)`$ to $`Z(0)2`$ suddenly at time $`t_n`$ when the $`\alpha `$ particle leaves the nucleus and hence
$$Z_n(t)=Z(0)2\theta (tt_n).$$
(24)
Here suffix $`n`$ refers to the $`n`$-th trajectory. In order to obtain the nuclear charge that can be compared with that of the traditional theory, we have to consider the ensemble of all trajectories each with a weight $`|\psi (r_0,0)|^2`$. This weight in the present case is $`1/N`$ for each trajectory. At time $`t`$ such that $`t_n<t<t_{n+1}`$, the particles of trajectories of 1, 2, $`\mathrm{},n`$ have escaped. We thus obtain
$$Z(t)=Z(0)\frac{2n}{N}=Z(0)2[1P(t_n)],$$
(25)
where we have used Eq. (21). In the limit of $`N\mathrm{}`$, this $`Z(t)`$ converges to the $`Z(t)`$ of Eq. (23) of the traditional theory. This illustrates how a quantity that appears in the Bohmian description can be related to its counterpart of the traditional theory.
Consider a gedanken experiment in which one tries to detect an $`\alpha `$ particle that is emitted from a source consisting of a single $`\alpha `$ emitting nucleus. The event in which one detects an $`\alpha `$ particle corresponds to one of the Bohm trajectories. If one repeats this experiment many times one experiences events, each of which is described by one of the Bohm trajectories. Here it is understood that the source is prepared every time in an identical manner. In contrast to this, the Schrรถdinger wave function does not describe any of the individual events, rather it only describes an ensemble of a large number of such events. Instead of repeating the experiment on one system, we can think of experiments on many independent systems that are all identically prepared. In this sense, the word โevent(s)โ can be replaced with โsystem(s)โ.
Figure 4 shows 21 trajectories such that $`r(t=10)`$ ranges from 0.2 to 0.6 with the interval $`\mathrm{\Delta }r=0.02`$. We have obtained these trajectories by integrating Eq. (15) starting at $`t=10`$. For the starting points, we have chosen to keep $`\mathrm{\Delta }r`$ constant rather than $`\mathrm{\Delta }s`$. This choice is only a matter of convenience or simplicity of the calculation involved. Because of this choice, unlike in Fig. 1, the trajectory density in this figure is not proportional to the probability density. For example, in the figure the statistical weight is larger for a trajectory with a larger value of $`r(t=10)`$. The range in terms of probability $`s(r_0)`$ is from $`6.44\times 10^9`$ to $`1.047\times 10^7`$. As we stated before we can easily trace the history of the trajectories back to $`t=0`$. The values of $`r_0`$ at $`t=0`$ of the trajectories range from $`9.93\times 10^4`$ to $`2.52\times 10^3`$. These trajectories are starting almost from the origin. This is why they remain inside for a very long time.
The trajectories go back and forth across the barrier. As Winter pointed out many years ago, the current density $`j(r,t)`$ at the barrier fluctuates in this time interval <sup>3</sup><sup>3</sup>3There are other space-time regions in which similar fluctuations (with larger amplitudes) of the current density occur. See Fig. 5 of .. At times it becomes negative, i.e., inward. Figure 4 visualizes this feature. In such a situation we redefine the escape time as the time when the particle finally leaves the barrier. The exponential law does not hold in this time region any longer. Note also that, after leaking out through the potential barrier, the trajectories tend to remain close to the barrier. This situation is very different from that of Fig. 1. The time interval shown in Fig. 4 is the transition period between stages II and III. In the latter, the nonescape probability decreases like $`1/t^3`$.
Figure 5 shows the potential $`U(r,t)`$ which is equal to the quantum potential $`Q(r,t)`$ with its $`\delta `$-function part removed. There is no potential barrier in $`U(r,t)`$ and hence there is no tunneling phenomenon. As $`r`$ increases across the barrier at $`r=a=1`$, $`U`$ sharply drops but $`U`$ is continuous across the barrier (except at $`t=0`$). The behavior of $`U`$ is complicated for very small $`t`$ and also for very small $`r`$. The figure does not show the part of $`t<0.05`$ and $`r<0.001`$. Close scrutiny reveals that the trajectories rapidly fluctuate when $`t`$ and hence $`r(t)`$ are very small. Although we do not show it, the behavior of $`U(r,t)`$ is also complicated in the space-time region that corresponds to Fig. 4. For $`ra`$, $`U`$ becomes negligible.
The results shown above are all for the case of $`\lambda =6`$. We have also examined the case of larger values of $`\lambda `$. For example, when $`\lambda =100`$, the deviations from the exponential law are very small. Let us add that, if we are to simulate $`\alpha `$ decay processes, we have to assume much larger values of $`\lambda `$, for example, of the order of 10<sup>8</sup> for <sup>212</sup>Po. See Section 5 of .
## 4 Summary
For the model defined by Eqs. (1) and (3) we examined the decay process from Bohmโs point of view. We obtained Bohm trajectories with which we can interpret various features of the decay process. We see deviations from the exponential law at very small time and also at very large time. The decay process is slower in the beginning than the exponential law predicts. In the time interval of $`t=`$ 10 to 14, we showed that the trajectories go back and forth across the barrier. This corresponds to the current density fluctuations that Winter found. Beyond that time region, the exponential law is replaced with a power law. One can verify this by showing that for large times and very small $`r_0`$, $`t_n(Nn)^{1/3}`$, confirming that the trajectories in this region of very large time are consistent with $`P(t)1/t^3`$.
The results that we obtained in the Bohmian picture are complementary but not contradictory to the traditional quantum mechanics. Let us, however, emphasize the following point. In contrast to the Schrรถdinger wave function which describes an ensemble of a large number of events, each of the Bohm trajectories represents an individual event. Such information on individual events is masked when the uncertainty regarding the starting points of the trajectories is incorporated. In a situation in which one has to deal with an individual event, however, the Bohmian approach may lead to new insights.
The model that we have considered is one of the simplest models for the decay process. The method that we have used can be applied to other potential models.
This work was supported by the Natural Sciences and Engineering Research Council of Canada.
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# Low Frequency Gravitational Waves from White Dwarf MACHO Binaries
## 1 Introduction
The MACHO project has detected 13-17 gravitational microlensing events in the direction of the Large Magellanic Cloud to date (Alcock et al., 2000), while the EROS collaboration has detected two (Lasserre et al., 1999). One result of the analysis of the observations (Alcock et al., 2000) is that, independent of the halo model assumed, there are of order $`2\times 10^{11}`$ MACHOs of mean mass $`0.5M_{}`$ in the Galactic halo.<sup>1</sup><sup>1</sup>1It is noteworthy that this number is similar in magnitude to the number of stars in the Galaxy.
While the microlensing observations have confirmed the existence of MACHOs (Massive Astrophysical Compact Halo Objects) in the halo, the actual nature of these objects is still a matter of much debate. The mass suggests main sequence red dwarf stars, but these appear to be ruled out observationally (Bahcall et al., 1994; Graff & Freese, 1996a; Graff & Freese, 1996b). Until recently, white dwarf stars also seemed to be a highly unlikely MACHO candidate, since observations of the Galactic halo (see, e.g., Flynn et al., 1996) did not find a population of old, red, white dwarf stars. Since conventional baryonic stars seemed to be ruled out, highly speculative candidate objects have been proposed as MACHO models, such as primordial black holes and boson stars.
Recently, however, new models of white dwarf cooling processes indicate that old, cool, white dwarf atmospheres form molecular hydrogen, which can strongly absorb red light. This implies that aging white dwarfs will be blue objects, rather than red as was previously believed (Saumon & Jacobson, 1999; Hansen 1999). In light of these new cooling models, predictions have been made for the number of halo white dwarfs that should be seen in deep field surveys (Richer, 1999), and new analyses have detected candidate halo white dwarfs with blue colors in the Hubble Deep Field (Ibata et al., 1999; Mรฉndez & Minniti, 2000). In addition, local blue white dwarfs with large proper motions, indicating that they are members of the halo population, have been tentatively identified (Ibata et al., 2000).
In this Letter we assume that the halo MACHOs are white dwarf (WD) stars, and estimate the low frequency ($`10^5\mathrm{Hz}<f<10^1\mathrm{Hz}`$) gravitational wave (GW) background that would be produced from a halo population of white dwarf binaries. Such a background could be an important source for space-based laser interferometer gravitational wave detectors such as the proposed LISA mission (Bender et al., 1998). Gravitational waves from halo WD binaries could be a very interesting signal (if one is interested in characterizing the binary population of MACHOs) or a serious confusion noise source (if one is concerned that this background is masking signals from other sources of interest, such as a cosmological background of GW). Estimates of the gravitational wave signal from a halo composed of primordial black holes has shown that the signal from binaries in the halo could dominate the output of an interferometer such as LISA (Hiscock, 1998; Ioka et al., 1999).
In the absence of any observational evidence concerning the properties a halo population of binary WDs, we make the assumption that white dwarf binary properties in the halo mimic those of Galactic disk WD binaries (Hils, Bender & Webbink, 1990; Bender & Hils, 1997). This assumption is almost certainly incorrect; the halo WD population is generally felt to be much older than the disk population, and (based on the microlensing events) probably has a mass distribution that differs from the disk WD population. However, using the disk as a model is the best one can currently do.
## 2 Disk binaries as GW sources
The GW background generated by disk binaries (both Galactic and extragalactic) has been thoroughly studied by Hils, Bender, and Webbink (Hils, Bender & Webbink, 1990; Bender & Hils, 1997). They have made careful estimates of the GW signal due to Galactic disk WD binaries and also that due to extragalactic binaries. In recent work, they have combined these signal amplitudes with the planned properties of LISA to generate a simulation of LISAโs response to the combined Galactic and extragalactic signals. The key factor in this analysis is the width of a frequency bin in the LISA data analysis for periodic sources. With a one-year integration time, each frequency bin will have a width $`\mathrm{\Delta }f=(1\mathrm{yr})^13\times 10^8\mathrm{Hz}`$. At frequencies beginning at around $`1\times 10^3`$ Hz and higher, the number of Galactic binaries per bin will begin to drop to order unity. At this point, the properties of individual Galactic binaries can be determined and their signal removed from the record, so that the weaker combined signal of extragalactic binaries will begin to be observable in the open bins. Solving for a particular Galactic binary and removing it from the data record will typically require three bins of information (Hellings, 1996). The effective spectral amplitude observed by LISA, $`h_f`$, after taking the finite bin width into account, is given by (Bender & Hils, 1997)
$$h_f=\left(\frac{(h_f^e)^2\left[(h_f^G)^2+(h_f^e)^2(1e^{3r})\right]}{e^{3r}(h_f^G)^2+(h_f^e)^2(1e^{3r})}\right)^{1/2},$$
(1)
where $`r`$ is the number of Galactic binaries per frequency bin, $`h_f^G`$ is the spectral amplitude of the Galactic binary background, and $`h_f^e`$ is the spectral amplitude of the extragalactic binaries.
Using the relation between the spectral amplitude and the number of binaries per unit frequency $`dN/df`$, the GW luminosity of a binary $`L(f)`$, and the average inverse distance squared $`d^2`$,
$$h_f(f)=\frac{2}{\pi f}\left[\frac{1}{d^2}L(f)\frac{dN}{df}\right]^{1/2},$$
(2)
it is possible to extract an estimate of $`dN/df`$ from the Bender-Hils results, along with an approximation to the spectral amplitudes of the backgrounds due to Galactic disk binaries and extragalactic binaries. We find that the Bender-Hils results are well approximated by
$$\frac{dN}{df}4.47\times 10^3f^{11/3},$$
(3)
$$\mathrm{log}_{10}\left(h_f^G\right)=\left(\frac{7}{6}\right)\mathrm{log}_{10}(h_f)21.8,$$
(4)
and
$$\mathrm{log}_{10}\left(h_f^e\right)=\left(\frac{7}{6}\right)\mathrm{log}_{10}(h_f)23.0$$
(5)
where the frequency $`f`$ is measured in Hz. The frequency dependence in these equations is characteristic of a population of circular binaries that is evolving solely due to gravitational radiation reaction.
## 3 Rescaling the Disk to the Halo
A simple estimate of the gravitational wave background expected from binary WDs in the halo can be obtained by assuming that the WD binary population of the halo is similar in nature to that of the disk. One can estimate the halo GW background by rescaling the disk binary WD population, based on three factors:
1. The ratio of the total number of halo WD MACHOs to the total number of Galactic disk WDs, $`N_{\mathrm{halo}}/N_{\mathrm{disk}}`$.
2. The ratio of the average inverse distance squared of a halo MACHO to the average inverse distance squared of a disk WD, $`d^2_{\mathrm{halo}}/d^2_{\mathrm{disk}}`$.
3. The ratio of the fraction of white dwarfs in binaries in the halo to the fraction of white dwarfs in binaries in the disk, $`\alpha `$.
The number of WDs in the disk is computed by integrating over the standard cylindrical exponential model,
$$\rho =\rho _0\mathrm{exp}\left[\frac{r}{r_0}\right]\mathrm{exp}\left[\frac{|z|}{z_0}\right],$$
(6)
where $`r_0=3.5`$ kpc and $`z_o=90`$ pc are the exponential scale heights for the WD population in the disk (Hils, Bender & Webbink, 1990), and $`\rho _0=4.72\times 10^2\mathrm{pc}^3`$ is the number density of white dwarfs at the center of the galaxy (computed from the local density of white dwarfs in the solar neighborhood, $`\rho _{}=4.16\times 10^3\mathrm{pc}^3`$ given in Knox, Hawkins & Hambly, 1999). Integrating the disk WD density using the distribution in Eq. (6) yields
$$N_{\mathrm{disk}}=6.5\times 10^8$$
(7)
for the disk population.
We assume that the number of MACHOs in a $`50`$ kpc halo (to the LMC), is $`N_{\mathrm{halo}}^{50\mathrm{kpc}}=2\times 10^{11}`$, the halo-model-independent result obtained by the MACHO collaboration (Alcock et al., 2000). For a larger halo, extending some $`300`$ kpc (halfway to M31), this number can be scaled by assuming that the spatial distribution of white dwarf MACHOs follows the standard spherical flat rotational halo model given by
$$\rho =\widehat{\rho }\frac{R^2+a^2}{r^2+a^2},$$
(8)
where $`\widehat{\rho }`$ is the local density of dark matter, $`r`$ is the Galactocentric radius, $`R=8.5\mathrm{kpc}`$ is the Galactocentric radius of the Sun, and $`a=5.0\mathrm{kpc}`$ is the halo core radius. Integrating Eq. (8) out to $`50`$ kpc and setting the number of MACHOs equal to $`2\times 10^{11}`$, one obtains $`\widehat{\rho }=0.0094\mathrm{M}_{}\mathrm{pc}^3`$. Using this value and integrating Eq. (8) out to 300 kpc gives
$$N_{\mathrm{halo}}^{300\mathrm{kpc}}=1.3\times 10^{12}.$$
(9)
The average inverse distance squared between sources and the Sun, $`d^2`$, may be found by integrating over the source distributions given in Eqs. (8) and (6). Expressing the result as a distance, one finds that for the disk,
$$d^2^{1/2}=4.85\mathrm{kpc},$$
(10)
while
$$d^2^{1/2}=15.7\mathrm{kpc}$$
(11)
for a $`50`$ kpc halo, and
$$d^2^{1/2}=39.5\mathrm{kpc}$$
(12)
for a $`300`$ kpc halo.
Nothing is presently known about the ratio of the fraction of white dwarfs in binaries in the halo to the fraction of white dwarfs in binaries in the disk, $`\alpha `$, so we leave this as a free parameter in our model, and examine the consequences of different values for $`\alpha `$.
In rescaling the disk background GW spectrum, there are two separate effects associated with the potentially larger number of WD binaries in the halo. First, the larger number tends to raise the overall level of the halo background relative to that of the disk. Second, since there are more halo binaries per unit frequency interval, this pushes the point in the spectrum where one first encounters open frequency bins (one or fewer Galactic halo binaries per bin), to higher frequencies than in the disk. Scaling the number of Galactic disk binaries by multiplying by the ratio of the total number of WDs in the halo to the number in the disk yields
$$\mathrm{log}_{10}\left(\frac{dN}{df}\right)_{\mathrm{halo}}=\left(\frac{11}{3}\right)\mathrm{log}_{10}(f)2.35+\mathrm{log}_{10}\left(\frac{\alpha N_{\mathrm{halo}}}{N_{\mathrm{disk}}}\right).$$
(13)
The overall level of the GW backgrounds from Galactic halo WD binaries and extragalactic halo binaries will scale according to
$$h_f^{\mathrm{halo}}=K(\alpha )h_f^{\mathrm{disk}},$$
(14)
where the scaling factor $`K(\alpha )`$ is defined as
$$K(\alpha )=\left[\alpha \frac{N_{\mathrm{halo}}}{N_{\mathrm{disk}}}\frac{d^2_{\mathrm{halo}}}{d^2_{\mathrm{disk}}}\right]^{1/2}.$$
(15)
If one assumes that the halo population of white dwarfs precisely mimics the disk population, then the same fraction of WDs will be in binaries in the halo as in the disk, and $`\alpha =1`$. In this case the scaling factor is given by
$$K(\alpha =1)=\{\begin{array}{cc}5.42\hfill & \text{for a 50 kpc halo}\hfill \\ 5.49\hfill & \text{for a 300 kpc halo}\hfill \end{array}.$$
(16)
An estimate of the response of LISA to a background of GW from halo WD binaries can now be obtained by rescaling the Galactic and extragalactic disk spectral amplitudes \[Eqs. (4) and (5)\] using Eq. (14) with Eq. (13) in Eq. (1). The resulting signal estimate for LISA is illustrated in Figure (1), along with the Bender-Hils estimate for disk binaries for comparison. The signal from the halo WD binaries is substantially stronger than that from the disk binaries. At lower frequencies, the predicted backgrounds for a 50 kpc and 300 kpc halo are indistinguishable in the figure, owing to the similarity in the values of $`K`$ in Eq.(16). The larger number of halo binaries compared to the disk fills the frequency bins to a substantially higher frequency before one encounters open bins, where weaker signals such as the extragalactic background may be observed. If only the disk background is present, then potential extragalactic sources can be detected above a critical frequency of about $`2\times 10^3`$ Hz, where frequency bins cease to be cluttered with many Galactic binaries. With a 50 kpc halo, the greater number of Galactic binaries increases this critical frequency to about $`1\times 10^2`$ Hz, while for a larger 300 kpc halo, the critical frequency is further increased to about $`2\times 10^2`$ Hz. In the latter case, the frequency at which bins begin to open up and allow weaker extragalactic signals to be detected is roughly equal to where LISAโs instrumental noise curve begins to rise.
Of course, it is presently unknown whether the fraction of WDs in binaries in the halo is comparable to that in the disk. A priori, it could be larger ($`\alpha >1`$) or smaller ($`\alpha <1`$). One question which can be posed within the present simple model is to ask what value of $`\alpha `$ will result in the halo GW signal being similar in magnitude and shape to that of the disk. This determines a minimum value for $`\alpha `$, above which the halo signal will dominate over that of the disk binaries. Reducing $`\alpha `$ in Eqs.(13) and 14) one finds that the signal from a $`50`$ kpc halo will be dominant if $`\alpha >10^2`$, while the signal strength from a $`300`$ kpc halo would exceed that of the disk if $`\alpha >5\times 10^3`$. This implies that even if the fraction of halo WDs in binaries is as low as 1% of the fraction of disk WDs in binaries, the WD MACHO binaries will be โbrightโ enough to stand out from the expected signal of the disk binaries. While the numbers here are highly sensitive to the specific details of the halo binary MACHO population, they illustrate that the halo will be a significant source of a low frequency confusion background of gravitational waves unless the fraction of MACHO WDs in binaries is orders of magnitude lower than in the disk.
## 4 Discussion
While this simple scaled model is certainly not an accurate representation of the Galactic halo binaries, it does illustrate that the GW background from a halo population of white dwarf binaries could easily dominate the signal in a space-based interferometer such as LISA. Further, this result, together with other studies that have considered the GW signal from MACHOs if they were identified with primordial black holes (Nakamura et al. 1997; Hiscock, 1998; Ioka et al., 1999; Ioka et al., 2000), demonstrates that whatever the nature of MACHOs, if even a small fraction of them are in binary systems, then they will create a strong confusion noise background which could saturate the frequency range in which LISA is most sensitive ($`2\times 10^3\mathrm{Hz}<f<2\times 10^2\mathrm{Hz}`$).
Some may consider such a prospect discouraging, as the combined signal from abundant halo binaries could mask other weak signals and make them undetectable. This has previously been a serious concern with respect to the disk binariesโhence the name โconfusion noiseโ for a signal that actually describes the short-period binary population of the Galaxy. There has been hope that the confusion noise from Galactic disk binaries could be extracted from the general stochastic (e.g., cosmological) background by utilizing the anisotropic nature of the disk signal (Giampieri & Polnarev, 1997). If there is a substantial signal associated with halo binaries, however, then this scheme will not work. The solar position is sufficiently near the center of a spherical Galactic halo that it would seem difficult or impossible to be able to subtract out the halo confusion noise signal based on its very small anisotropy.
On the other hand, even our simple analysis shows that the GW signal from the halo binaries could be a powerful tool to determine the properties of the MACHO binary population, as well as general properties of the halo itself, such as its size.
The authors wish to thank P. Bender and R. Hellings for helpful discussions. This work was supported in part by NSF Grant No. PHY-9734834 and NASA Cooperative Agreement NCC5-410.
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# The redshift space power spectrum in the halo model
## 1 Introduction
The power spectrum of the mass fluctuations in the universe is one of the most fundamental quantities in large-scale structure. It is robust, but sensitive to several cosmological parameters such as the Hubble constant, the matter density and of course the primordial power spectrum (usually parameterized by an amplitude and a slope). While the theory behind the power spectrum in the linear regime is quite straightforward, analytically handling clustering in the non-linear regime has proven quite difficult. Recently several authors have developed a new way of looking at the non-linear power spectrum which imagines all the mass in the universe lies in a halo of some mass . They postulate that on large scales the halos cluster according to linear theory while on small scales the power is dominated by the halo profiles . Specifically most pairs of dark matter particles with small interparticle separations lie within the same halo, and thus their correlations can be predicted by the halo profile. While this model requires many ingredients to be fixed by numerical experiments (typically N-body simulations) it provides a useful structure for thinking about gravitational clustering which gives insights into several outstanding problems .
The work to date has all focussed on the clustering of the matter or galaxies in โreal spaceโ, whereas many if not most observations of clustering take place in โredshift spaceโ. It is the purpose of this work to show that redshift space distortions can be handled naturally in the halo picture and that doing so provides insight into some well studied phenomena.
## 2 The halo model
The model for non-linear clustering is based on the Press-Schechter theory, in which all of the mass in the universe resides in a virialized halo of a certain mass. The extensions to this theory introduced recently allow one to calculate the power spectra and cross correlations between both the mass and the galaxies, given a suitable prescription for how galaxies populate dark matter halos. Since at present such prescriptions are somewhat ad hoc we shall concentrate here on the mass power spectrum, though there is no obstacle in principle to extending the method to galaxies.
In this model the power spectrum, $`P(k)`$, is the sum of two pieces. The first is that due to a system of (smooth) halos of profile $`y(k)`$ laid down with inter-halo correlations assumed to be a biased sampling of $`P_{\mathrm{lin}}(k)`$. Since the real space convolution is simply a Fourier space multiplication this contribution is
$$P^{2\mathrm{halo}}(k)=P_{\mathrm{lin}}(k)\left[f(\nu )๐\nu b(\nu )y(k;M)\right]^2$$
(1)
where $`b(\nu )`$ is the (linear) bias of a halo of mass $`M(\nu )`$ and $`f(\nu )`$ is the multiplicity function. The peak height $`\nu `$ is related to the mass of the halo through
$$\nu \left(\frac{\delta _c}{\sigma (M)}\right)^2$$
(2)
where $`\delta _c=1.69`$ and $`\sigma (M)`$ is the rms fluctuation in the matter density smoothed with a top-hat filter on a scale $`R^3=3M/4\pi \overline{\rho }`$. Both $`b`$ and $`f`$ come from fits to N-body simulations. We use
$$b(\nu )=1+\frac{\nu 1}{\delta _c}+\frac{2p}{\delta _c(1+\nu ^p)}$$
(3)
and
$$\nu f(\nu )=A(1+\nu ^p)\nu ^{1/2}e^{\nu ^{}/2}$$
(4)
where $`p=0.3`$ and $`\nu ^{}=0.707\nu `$. The normalization constant $`A`$ is fixed by the requirement that all of the mass lie in a given halo
$$f(\nu )๐\nu =1.$$
(5)
We assume that the halos all have spherical profiles depending only on the mass. We neglect any substructure or halos-within-halos as this seems to be unimportant for the power spectrum. We take the โNFWโ form
$$\rho (r)=\frac{\rho _0}{x(1+x)^2},$$
(6)
where $`x=r/r_s`$ is a scaled radius, though our results are not particularly sensitive to this choice . The mass of this halo is defined to be the virial mass $`M=(4\pi /3)\delta _{\mathrm{vir}}\overline{\rho }r_{\mathrm{vir}}^3`$ where we take $`\delta _{\mathrm{vir}}=200`$ and define the virial radius, $`r_{\mathrm{vir}}`$, as the radius within which the mean density enclosed is $`\delta _{\mathrm{vir}}`$ times the background density. The concentration parameter $`c=r_{\mathrm{vir}}/r_s`$. The profile can then be described by its virial mass and concentration. N-body simulations suggest the two are related and we follow in using
$$c(M)=10\left(\frac{M}{M_{}}\right)^{0.2}$$
(7)
where $`\sigma (M_{})=1`$.
Finally then we need $`y(k)`$ which is the Fourier transform of Eq. (6) normalized to unit mass. The Fourier transform can be done analytically
$$\begin{array}{cc}\stackrel{~}{\rho }(k)=4\pi \rho _0r_s^3[\hfill & \mathrm{cos}z\left\{\mathrm{Ci}([1+c]z)\mathrm{Ci}(z)\right\}+\hfill \\ & \mathrm{sin}z\{\mathrm{Si}([1+c]z)\mathrm{Si}(z)\}\frac{\mathrm{sin}cz}{(1+c)z}]\hfill \end{array}$$
(8)
where $`zkr_s`$ and the total mass is
$$M=4\pi \rho _0r_s^3\left[\mathrm{log}(1+c)\frac{c}{1+c}\right]$$
(9)
Eq. (1) is the dominant contribution to the power spectrum on large scales. On small scales we are dominated by pairs lying within a single halo
$$P^{1\mathrm{halo}}(k)=\frac{1}{(2\pi )^3}f(\nu )๐\nu \frac{M(\nu )}{\overline{\rho }}|y(k)|^2$$
(10)
We show an example of how well this formalism predicts the matter power spectrum at $`z=0`$ in Fig. 1. We have chosen a particular $`\mathrm{\Lambda }`$CDM model with $`\mathrm{\Omega }_\mathrm{m}=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, $`h=0.7`$, $`\mathrm{\Omega }_\mathrm{B}h^2=0.02`$ and $`n=1`$. The model has been normalized to the COBE 4-year data using the method of Bunn & White . The agreement in the linear regime is good, as is to be expected. On smaller scales the model is a surprisingly good fit to the N-body data given the simplicity of the assumptions. The slight shortfall in power at higher $`k`$ can be remedied by modifying the prescription somewhat, but we will here stick to the parameters outlined in โ this model is valuable more for its pedagogical value than as a substitute for direct calculations.
## 3 Redshift space distortions
The model described above naturally lends itself to a treatment of redshift space distortions. There are two effects which come in when moving to redshift space. The first is a boost of power on large scales due to streaming of matter into overdense regions. The second is a reduction of power on small scales due to virial motions within an object. In the halo model the large-scale and small-scale effects can be separated out in a simple fashion.
Kaiser first showed that on large scales one expects an enhancement of the power spectrum in redshift space. In linear theory a density perturbation $`\delta _k`$ generates a velocity perturbation $`\dot{\delta }=ikv`$ with $`\stackrel{}{v}`$ parallel to $`\stackrel{}{k}`$. Using density conservation to linear order and making the distant observer approximation ($`kr1`$, we shall work throughout in the plane-parallel limit, see for a large-angle formalism) the redshift space density contrast can be written as
$$\delta _{\mathrm{redshift}}=\delta _{\mathrm{real}}\left(1+f\mu ^2\right)$$
(11)
where $`f(\mathrm{\Omega })d\mathrm{log}\delta /d\mathrm{log}a\mathrm{\Omega }^{0.6}`$, $`a`$ is the scale-factor and $`\mu =\widehat{r}\widehat{k}`$. Thus on large-scales the power spectrum is increased by
$$\frac{1}{2}_1^{+1}๐\mu \left(1+f\mu ^2\right)^2=1+\frac{2}{3}f+\frac{1}{5}f^2$$
(12)
which is approximately $`1.37`$ in our model.
On small scales virial motions within collapsed objects reduce power in redshift space. If we assume that our halos are isotropic, virialized and isothermal with $`1D`$ velocity dispersion $`\sigma `$ then the peculiar motions within the halo add a Gaussian noise to the redshift space radial coordinate. Once again the real space convolution becomes a Fourier space multiplication and the inferred density contrast is
$$\delta _{\mathrm{redshift}}=\delta _{\mathrm{real}}e^{(k\sigma \mu )^2/2}$$
(13)
Integrating this over $`\mu `$ gives a suppression
$$_1(y=k\sigma )=\sqrt{\frac{\pi }{2}}\frac{\mathrm{erf}(y/\sqrt{2})}{y}$$
(14)
As pointed out by Peacock & Dodds however, the โfullโ effect of redshift space distortions includes both the enhancement and the suppression of power, and one must include the $`\mu `$ dependence of both factors before doing the integral. Including both terms the redshift space distortion becomes, upon integrating over $`\mu `$,
$`_2(y=k\sigma )`$ $`=`$ $`{\displaystyle \frac{\sqrt{\pi }}{8}}{\displaystyle \frac{\mathrm{erf}(y)}{y^5}}\left[3f^2+4fy^2+4y^4\right]`$ (15)
$``$ $`{\displaystyle \frac{e^{y^2}}{4y^4}}\left[f^2(3+2y^2)+4fy^2\right]`$
Thus to predict the redshift space power spectrum in the halo model we modify Eqs. (1, 10) to
$`P(k)`$ $`=`$ $`(1+{\displaystyle \frac{2}{3}}f+{\displaystyle \frac{1}{5}}f^2)P_{\mathrm{lin}}(k)\times `$ (16)
$`\left[{\displaystyle f(\nu )๐\nu b(\nu )_1(k\sigma )y(k;M)}\right]^2`$
$`+`$ $`{\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle f(\nu )๐\nu \frac{M(\nu )}{\overline{\rho }}_2(k\sigma )|y(k)|^2}`$
In the first term we have broken out the enhancement due to halo motions, treated in linear theory, and the $`_1`$ term from the virial motion. Technically we should do the integral over $`\nu `$ including the Gaussian from Eq. (13) first and then integrate over $`\mu `$. We can see in Fig. 2 however that this numerically simpler approximation works quite well. Note that in contrast to the real-space power spectrum the clustering term remains significant to larger $`k`$.
We assume that the halos are isothermal. From the mass within the virial radius the the 1D velocity dispersion of a halo of mass $`M`$ is
$`\sigma ^2`$ $`=`$ $`GM/2r_{\mathrm{vir}}`$ (17)
$`=`$ $`G\left({\displaystyle \frac{\pi }{6}}M^2\overline{\rho }\delta _{\mathrm{vir}}\right)^{1/3}`$ (18)
We find that our results are almost entirely unchanged if we estimate $`\sigma `$ from the circular velocity interior to $`r_s`$ instead of $`r_{\mathrm{vir}}`$.
Putting the pieces together gives the lower solid line in Fig. 1. Notice that the model provides an adequate description of the redshift space power spectrum. The clustering term remains significant to smaller scales, which may explain why perturbation theory results seem to work better in redshift space than in real space and why the redshift space power spectrum approximates the linear theory prediction over such a wide range of scales. At small-$`k`$ this model makes the same predictions as linear theory with a constant โeffectiveโ bias
$$bf(\nu )๐\nu b(\nu )$$
(19)
The high-$`k`$ behavior of the model qualitatively reproduces the suppression seen in the N-body simulations. This is not surprising since Sheth and Diaferio & Geller have shown that a sum of Gaussian random velocities weighted by the Press-Schechter mass function provides a good description of the exponential distribution of velocities seen in N-body simulations.
The model underestimates the N-body results in Fig. 1 in both real and redshift space. We show the prediction for the ratio of redshift-space to real-space power in Fig. 2 compared to the same ratio from the N-body simulations. Here we see that the agreement is very good over more than 2 decades in length scale.
## 4 Conclusions
Many of the features of the power spectrum of density fluctuations in the universe can be simply understood in a model based virialized halos. Important ingredients in the model are that the halos be biased tracers of the linear power spectrum and have a uniform profile with a correlation between the internal structure and the mass which should span a wide range sampling a a Press-Schechter like mass function. Within this model it is easy to account accurately for redshift space distortions which alter the power on large and small scales. The effects of virialized motions within halos suppress the โPoissonโ or 1-halo term, so the redshift space power spectrum traces the linear theory result to smaller physical scales. This effect could be responsible for the fact that perturbation theory is known to work better in redshift space than real space when compared to numerical simulations.
## Acknowledgments
I thank U. Seljak and R. Sheth for useful comments on the manuscript. This work was supported by a grant from the US National Science Foundation.
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# Untitled Document
Banach space representations and Iwasawa theory
P. Schneider, J. Teitelbaum
The lack of a $`p`$-adic Haar measure causes many methods of traditional representation theory to break down when applied to continuous representations of a compact $`p`$-adic Lie group $`G`$ in Banach spaces over a given $`p`$-adic field $`K`$. For example, the abelian group $`G=ZZ_p`$ has an enormous wealth of infinite dimensional, topologically irreducible Banach space representations, as may be seen in the paper by Diarra \[Dia\]. We therefore address the problem of finding an additional โfinitenessโ condition on such representations that will lead to a reasonable theory. We introduce such a condition that we call โadmissibilityโ. We show that the category of all admissible $`G`$-representations is reasonable โ in fact, it is abelian and of a purely algebraic nature โ by showing that it is anti-equivalent to the category of all finitely generated modules over a certain kind of completed group ring $`K[[G]]`$.
In the first part of our paper we deal with the general functional-analytic aspects of the problem. We first consider the relationship between $`K`$-Banach spaces and compact, linearly topologized $`o`$-modules where $`o`$ is the ring of integers in $`K`$. As a special case of ideas of Schikhof \[Sch\], we recall that there is an anti-equivalence between the category of $`K`$-Banach spaces and the category of torsionfree, linearly compact $`o`$-modules, provided one tensors the Hom-spaces in the latter category with $`\text{ }Q`$. In addition we have to investigate how this functor relates certain locally convex topologies on the Hom-spaces in the two categories. This will enable us then to derive a version of this anti-equivalence with an action of a profinite group $`G`$ on both sides relating $`K`$-Banach space representations of $`G`$ and certain topological modules for the ring $`K[[G]]:=K_oo[[G]]`$.
Having established these topological results we assume that $`G`$ is a compact $`p`$-adic Lie group and focus our attention on the Banach representations of $`G`$ that correspond under the anti-equivalence to finitely generated modules over the ring $`K[[G]]`$. We characterize such Banach space representations intrinsically. We then show that the theory of such โadmissibleโ representations is purely algebraic โ one may โforgetโ about topology and instead study finitely generated modules over the noetherian ring $`K[[G]]`$.
As an application of our methods we determine the topological irreducibility as well as the intertwining maps for representations of $`GL_2(ZZ_p)`$ obtained by induction of a continuous character from the subgroup of lower triangular matrices. Let us stress the fact that topological irreducibility for an admissible Banach space representation corresponds to the algebraic simplicity of the dual $`K[[G]]`$-module. It is indeed the latter which we will analyze. These results are a complement to the treatment of the locally analytic principal series representations studied in \[ST1\].
Throughout this paper $`K`$ is a finite extension of $`\text{ }Q_p`$ with ring of integers $`oK`$ and absolute value $`||`$. A topological $`o`$-module is called linear-topological if the zero element has a fundamental system of open neighbourhoods consisting of $`o`$-submodules. We let
$$\begin{array}{cc}\mathrm{Mod}_{\mathrm{top}}(o):=& \text{category of all Hausdorff linear-topological}o\text{-modules}\\ & \text{with morphisms being all continuous}o\text{-linear maps.}\end{array}$$
1. A duality for Banach spaces
In this section we will recall a certain duality theory for $`K`$-Banach spaces due to Schikhof (\[Sch\]). Because of the fundamental role it will play in our later considerations and since it is quite easy over locally compact fields we include the proofs. We set
$$\begin{array}{cc}\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o):=& \text{the full subcategory in}\mathrm{Mod}_{\mathrm{top}}(o)\text{of all}\\ & \text{torsionfree and compact linear-topological}\\ & o\text{-modules.}\end{array}$$
Remark 1.1:
i. An $`o`$-module is torsionfree if and only if it is flat;
ii. a compact linear-topological $`o`$-module $`M`$ is flat if and only if $`M\underset{iI}{}o`$ for some set $`I`$.
Proof: i. \[B-CA\] Chap. I ยง2.4 Prop. 3(ii). ii. \[SGA3\] $`\mathrm{Exp}.\mathrm{VII}_B`$ (0.3.8).
For later purposes let us note that any $`o`$-module $`M`$ in $`\mathrm{Mod}_{\mathrm{top}}(o)`$ has a unique largest quotient module $`M_{\mathrm{cot}}`$ which is Hausdorff and torsionfree: If $`(M_j)_{jJ}`$ is the family of all torsionfree Hausdorff quotient modules of $`M`$ then $`M_{\mathrm{cot}}`$ is the coimage of the natural map $`M\underset{jJ}{}M_j`$.
For any $`o`$-module $`M`$ in $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)`$ we can construct the $`K`$-Banach space
$$M^d:=\mathrm{Hom}_o^{\mathrm{cont}}(M,K)\text{with norm}\mathrm{}:=\underset{mM}{max}|\mathrm{}(m)|.$$
This defines a contravariant additive functor
$$\begin{array}{ccc}\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)& & \mathrm{Ban}(K)\\ \\ M& & M^d\end{array}$$
into the category $`\mathrm{Ban}(K)`$ of all $`K`$-Banach spaces with morphisms being all continuous $`K`$-linear maps. Actually all maps in the image of this functor are norm decreasing. The groups of homomorphisms in $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)`$ are $`o`$-modules whereas in $`\mathrm{Ban}(K)`$ they are $`K`$-vector spaces. The above functor therefore extends naturally to a contravariant additive functor
$$\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)_{\text{ }Q}\mathrm{Ban}(K).$$
Here $`๐_{\text{ }Q}`$, for any additive category $`๐`$, denotes the additive category with the same objects as $`๐`$ and such that
$$\mathrm{Hom}_{๐_{\text{ }Q}}(A,B):=\mathrm{Hom}_๐(A,B)\text{ }Q$$
for any two objects $`A,B`$ in $`๐`$ with the composition of morphisms in $`๐_{\text{ }Q}`$ being the $`\text{ }Q`$-linear extension of the composition in $`๐`$.
Theorem 1.2:
The functor
$$\begin{array}{ccc}\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)_Q& \stackrel{}{}& \mathrm{Ban}(K)\\ \\ M& & M^d\end{array}$$
is an anti-equivalence of categories.
Proof: Let $`\mathrm{Ban}(K)^1`$ denote the category of all $`K`$-Banach spaces $`(E,)`$ such that $`E|K|`$ with morphisms being all norm decreasing $`K`$-linear maps. Clearly our functor factorizes into
$$\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)\stackrel{(.)^d}{}\mathrm{Ban}(K)^1\stackrel{\mathrm{forget}}{}\mathrm{Ban}(K).$$
For any $`K`$-Banach space $`(E,)`$ we may define by $`v^{}:=\mathrm{inf}\{r|K|:rv\}`$ another norm $`^{}`$ on $`E`$ satisfying $`E^{}|K|`$. Because of $`|\pi |v/v^{}1`$ for $`v0`$, where $`\pi `$ is a prime element of $`K`$, the two norms $``$ and $`^{}`$ are equivalent. It follows that the right hand functor above induces an equivalence of categories
$$(\mathrm{Ban}(K)^1)_{\text{ }Q}\stackrel{}{}\mathrm{Ban}(K).$$
We therefore are reduced to show that
$$\begin{array}{ccc}\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)& & \mathrm{Ban}(K)^1\\ \\ M& & M^d\end{array}$$
is an anti-equivalence of categories. Let $`(E,)`$ be a $`K`$-Banach space and denote by $`E^{}:=\{vE:v1\}`$ its unit ball. Then
$$E^d:=\mathrm{Hom}_o(E^{},o)\text{with the topology of pointwise convergence}$$
is a linear-topological $`o`$-module which is torsionfree and complete. In fact, $`E^d`$ is the unit ball of the dual Banach space $`E^{}`$ but equipped with the weak topology. Since
$$\begin{array}{ccc}E^d& & \underset{vE^{}}{}o\\ \\ \lambda & & (\lambda (v))_v\end{array}$$
is a topological embedding we see that $`E^d`$ is compact. This defines a functor
$$\begin{array}{ccc}\mathrm{Ban}(K)^1& & \mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)\\ \\ (E,)& & E^d.\end{array}$$
It is an immediate consequence of Remark 1.1 that, for an $`o`$-module $`M`$ in $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)`$, the $`o`$-linear map
$$\begin{array}{ccc}\iota _M:M& & (M^d)_s^{}\\ \\ m& & [\mathrm{}\mathrm{}(m)]\end{array}$$
into the weak dual $`(M^d)_s^{}`$ of the Banach space $`M^d`$ is injective. Since it is easily seen to be continuous the compactness of $`M`$ implies that $`\iota _M`$ is a closed embedding. By definition the image of $`\iota _M`$ is contained in $`M^{dd}`$. Assume now that there is a $`\lambda M^{dd}\mathrm{im}(\iota _M)`$. Since $`\mathrm{im}(\iota _M)`$ is closed in $`(M^d)_s^{}`$ there is, by Hahn-Banach (\[Mon\] V.1.2 Thm. 5(ii) or \[NFA\] 13.3), a continuous linear form on $`(M^d)_s^{}`$ which in absolute value is $`1`$ on $`\lambda `$ and is $`<1`$ on $`\mathrm{im}(\iota _M)`$. But, as another consequence of Hahn-Banach (\[NFA\] 9.7), any continuous linear form on $`(M^d)_s^{}`$ is given by evaluation in a vector in $`M^d`$. Hence we find an $`\mathrm{}M^d`$ such that $`|\lambda (\mathrm{})|1`$ and $`|\mathrm{}(M)|<1`$. The latter implies $`\mathrm{}<1`$ so that $`\lambda (\mathrm{})\lambda \mathrm{}|<1`$ which is a contradiction. We obtain that $`\iota _M:M\stackrel{}{}M^{dd}`$, in fact, is a topological isomorphism. This means that the $`\iota _M`$ constitute a natural isomorphism between the identity functor and the functor $`(.)^{dd}`$ on $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)`$. On the other hand any $`(E,)`$ in $`\mathrm{Ban}(K)^1`$ is isometric to a Banach space $`c_0(I)`$ for some set $`I`$ (\[Mon\] IV.3 Cor. 1 or \[NFA\] 10.1). A straightforward explicit computation shows that $`c_0(I)^{dd}=c_0(I)`$. The functor $`MM^d`$ therefore is fully faithful as well as essentially surjective and consequently an equivalence.
The exactness properties of this functor are as follows.
Proposition 1.3:
For any map $`f:MN`$ in $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)`$ we have:
i. $`ker(f)^d=M^d/\overline{f^d(N^d)}`$;
ii. $`[coker(f)_{\mathrm{cot}}]^d=ker(f^d)`$ ;
iii. $`f`$ is surjective if and only if $`f^d`$ is an isometry.
Proof: i. The submodules $`ker(f)`$ and $`im(f)`$ lie again in $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)`$. It follows from \[SGA3\] $`\mathrm{Exp}.\mathrm{VII}_B`$ (0.3.7) that the surjection $`Mim(f)`$ splits, i.e., we have $`Mker(f)im(f)`$ in $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)`$. It suffices therefore to consider the case where $`f`$ is injective and to show that then the image of $`f^d`$ is dense in $`M^d`$. If not we find by Hahn-Banach a nonzero continuous linear form $`\lambda `$ on $`M^d`$ which vanishes on the image of $`f^d`$. Up to scaling we may assume that $`\lambda M^{dd}`$, i.e., that there is a nonzero $`mM`$ such that $`\lambda (\mathrm{})=\mathrm{}(m)`$. The vanishing property of $`\lambda `$ means of course that $`f(m)=0`$ which is a contradiction.
iii. If $`f`$ is surjective then $`f^d`$ is an isometry by construction. Suppose now that $`f^d`$ is an isometry. Let $`nN`$; we view $`n`$ as a linear form in the unit ball of the dual Banach space $`(N^d)^{}`$. By Hahn-Banach $`n`$ extends (via $`f^d`$) to a linear form in the unit ball of $`(M^d)^{}`$; this means of course that we find an $`mM`$ such that $`f(m)=n`$.
ii. Let $`E`$ denote the kernel of $`f^d`$. Then $`E^d`$ is, by iii., a torsionfree Hausdorff quotient of $`coker(f)`$. On the other hand $`[coker(f)_{\mathrm{cot}}]^d`$ clearly is a subspace of $`ker(f^d)`$.
Let $`M`$ be a module in $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)`$. Since $`M`$ is torsionfree it is an $`o`$-submodule of the $`K`$-vector space $`M_K:=M_oK`$. Thm. 1.2 tells us that there is a natural identification
$$M_K=\mathrm{Hom}_o^{\mathrm{cont}}(o,M)\text{ }Q=\mathrm{Hom}_K^{\mathrm{cont}}(M^d,K)=(M^d)^{}$$
between $`M_K`$ and the continuous dual $`(M^d)^{}`$ of the Banach space $`M^d`$. We always equip $`M_K`$ with the finest locally convex topology such that the inclusion $`MM_K`$ is continuous. An $`o`$-submodule $`LM_K`$ is open if and only if $`\alpha LM`$ is open in $`M`$ for any $`0\alpha o`$. By construction this topology has the property that
$$\mathrm{Hom}_o^{\mathrm{cont}}(M,V)=(M_K,V)$$
for any locally convex $`K`$-vector space $`V`$ where, following a common convention, we let $`(.,.)`$ denote the vector space of continuous linear maps between two locally convex $`K`$-vector spaces. In particular $`M^d`$ at least as a vector space is the continuous dual $`(M_K)^{}`$. Since under the identification $`M_K=(M^d)^{}`$ the topology of $`M`$ is induced by the weak topology on $`(M^d)^{}`$ we also see that the identification map $`M_K\stackrel{=}{}(M^d)_s^{}`$ is continuous. This shows that $`M_K`$ is Hausdorff and that $`M_K`$ also induces the given topology on $`M`$.
Lemma 1.4
The locally convex $`K`$-vector space $`M_K`$, for any $`M`$ in $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)`$, is complete.
Proof: Fix a prime element $`\pi `$ of $`K`$. Let $`\mathrm{}`$ be a minimal Cauchy filter on $`M_K`$. We first show that there is a $`mIN`$ such that
$$F\pi ^mM\mathrm{}\text{for all}F\mathrm{}.$$
Otherwise there exists for any $`nIN`$ a $`F_n\mathrm{}`$ with $`F_n\pi ^nM=\mathrm{}`$. By the minimality of $`\mathrm{}`$ we may assume that
$$F_n=F_n+L_n\text{for some open }o\text{-submodule}L_nM_K$$
(\[B-GT\] Chap. II ยง3.2 Prop. 5). We also may assume that the $`L_n`$ form a decreasing sequence $`L_1L_2\mathrm{}`$. The $`o`$-submodule
$$L:=\underset{nIN}{}(L_n\pi ^nM)$$
is open in $`M_K`$ since $`L\pi ^nML_n\pi ^nM`$ for all $`nIN`$. The $`L_n`$ being decreasing and the $`\pi ^nM`$ being increasing it is clear that
$$LL_n+\pi ^nM\text{for all}nIN.$$
As a Cauchy filter $`\mathrm{}`$ must contain a coset $`v+L`$ for some $`vM_K`$. If $`n_0IN`$ is chosen in such a way that $`v\pi ^{n_0}M`$ we have $`L_{n_0}+\pi ^{n_0}M\mathrm{}`$. Both sets $`F_{n_0}`$ and $`L_{n_0}+\pi ^{n_0}M`$ belonging to the filter $`\mathrm{}`$ we obtain $`F_{n_0}(L_{n_0}+\pi ^{n_0}M)\mathrm{}`$, i.e., $`F_{n_0}\pi ^{n_0}M=(F_{n_0}+L_{n_0})\pi ^{n_0}M\mathrm{}`$ which is a contradiction. We see that
$$\mathrm{}_m:=\{F\pi ^mM:F\mathrm{}\}$$
for an appropriate $`mIN`$ is a filter on $`\pi ^mM`$. Since $`\pi ^mM`$ is compact in $`M_K`$ the filter $`\mathrm{}_m`$ being also a Cauchy filter has to be convergent. By \[B-GT\] Chap. II ยง3.2 Cor. 3 then $`\mathrm{}`$ is convergent, too. This proves that $`M_K`$ is complete.
Lemma 1.5:
For any two $`o`$-modules $`M`$ and $`N`$ in $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)`$ we have:
i. For any compact subset $`CN_K`$ the closed $`o`$-submodule in $`N_K`$ topologically generated by $`C`$ is compact as well;
ii. for any compact subset $`CN_K`$ there is a $`0\alpha o`$ such that $`\alpha CN`$;
iii. $`\mathrm{Hom}_o^{\mathrm{cont}}(M,N)\text{ }Q=(M_K,N_K)`$;
iv. passing to the transpose induces a $`K`$-linear isomorphism
$$(M_K,N_K)\stackrel{}{}(N^d,M^d).$$
Proof: i. Let $`<C>`$ denote the $`o`$-submodule generated by $`C`$. Let $`LN_K`$ be any open and therefore also closed $`o`$-submodule. Since $`C`$ is compact we find finitely many $`c_1,\mathrm{},c_nC`$ with $`C(c_1+L)\mathrm{}(c_n+L)`$. Then $`<C>`$ is contained in $`oc_1+\mathrm{}+oc_n+L`$. But $`oc_1+\mathrm{}+oc_n`$ is compact, too, so that we again find finitely many $`a_1,\mathrm{},a_moc_1+\mathrm{}+oc_n`$ with
$$oc_1+\mathrm{}+oc_n(a_1+L)\mathrm{}(a_m+L).$$
Together we obtain
$$<C>\underset{1im}{}a_i+L$$
and since the right hand side is closed the closure $`\overline{<C>}`$ of $`<C>`$ also satisfies
$$\overline{<C>}\underset{1im}{}a_i+L.$$
Since $`L`$ was arbitrary this implies by \[B-GT\] Chap. II ยง4.2 Thm. 3 that $`\overline{<C>}`$ is precompact. On the other hand, as a consequence of Lemma 1.4, $`\overline{<C>}`$ is Hausdorff and complete. Hence $`\overline{<C>}`$ is compact.
ii. By i. we may assume that $`C`$ is a compact $`o`$-submodule of $`N_K`$. Fix a prime element $`\pi `$ of $`K`$ and put $`C_n:=C\pi ^nN`$ for any $`nIN`$. These $`C_n`$ form an increasing sequence $`C_1C_2\mathrm{}`$ of compact $`o`$-submodules of $`C`$ such that $`C=\underset{nIN}{}C_n`$. We have to show that $`C_m=C`$ holds for some $`mIN`$. Being empty the subset $`\underset{nIN}{}(C\backslash C_n)`$ is not dense in $`C`$. As a compact space $`C`$ in particular is a Baire space (\[B-GT\] Chap. IX ยง5.3 Thm. 1) so that already some $`C\backslash C_n`$ is not dense in $`C`$. This means that $`C_n`$ contains a non-empty open subset of $`C`$. It is then itself an open $`o`$-submodule and therefore has to be of finite index in $`C`$. Our claim obviously follows from that.
iii. We have $`\mathrm{Hom}_o^{\mathrm{cont}}(M,N)\text{ }Q=\mathrm{Hom}_o^{\mathrm{cont}}(M,N)_oK=\mathrm{Hom}_o^{\mathrm{cont}}(M,N_K)=(M_K,N_K)`$ where the second identity is a consequence of the second assertion.
iv. This follows from iii. and Thm. 1.2.
The assertion iii. in Lemma 1.5 in particular means that $`M_K`$ and $`N_K`$ are isomorphic in $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)_{\text{ }Q}`$ if and only if they are isomorphic as locally convex vector spaces.
For any two $`o`$-modules $`M`$ and $`N`$ in $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)`$ we always view $`\mathrm{Hom}_o^{\mathrm{cont}}(M,N)`$ as a linear-topological $`o`$-module by equipping it with the topology of compact convergence. As a consequence of Lemma 1.5 this topology is induced by the topology of compact convergence on the vector space $`(M_K,N_K)`$. We write $`_{cc}(M_K,N_K)`$ for $`(M_K,N_K)`$ equipped with the finest locally convex topology such that the inclusion $`\mathrm{Hom}_o^{\mathrm{cont}}(M,N)_{cc}(M_K,N_K)`$ is continuous. By a similar argument as before $`_{cc}(M_K,N_K)`$ is Hausdorff and the latter inclusion is a topological embedding. Moreover, by Lemma 1.5, $`_{cc}(M_K,N_K)`$ is, in both variables, a functor on $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)_{\text{ }Q}`$.
Given two $`K`$-Banach spaces $`E_1`$ and $`E_2`$ we write, following traditional usage, $`_s(E_1,E_2)`$ for the vector space $`(E_1,E_2)`$ equipped with the locally convex topology of pointwise convergence. We write $`_{bs}(E_1,E_2)`$ for $`(E_1,E_2)`$ equipped with the finest locally convex topology which coincides with the topology of pointwise convergence on any equicontinuous subset in $`(E_1,E_2)`$. Corresponding to any choice of defining norms $`_i`$ on $`E_i`$ for $`i=1,2`$ we have the operator norm $``$ on $`(E_1,E_2)`$. A subset in $`(E_1,E_2)`$ is equicontinuous if and only if it is bounded with respect to $``$. Hence the topology of $`_{bs}(E_1,E_2)`$ can equivalently be characterized as being the finest locally convex topology which induces the topology of pointwise convergence on the unit ball with respect to $``$ in $`(E_1,E_2)`$.
Proposition 1.6:
Passing to the transpose induces, for any $`o`$-modules $`M`$ and $`N`$ in $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)`$, an isomorphism of locally convex $`K`$-vector spaces
$$_{cc}(M_K,N_K)\stackrel{}{}_{bs}(N^d,M^d).$$
Proof: It is clear from our preliminary discussion that it suffices to show that the $`o`$-linear isomorphism
$$\mathrm{Hom}_o^{\mathrm{cont}}(M,N)\stackrel{}{}\{f(N^d,M^d):f1\}$$
given by the transpose is topological provided the left, resp. right, hand side carries the topology of compact, resp. pointwise convergence. We recall from the proof of Thm. 1.2 that $`M`$ is the unit ball in the dual Banach space $`(M^d)^{}`$ equipped with the weak topology; we also have seen there that the closed equicontinuous subsets of the weak dual $`(M^d)_s^{}`$ are compact. By the Banach-Steinhaus theorem (\[Tie\] Thm. 4.3) a subset of $`(M^d)_s^{}`$ is equicontinuous if and only if it is bounded. Clearly any compact subset is bounded. It follows that for a closed subset of $`(M^d)_s^{}`$ the following properties are equivalent: Bounded, equicontinuous, bounded for the dual Banach norm, compact. This shows that the topology of compact convergence on $`\mathrm{Hom}_o^{\mathrm{cont}}(M,N)`$ is induced by the strong topology on $`((M^d)_s^{},(N^d)_s^{})`$. Our assertion therefore will be a consequence of the quite general fact that for any two $`K`$-Banach spaces $`E_1`$ and $`E_2`$ the transpose induces a topological isomorphism
$$_s(E_1,E_2)\stackrel{}{}_b((E_2)_s^{},(E_1)_s^{})$$
where on the right hand side the subscript $`b`$ indicates, as usual, the strong topology. This is straightforward from the definitions and the fact that
$$\begin{array}{ccc}\text{set of all open}& \stackrel{}{}& \text{set of all closed equicontinuous}\\ o\text{-submodules in}E_2& & o\text{-submodules in}(E_2)_s^{}\\ \\ L& & L^p:=\{\mathrm{}(E_2)^{}:|\mathrm{}(v)|1\text{for any}vL\}\end{array}$$
is a bijection which is a direct consequence of the Hahn-Banach theorem.
2. Iwasawa modules and representations
ยฟFrom now on we let $`G`$ denote a fixed profinite group. The completed group ring of $`G`$ (over $`o`$) is defined to be
$$o[[G]]:=\underset{H\mathrm{๐ฉ}}{\underset{}{lim}}o[G/H]$$
where $`\mathrm{๐ฉ}=\mathrm{๐ฉ}(G)`$ denotes the family of all open normal subgroups of $`G`$. In a natural way $`o[[G]]`$ is a torsionfree and compact linear-topological $`o`$-module; the ring multiplication is continuous. The surjections $`o[G]o[G/H]`$ for $`H\mathrm{๐ฉ}`$ induce in the limit a ring homomorphism
$$o[G]o[[G]]$$
whose image is dense and which is injective (\[Laz\] II.2.2.3.1). Being the projective limit of the inclusions $`G/Ho[G/H]`$ the composed map
$$G\stackrel{}{}o[G]o[[G]]$$
is continuous and hence, by compactness, a homeomorphism onto its image.
Consider now a module $`M`$ in $`\mathrm{Mod}_{\mathrm{top}}(o)`$ and let $`C(G,M)`$ denote the $`o`$-module of all continuous maps from $`G`$ into $`M`$. It follows from the above discussion that the $`o`$-linear map
$$\begin{array}{ccc}\mathrm{Hom}_o^{\mathrm{cont}}(o[[G]],M)& & C(G,M)\\ \\ f& & f|G\end{array}$$
is well defined and injective.
Lemma 2.1:
For any complete $`o`$-module $`M`$ in $`\mathrm{Mod}_{\mathrm{top}}(o)`$ the map
$$\begin{array}{ccc}\mathrm{Hom}_o^{\mathrm{cont}}(o[[G]],M)& \stackrel{}{}& C(G,M)\\ \\ f& & f|G\end{array}$$
is a bijection.
Proof: We extend a given $`\phi C(G,M)`$ $`o`$-linearly to $`o[G]`$. By the completeness assumption and the density of $`o[G]`$ in $`o[[G]]`$ it suffices to show that this extension, which we again denote by $`\phi `$, is continuous with respect to the topology induced by $`o[[G]]`$. Fix an open $`o`$-submodule $`LM`$. By the uniform continuity of $`\phi `$ on $`G`$ we find an $`H\mathrm{๐ฉ}`$ such that
$$\phi (g_iH)\phi (g_i)+L$$
for all $`g_i`$ in a system of representatives for the left cosets of $`H`$ in $`G`$ (compare \[Laz\] II.2.2.5). Let $`\alpha o`$ be some element such that $`\alpha \phi (g_i)L`$ for all $`i`$. The $`o`$-submodule
$$L^{}:=\underset{i}{}\{\underset{gg_iH}{}r_gg:\underset{g}{}r_g\alpha o\}$$
then is open in $`o[G]`$ and we have
$$\phi (L^{})\underset{i}{}(\alpha o\phi (g_i)+L)L.$$
We set $`K[[G]]:=o[[G]]_K`$. This is a locally convex vector space as well as a $`K`$-algebra such that the multiplication is separately continuous.
Corollary 2.2
For any quasi-complete Hausdorff locally convex $`K`$-vector space $`V`$ we have the $`K`$-linear isomorphism
$$\begin{array}{ccc}(K[[G]],V)& \stackrel{}{}& C(G,V)\\ \\ f& & f|G.\end{array}$$
Proof: The map is clearly well defined and injective. For the surjectivity let $`\phi C(G,V)`$. Define $`M`$ to be the closed $`o`$-submodule of $`V`$ topologically generated by $`\phi (G)`$. This $`M`$ lies in $`\mathrm{Mod}_{\mathrm{top}}(o)`$. Since $`G`$ is compact $`M`$ is bounded in $`V`$. The quasi-completeness of $`V`$ therefore ensures that $`M`$ is complete. Hence we have, by Lemma 2.1, a continuous $`o`$-linear map $`f:o[[G]]MV`$ such that $`f|G=\phi `$. The $`K`$-linear extension of $`f`$ then is the preimage of $`\phi `$ we were looking for.
We will apply these results to obtain a $`G`$-equivariant version of the duality theorem of the previous section.
Definition:
A $`K`$-Banach space representation $`E`$ of $`G`$ is a $`K`$-Banach space $`E`$ together with a $`G`$-action by continuous linear automorphisms such that the map $`G\times EE`$ describing the action is continuous.
We define
$$\begin{array}{cc}\mathrm{Ban}_G(K):=& \text{category of all}K\text{-Banach representations of}G\text{with}\\ & \text{morphisms being all}G\text{-equivariant continuous linear maps.}\end{array}$$
As a consequence of the Banach-Steinhaus theorem (\[Tie\] Thm. 4.1.1), to give a $`K`$-Banach representation of $`G`$ on the $`K`$-Banach space $`E`$ is the same as to give a continuous homomorphism $`G_s(E,E)`$. But $`_s(E,E)`$ is quasi-complete and Hausdorff (\[B-TVS\] III.27 Cor. 4 or \[NFA\] 7.14). Hence such a homomorphism extends, by Cor. 2.2, uniquely to a continuous $`K`$-linear map $`K[[G]]_s(E,E)`$. By a density argument the latter map is a $`K`$-algebra homomorphism. This shows that a $`K`$-Banach space representation of $`G`$ on $`E`$ is the same as a separately continuous action $`K[[G]]\times EE`$ of the algebra $`K[[G]]`$ on $`E`$.
Since the image of $`o[[G]]`$ in $`_s(E,E)`$ under the above homomorphism is compact and hence (by Banach-Steinhaus) equicontinuous we also have that a $`K`$-Banach space representation of $`G`$ on $`E`$ is the same as a continuous (unital) homomorphism of $`K`$-algebras $`K[[G]]_{bs}(E,E)`$.
Definition:
An Iwasawa $`G`$-module over $`o`$ is an $`o`$-module $`M`$ in $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)`$ together with a continuous (left) action $`o[[G]]\times MM`$ of the compact $`o`$-algebra $`o[[G]]`$ on $`M`$ such that the induced $`o`$-action on $`M`$ is the given $`o`$-module structure.
Let
$$\begin{array}{cc}\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o[[G]]):=& \text{category of all Iwasawa}G\text{-modules over}o\\ & \text{with morphisms being all continuous}o[[G]]\text{-}\\ & \text{module homomorphisms.}\end{array}$$
A continuous (unital) homomorphism of $`K`$-algebras
$$K[[G]]_{cc}(M_K,M_K)$$
$`()`$
for some $`M_K`$ in $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)_{\text{ }Q}`$ induces a continuous map $`o[[G]]\mathrm{Hom}_o^{\mathrm{cont}}(M,M_K)`$ where the right hand side carries the topology of compact convergence. By \[B-GT\] Chap. X ยง3.4 Thm. 3 this is the same as a continuous map $`o[[G]]\times MM_K`$. According to Lemma 1.5.ii the image of this latter map is contained in $`\alpha ^1M`$ for some $`0\alpha o`$. If $`N`$ denotes the closed $`o`$-submodule of $`M_K`$ topologically generated by this image we therefore have $`N_K=M_K`$, and the above homomorphism of $`K`$-algebras $`()`$ is the tensor product with $`K`$ of a continuous (unital) homomorphism of $`o`$-algebras
$$o[[G]]\mathrm{Hom}_o^{\mathrm{cont}}(N,N).$$
Again by \[B-GT\] loc. cit. this is the same as a continuous action $`o[[G]]\times NN`$ of the compact $`o`$-algebra $`o[[G]]`$ on the $`o`$-module $`N`$ in $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o)`$ which extends the $`o`$-module structure.
Hence we see that to give a continuous (unital) homomorphism of $`K`$-algebras $`()`$ is the same as to give an object in the category $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o[[G]])_{\text{ }Q}`$. By combining this discussion with Prop. 1.6 we arrive at the following equivariant version of Thm. 1.2.
Theorem 2.3:
The functor
$$\begin{array}{ccc}\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o[[G]])_Q& \stackrel{}{}& \mathrm{Ban}_G(K)\\ \\ M& & M^d\end{array}$$
is an anti-equivalence of categories.
3. Admissible representations
In order to obtain a reasonable theory of Banach space representations it seems necessary to impose certain additional finiteness conditions. The first idea is to consider only those $`K`$-Banach space representations of $`G`$ which correspond, under the duality of the previous section, to finitely generated $`K[[G]]`$-modules. As a consequence of the compactness of the ring $`o[[G]]`$ it will turn out that the theory of these representations in fact is completely algebraic in nature. In order to obtain an intrinsic characterization we will assume in this section that $`G`$ is a compact $`p`$-adic Lie group. We then have:
\- The subfamily of all topologically finitely generated pro-$`p`$-groups in $`\mathrm{๐ฉ}=\mathrm{๐ฉ}(G)`$ is cofinal (\[B-GAL\] Chap. III ยง1.1 Prop. 2(iii) and ยง7.3 and 4 and \[Laz\] III.2.2.6 and III.3.1.3).
\- The ring $`o[[G]]`$ is left and right noetherian (\[Laz\] V.2.2.4).
The ring $`K[[G]]`$ then is left and right noetherian as well.
Definition:
A $`K`$-Banach space representation $`E`$ of $`G`$ is called admissible if there is a $`G`$-invariant bounded open $`o`$-submodule $`LE`$ such that, for any $`H\mathrm{๐ฉ}`$, the $`o`$-submodule $`(E/L)^H`$ of $`H`$-invariant elements in the quotient $`E/L`$ is of cofinite type.
We recall that an $`o`$-module $`N`$ is called of cofinite type if its Pontrjagin dual $`\mathrm{Hom}_\mathrm{o}(N,K/o)`$ is a finitely generated $`o`$-module. We also point out that an arbitrary open $`o`$-submodule $`LE`$ contains the $`G`$-invariant open $`o`$-submodule $`\underset{gG}{}gL`$.
Let
$$\begin{array}{cc}\mathrm{Ban}_G^{\mathrm{adm}}(K):=& \text{the full subcategory in}\mathrm{Ban}_G(K)\\ & \text{of all admissible representations. }\end{array}$$
On the other hand we let $`\mathrm{Mod}_{\mathrm{fg}}^{\mathrm{fl}}(o[[G]])`$, resp. $`\mathrm{Mod}_{\mathrm{fg}}(K[[G]])`$, denote the category of all finitely generated and $`o`$-torsionfree (left unital) $`o[[G]])`$-modules, resp. of all finitely generated (left unital) $`K[[G]])`$-modules. It is clear that
$$\mathrm{Mod}_{\mathrm{fg}}^{\mathrm{fl}}(o[[G]])_{\text{ }Q}=\mathrm{Mod}_{\mathrm{fg}}(K[[G]]).$$
Since $`K[[G]]`$ is noetherian the category $`\mathrm{Mod}_{\mathrm{fg}}(K[[G]])`$ is abelian.
Proposition 3.1:
i. A finitely generated $`o[[G]]`$-module $`M`$ carries a unique Hausdorff topology - its canonical topology - such that the action $`o[[G]]\times MM`$ is continuous;
ii. any submodule of a finitely generated $`o[[G]]`$-module is closed in the canonical topology;
iii. any $`o[[G]]`$-linear map between two finitely generated $`o[[G]]`$-modules is continuous for the canonical topologies.
Proof: Since $`o[[G]]`$ is compact and noetherian this is an easy exercise. But we point out that the assertions hold for any compact ring by \[AU\] Cor. 1.10.
It follows that equipping a module in $`\mathrm{Mod}_{\mathrm{fg}}^{\mathrm{fl}}(o[[G]])`$ with its canonical topology induces a fully faithful embedding $`\mathrm{Mod}_{\mathrm{fg}}^{\mathrm{fl}}(o[[G]])\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o[[G]])`$. This then in turn induces a fully faithful embedding $`\mathrm{Mod}_{\mathrm{fg}}(K[[G]])\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o[[G]])_{\text{ }Q}`$. In other words we can and will view $`\mathrm{Mod}_{\mathrm{fg}}(K[[G]])`$ as a full subcategory of $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o[[G]])_{\text{ }Q}`$.
For each $`H\mathrm{๐ฉ}`$ let $`I_H`$ denote the kernel of the projection map $`o[[G]]o[G/H]`$. This is a family of 2-sided ideals in $`o[[G]]`$ which converges to zero. As a left (or right) ideal $`I_H`$ is generated by the elements $`h1`$ for $`hH`$. For the sake of completeness we include a proof of the following well known fact.
Lemma 3.2:
Let $`H\mathrm{๐ฉ}`$ be a pro-$`p`$-group; then the ideal powers $`I_H^n`$, for $`nIN`$, converge to zero.
Proof: We may assume that $`G`$ is finite. Let $`\pi `$ denote a prime element in $`o`$ and $`k:=o/\pi o`$ the residue field of $`o`$. By Cliffordโs theorem (\[CR\] (49.2)) and \[Ser\] IXยง1 the ideal $`\mathrm{ker}(k[G]k[G/H])`$ is contained in the radical of the ring $`k[G]`$. Since this radical is nilpotent we have $`I_H^m\pi o[G]`$ for some $`mIN`$.
Lemma 3.3:
Let $`H\mathrm{๐ฉ}`$ be a pro-$`p`$-group; a module $`M`$ in $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o[[G]])`$ is finitely generated over $`o[[G]]`$ if and only if $`M/I_HM`$ is finitely generated over $`o`$.
Proof: This is the well known Nakayama lemma; compare \[BH\] for a thorough discussion.
Lemma 3.4:
A $`K`$-Banach space representation $`E`$ of $`G`$ is admissible if and only if the dual space $`E^{}`$ is finitely generated over $`K[[G]]`$.
Proof: Let us first assume that $`E^{}`$ is finitely generated over $`K[[G]]`$. There is then a finitely generated $`o[[G]]`$-submodule $`ME^{}`$ such that $`E^{}=M_K`$. After equipping $`M`$ with its canonical topology we have $`E=M^d`$. Moreover $`L:=\mathrm{Hom}_o^{\mathrm{cont}}(M,o)`$ is a $`G`$-invariant bounded open $`o`$-submodule in $`E`$. It follows from Remark 1.1 that $`E/L=\mathrm{Hom}_o^{\mathrm{cont}}(M,K/o)`$ (where $`K/o`$ carries the discrete topology) and hence that
$$(E/L)^H=\mathrm{Hom}_o^{\mathrm{cont}}(M,K/o)^H=\mathrm{Hom}_o^{\mathrm{cont}}(M/I_HM,K/o)$$
$`()`$
for any $`H\mathrm{๐ฉ}`$. Hence $`(E/L)^H`$ is of cofinite type.
On the other hand let now $`H\mathrm{๐ฉ}`$ be a pro-$`p`$-group and $`LE`$ be a $`G`$-invariant bounded open $`o`$-submodule such that $`(E/L)^H`$ is of cofinite type. In the proof of Prop. 1.6 we had recalled that the $`G`$-invariant $`o`$-submodule $`M:=L^p`$ in $`E_s^{}`$ is compact. Since $`L`$ is bounded we have $`E^{}=M_K`$. So the identities $`()`$ apply correspondingly and we obtain that $`\mathrm{Hom}_o^{\mathrm{cont}}(M/I_HM,K/o)`$ is of cofinite type. But since $`I_H`$ is finitely generated as a right ideal the submodule $`I_HM`$ is the image of finitely many copies $`M\times \mathrm{}\times M`$ under a continuous map and hence is closed in $`M`$. By Pontrjagin duality and the Nakayama lemma over $`o`$ applied to the compact $`o`$-module $`M/I_HM`$ the latter therefore is finitely generated over $`o`$. Lemma 3.3 then implies that $`M`$ is finitely generated over $`o[[G]]`$ and hence that $`E^{}`$ is finitely generated over $`K[[G]]`$.
The above proof shows that the defining condition for admissibility only needs to be tested for a single pro-$`p`$-group $`H\mathrm{๐ฉ}`$. On the other hand assume $`E`$ to be an admissible representation of $`G`$ and let $`LE`$ be as in the above definition. Consider an $`H\mathrm{๐ฉ}`$ and an arbitrary $`G`$-invariant open $`o`$-submodule $`L_\mathrm{o}E`$. We claim that $`(E/L_\mathrm{o})^H`$ is of cofinite type. Replacing $`L`$ by $`\alpha L`$ for some appropriate $`0\alpha o`$ we may assume that $`LL_\mathrm{o}`$. As we have seen in the above proof $`M:=L^p`$ is a finitely generated $`o[[G]]`$-module. Since $`o[[G]]`$ is noetherian the $`o[[G]]`$-submodule $`M_\mathrm{o}:=L_\mathrm{o}^p`$ of $`M`$ also is finitely generated. As we have seen this implies that $`(E/L_\mathrm{o})^H`$ is of cofinite type.
Theorem 3.5:
The functor
$$\begin{array}{ccc}\mathrm{Mod}_{\mathrm{fg}}(K[[G]])& \stackrel{}{}& \mathrm{Ban}_G^{\mathrm{adm}}(K)\\ \\ M& & M^d\end{array}$$
is an anti-equivalence of categories.
Proof: Since $`\mathrm{Mod}_{\mathrm{fg}}(K[[G]])`$ is a full subcategory of $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o[[G]])_{\text{ }Q}`$ by Prop. 3.1 this follows from Thm. 2.3 and Lemma 3.4.
As an immediate consequence we obtain that the category $`\mathrm{Ban}_G^{\mathrm{adm}}(K)`$ is abelian.
Corollary 3.6:
The functor $`EE^{}`$ induces a bijection
$$\begin{array}{ccc}\text{set of isomorphism classes}& & \\ \text{of topologically irreducible}& \stackrel{}{}& \text{set of isomorphism classes}\\ \text{admissible}K\text{-Banach space}& & \text{of simple}K[[G]]\text{-modules.}\\ \text{representations of }G\end{array}$$
Proof: For any proper closed $`G`$-invariant subspace $`\{0\}E_\mathrm{o}\underset{}{}E`$ we have, by Hahn-Banach, the exact sequence of dual vector spaces $`0(E/E_\mathrm{o})^{}E^{}E_\mathrm{o}^{}0`$ in which all three terms are nonzero. If the $`K[[G]]`$-module $`E^{}`$ is simple the representation $`E`$ therefore must be topologically irreducible. On the other hand write $`E^{}=M_K`$ for some module $`M`$ in $`\mathrm{Mod}_{\mathrm{fg}}^{\mathrm{fl}}(o[[G]])`$ and let $`\{0\}V\underset{}{}M_K`$ be a proper $`K[[G]]`$-submodule. By Prop. 3.1.ii the nonzero $`o[[G]]`$-submodule $`N:=VM`$ is closed in $`M`$ and hence lies in $`\mathrm{Mod}_{\mathrm{comp}}^{\mathrm{fl}}(o[[G]])`$. Since the quotient $`(M/N)_{\mathrm{cot}}=M/N`$ is nonzero as well it follows from Prop. 1.3 that the kernel of the dual map $`E=M^dN^d`$ is a nonzero proper closed $`G`$-invariant subspace of $`E`$.
One of the typical pathologies of general Banach space representations of $`G`$ is avoided by the admissibility requirement as the following result shows.
Corollary 3.7:
Any nonzero $`G`$-equivariant continuous linear map between two topologically irreducible admissible $`K`$-Banach space representations of $`G`$ is an isomorphism.
Proof: This is immediate from Thm. 3.5 and Cor. 3.6.
The simplest group to which the results of this section apply is the group $`G=ZZ_p`$ of $`p`$-adic integers. As is shown in \[Dia\] already this group has an extreme wealth of topologically irreducible $`K`$-Banach space representations. On the other hand for a commutative group all โreasonableโ topologically irreducible $`K`$-Banach space representations should be finite dimensional. This is achieved by the admissibility requirement. The ring $`o[[ZZ_p]]`$ is the ring considered in classical Iwasawa theory; it is isomorphic to the power series ring $`o[[T]]`$ in one variable over $`o`$ (\[Was\] 7.1). It follows (\[Was\] ยง13.2) that $`K[[G]]`$ is a principal ideal domain in which every maximal ideal is of finite codimension.
Remark: In \[ST2\] we have introduced the notion of an analytic module over the algebra $`D(G,K)`$ of $`K`$-valued distributions on $`G`$ and we have advocated the conjecture that any $`D(G,K)`$-module of finite presentation is analytic. Since $`K[[G]]`$ is naturally a subalgebra of $`D(G,K)`$ base change would (assuming this conjecture) induce a functor from $`\mathrm{Mod}_{\mathrm{fg}}(K[[G]])`$ into the category of analytic $`D(G,K)`$-modules. Since the latter are dual to a certain class of locally analytic $`G`$-representations this functor should correspond to the passage from a $`K`$-Banach space representation to the subspace of locally analytic vectors. The next basic question in this context then would be whether the ring extension $`K[[G]]D(G,K)`$ is faithfully flat. This is in the spirit of whether every admissible $`K`$-Banach space representation of $`G`$ contains a locally analytic vector.
4. The group $`G=GL_2(ZZ_p)`$
In this section we will analyze a certain infinite series of Iwasawa modules for the group $`G:=GL_2(ZZ_p)`$. Let $`BG`$ denote the Iwahori subgroup of all matrices which are lower triangular modulo $`p`$. In $`B`$ we consider the subgroups $`P,P^{}`$, and $`T`$ of lower triangular, upper triangular, and diagonal matrices, respectively. We also need the subgroups $`U`$ and $`U^{}`$ of unipotent matrices in $`P`$ and $`P^{}`$, respectively. We fix a continuous character $`\chi :To^\times `$. By Cor. 2.2 it extends uniquely to a continuous homomorphism of $`K`$-algebras $`\chi :K[[T]]K`$. The inclusions $`PBG`$, resp. the projection $`PT`$, induce continuous algebra monomorphisms $`K[[P]]K[[B]]K[[G]]`$, resp. a continuous algebra epimorphism $`K[[P]]K[[T]]`$. We denote by $`K^{(\chi )}`$ the one dimensional $`K[[P]]`$-module given by the composed homomorphism $`K[[P]]K[[T]]\stackrel{\chi }{}K`$. Our aim is to study the finitely generated $`K[[G]]`$\- and $`K[[B]]`$-modules
$$M_\chi :=K[[G]]\underset{K[[P]]}{}K^{(\chi )}\text{and}N_\chi :=K[[B]]\underset{K[[P]]}{}K^{(\chi )},$$
respectively. In a similar way (and by a slight abuse of notation) we have the finitely generated $`K[[B]]`$-module
$$N_\chi ^{}:=K[[B]]\underset{K[[P^{}]]}{}K^{(\chi )}.$$
Put $`w:=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)G`$ and $`w\chi (t):=\chi (w^1tw)`$. As a consequence of the Bruhat decomposition $`G=B\dot{}BwP`$ the module $`M_\chi `$, as a $`K[[B]]`$-module, decomposes into
$$M_\chi N_\chi N_{w\chi }^{}.$$
For later use we note that this decomposition is not $`K[[G]]`$-equivariant since obviously $`wN_\chi N_{w\chi }^{}`$.
The module theoretic properties of the series of modules $`N_\chi `$ and $`M_\chi `$ are governed by one numerical invariant $`c(\chi )K`$ of the character $`\chi `$ which is defined by the expansion
$$\chi (\left(\begin{array}{cc}a^1& 0\\ 0& a\end{array}\right))=\mathrm{exp}(c(\chi )\mathrm{log}(a))$$
for $`a`$ sufficiently close to 1 (the existence follows from the topological cyclicity of the group $`1+pZZ_p`$).
In order to investigate the module $`N_\chi `$ we use the Iwahori decomposition which says that multiplication induces a homeomorphism $`U^{}\times P\stackrel{}{}B`$. It implies that $`o[[B]]=o[[U^{}]]\widehat{}o[[P]]`$ where $`\widehat{}`$ is the completed tensor product for linear-topological $`o`$-modules (\[SGA3\] $`\mathrm{Exp}.\mathrm{VII}_B`$ (0.3)). The inclusion $`K[[U^{}]]K[[B]]`$ therefore induces an isomorphism of $`K[[U^{}]]`$-modules
$$K[[U^{}]]\stackrel{}{}N_\chi .$$
$`()`$
In particular, any $`K[[B]]`$-submodule of $`N_\chi `$ corresponds to a certain ideal in the ring $`K[[U^{}]]`$. Since the matrix $`\gamma :=\left(\begin{array}{cc}1& p\\ 0& 1\end{array}\right)`$ is a topological generator of $`U^{}`$ the ring $`K[[U^{}]]`$ is the ring of formal power series in $`\gamma 1`$ whose coefficients are bounded. As already recalled earlier this is a principal ideal domain and each ideal is generated by a polynomial in $`\gamma 1`$ with all its zeros lying in the open unit disk (\[Was\] ยง7.1).
Proposition 4.1:
If $`c(\chi )IN_0`$ then $`N_\chi `$ is a simple $`K[[B]]`$-module.
Proof: Let $`NN_\chi `$ be a nonzero $`K[[B]]`$-submodule, $`IK[[U^{}]]`$ be the ideal which corresponds to $`N`$ under the above isomorphism $`()`$, and $`F_I(\gamma 1)`$ be a polynomial which generates $`I`$ and has all its zeros in the open unit disk. The action of the element $`t_a:=\left(\begin{array}{cc}a& 0\\ 0& 1\end{array}\right)T`$ on $`N_\chi `$ is given on the left hand side of $`()`$ by
$$F(\gamma 1)\chi (t_a)F(\gamma ^a1)$$
for any bounded power series $`F(x)`$. Using the bounded power series
$$\omega _a(x):=(x+1)^a1=\underset{nIN}{}\left(\genfrac{}{}{0pt}{}{a}{n}\right)x^n$$
this can be rewritten as
$$F(\gamma 1)\chi (t_a)F(\omega _a(\gamma 1)).$$
Since this action preserves the ideal $`I`$ it follows that with $`z`$ every $`\omega _a(z)`$, for $`aZZ_{p}^{}{}_{}{}^{\times }`$, is a zero of the polynomial $`F_I(x)`$. This is only possible if $`z+1`$ is a $`p^m`$-th root of unity for some $`mIN`$. We therefore see that there are natural numbers $`k_\mathrm{o}`$ and $`\mathrm{}`$ such that $`F_I(x)`$ divides $`\omega _{p^{k_\mathrm{o}}}(x)^{\mathrm{}}`$. In particular, for any natural number $`kk_\mathrm{o}`$, the polynomial $`\omega _{p^k}(x)^{\mathrm{}}`$ lies in $`I`$. We now look at the action of the element $`u:=\left(\begin{array}{cc}1& 0\\ 1& 1\end{array}\right)`$ on $`N_\chi `$. It is straightforward to check that on the left hand side of $`()`$ we have
$$u(\gamma ^n)=\chi (\left(\begin{array}{cc}(1+np)^1& 0\\ 0& 1+np\end{array}\right))\gamma ^{n/(1+np)}\text{for any}nIN_0.$$
It follows that, for $`kk_\mathrm{o}`$, with $`\omega _{p^k}(x)^{\mathrm{}}`$ also
$$u((\gamma ^{p^k}1)^{\mathrm{}})=\underset{j=0}{\overset{\mathrm{}}{}}(1)^j\left(\begin{array}{c}\mathrm{}\\ j\end{array}\right)\chi (\left(\begin{array}{cc}(1+jp^{k+1})^1& 0\\ 0& 1+jp^{k+1}\end{array}\right))\gamma ^{jp^k/(1+jp^{k+1})}$$
lies in the ideal $`I`$. If $`\omega _{p^k}(x)^{\mathrm{}}`$ and its image under $`u`$, for some $`kk_\mathrm{o}`$, have no zero in common then $`I`$ has to be the unit ideal which means that $`N=N_\chi `$. In the opposite case we obtain
$$\begin{array}{ccc}0& =& \underset{j=0}{\overset{\mathrm{}}{}}(1)^j\left(\begin{array}{c}\mathrm{}\\ j\end{array}\right)\chi (\left(\begin{array}{cc}(1+jp^{k+1})^1& 0\\ 0& 1+jp^{k+1}\end{array}\right))\\ \\ \\ & =& \underset{j=0}{\overset{\mathrm{}}{}}(1)^j\left(\begin{array}{c}\mathrm{}\\ j\end{array}\right)\mathrm{exp}(c(\chi )\mathrm{log}(1+jp^{k+1}))\end{array}$$
for any sufficiently big $`kk_\mathrm{o}`$. This implies that the function
$$\underset{j=0}{\overset{\mathrm{}}{}}(1)^j\left(\begin{array}{c}\mathrm{}\\ j\end{array}\right)\mathrm{exp}(c(\chi )\mathrm{log}(1+jy))$$
of the variable $`y`$ which is analytic in a sufficiently small open disk around zero has infinitely many zeros and hence vanishes identically. In order to prove our assertion we therefore have to show that this is only possible if $`c(\chi )IN_0`$. But if $`c(\chi )IN_0`$ then evaluating all higher derivatives of the above function in zero would lead to the identities
$$\underset{j=1}{\overset{\mathrm{}}{}}(1)^j\left(\begin{array}{c}\mathrm{}\\ j\end{array}\right)j^m=0\text{for any}mIN.$$
This is clearly impossible.
The proof of the following companion result is completely analogous and is therefore omitted.
Proposition 4.2:
If $`c(\chi )IN_0`$ then $`N_\chi ^{}`$ is a simple $`K[[B]]`$-module.
Lemma 4.3:
$`\mathrm{Hom}_{K[[B]]}(N_\chi ^{},N_\chi ^{})=\mathrm{Hom}_{K[[B]]}(N_\chi ^{},N_\chi ^{})=0`$ for any two continuous characters $`\chi `$ and $`\chi ^{}:To^\times `$.
Proof: We compute
$$\begin{array}{c}\mathrm{Hom}_{K[[B]]}(N_\chi ^{},N_\chi ^{})=\mathrm{Hom}_{K[[P]]}(K^{(\chi ^{})},N_\chi ^{})\\ \\ \mathrm{Hom}_{K[[U]]}(K,N_\chi ^{})=\mathrm{Hom}_{K[[U]]}(K,K[[U]])=0.\end{array}$$
The other vanishing follows by a completely symmetric computation.
Theorem 4.4:
If $`c(\chi )IN_0`$ then $`M_\chi `$ is a simple $`K[[G]]`$-module.
Proof: By our above results the decomposition $`M_\chi N_\chi N_{w\chi }^{}`$ is a $`K[[B]]`$-invariant decomposition into two nonisomorphic simple $`K[[B]]`$-modules. But as noted already at the beginning it is not $`K[[G]]`$-invariant. Hence $`M_\chi `$ must be a simple $`K[[G]]`$-module.
The simple $`K[[G]]`$-modules which we have exhibited above are all nonisomorphic as the following result implies.
Proposition 4.5:
We have $`\mathrm{Hom}_{K[[G]]}(M_\chi ^{},M_\chi )=0`$ for any two continuous characters $`\chi \chi ^{}:To^\times `$.
Proof: Because of Lemma 4.3 it is sufficient to show that $`\mathrm{Hom}_{K[[B]]}(N_\chi ^{},N_\chi )=\mathrm{Hom}_{K[[B]]}(N_\chi ^{}^{},N_\chi ^{})=0`$. Since the arguments are completely symmetric we only discuss the vanishing of the first space. Making as usual the identification $`()`$ we have
$$\begin{array}{c}\mathrm{Hom}_{K[[B]]}(N_\chi ^{},N_\chi )=\mathrm{Hom}_{K[[P]]}(K^{(\chi ^{})},N_\chi )\\ \\ =\{FK[[U^{}]]:g(F)=\chi ^{}(g)F\text{for any}gP\}.\end{array}$$
Assume now that there is a nonzero $`F`$ in this latter space. Since any central matrix $`g=\left(\begin{array}{cc}b& 0\\ 0& b\end{array}\right)`$ in $`T`$ acts by multiplication with $`\chi (g)`$ on $`N_\chi `$ it follows immediately that $`\chi `$ and $`\chi ^{}`$ have to coincide on those matrices. On the other hand the action of an element $`t_aT`$ as described in the proof of Prop. 4.1 gives rise to the equation
$$\chi (t_a)F((1+x)^a1)=\chi ^{}(t_a)F(x)$$
between bounded power series over $`K`$. It was shown in the proof of \[ST1\] Prop. 5.5 that this implies $`c(\chi ^{})c(\chi )2IN_0`$ and $`F(x)=[\mathrm{log}(1+x)]^{(c(\chi ^{})c(\chi ))/2}`$. Since the power series $`\mathrm{log}(1+x)`$ is not bounded we in fact obtain $`c(\chi ^{})=c(\chi )`$ and $`F(x)=1`$. Going back to the above equation it follows that $`\chi (t_a)=\chi ^{}(t_a)`$. Hence we have shown that the existence of a nonzero $`F`$ forces the characters $`\chi `$ and $`\chi ^{}`$ to coincide.
To finish we briefly explain the dual picture. In the Banach space $`C(G,K)`$ of all $`K`$-valued continuous functions on $`G`$ we have the closed subspace
$$\mathrm{Ind}_P^G(\chi ):=\{fC(G,K):f(gp)=\chi (p^1)f(g)\text{for any}gG,pP\}.$$
Via the left translation action this is a $`K`$-Banach space representation of $`G`$ (a โprincipal seriesโ representation). By the interpretation of $`K[[G]]`$ as the space of bounded $`K`$-valued measures on $`G`$ we have that $`K[[G]]`$ is the continuous dual of $`C(G,K)`$. It easily follows that
$$\mathrm{Ind}_P^G(\chi )^{}=M_{\chi ^1}.$$
In particular, by Lemma 3.4, $`\mathrm{Ind}_P^G(\chi )`$ is an admissible $`G`$-representation. As a consequence of Cor. 3.6 and Thm. 4.4 we see that $`\mathrm{Ind}_P^G(\chi )`$ is topologically irreducible if $`c(\chi )IN_0`$. This latter fact (for a slightly restricted class of $`\chi `$) was proved in a direct and completely different way in \[Tru\].
References
\[AU\] Arnautov, V.I., Ursul, M.I.: On the uniqueness of topologies for some constructions of rings and modules. Siberian Math. J. 36, 631-644 (1995)
\[BH\] Ballister, P.N., Howson, S.: Note on Nakayamaโs Lemma for Compact $`\mathrm{\Lambda }`$-modules. Asian J. Math. 1, 224-229 (1997)
\[B-CA\] Bourbaki, N.: Commutative Algebra. Paris: Hermann 1972
\[B-GT\] Bourbaki, N.: General Topology. Berlin-Heidelberg-New York: Sprin-ger 1989
\[B-GAL\] Bourbaki, N.: Groupes et algรจbres de Lie, Chap. 2 et 3. Paris: Hermann 1972
\[B-TVS\] Bourbaki, N.: Topological Vector Spaces. Berlin-Heidelberg-New York: Springer 1987
\[CR\] Curtis, C., Reiner, I.: Representation theory of finite groups and associative algebras. New York-London: Wiley 1962
\[SGA3\] Demazure, M., Grothendieck, A.: Schรฉmas en Groupes I. Lect. Notes Math. 151. Berlin-Heidelberg-New York: Springer 1970
\[Dia\] Diarra, B.: Sur quelques reprรฉsentations $`p`$-adiques de $`ZZ_p`$. Indagationes Math. 41, 481-493 (1979)
\[Laz\] Lazard, M.: Groupes analytiques $`๐ญ`$-adiques. Publ. Math. IHES 26 (1965)
\[Mon\] Monna, A.F.: Analyse non-archimedienne. Ergebnisse der Math. 56. Berlin-Heidelberg-New York: Springer 1970
\[Sch\] Schikhof, W.H.: A perfect duality between $`p`$-adic Banach spaces and compactoids. Indag. Math. 6, 325-339 (1995)
\[NFA\] Schneider, P.: Nichtarchimedische Funktionalanalysis. Course at Mรผn-ster (1997)
\[ST1\] Schneider, P., Teitelbaum, J.: Locally analytic distributions and $`p`$-adic representation theory, with applications to $`GL_2`$. Preprintreihe des SFB 478, Heft 86, Mรผnster 1999
\[ST2\] Schneider, P., Teitelbaum, J.: $`U(๐ค)`$-finite locally analytic representations. Preprint 2000
\[Ser\] Serre, J.-P.: Local Fields. Berlin-Heidelberg-New York: Springer 1979
\[Tie\] van Tiel, J.: Espaces localement $`K`$-convexes I-III. Indagationes Math. 27, 249-258, 259-272, 273-289 (1965)
\[Tru\] Trusov, A.V.: Representations of the groups $`GL(2,ZZ_p)`$ and $`GL(2,\text{ }Q_p)`$ in spaces over non-archimedean fields. Moscow Univ. Math. Bull. 36, 65-69 (1981)
\[Was\] Washington, L.C.: Introduction to Cyclotomic Fields. Berlin-Heidel-berg-New York: Springer 1982
Peter Schneider Mathematisches Institut Westfรคlische Wilhelms-Universitรคt Mรผnster Einsteinstr. 62 D-48149 Mรผnster, Germany pschnei@math.uni-muenster.de http://www.uni-muenster.de/math/u/schneider
Jeremy Teitelbaum Department of Mathematics, Statistics, and Computer Science (M/C 249) University of Illinois at Chicago 851 S. Morgan St. Chicago, IL 60607, USA jeremy@uic.edu http://raphael.math.uic.edu/$``$jeremy
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# Hierarchical organization of cities and nations
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Universality in the behavior of complex systems often reveals itself in the form of scale-invariant distributions that are essentially independent of the details of the microscopic dynamics. A representative paradigm of complex behavior in nature is cooperative evolution. The interaction of individuals gives rise to a wide variety of collective phenomena that strongly differ from individual dynamicsโsuch as demographic evolution, cultural and technological development, and economic activity. A striking example of such cooperative phenomena is the formation of urban aggregates which, in turn, can be considered to cooperate in giving rise to nations. We find that population and area distributions of nations follow an inverse power-law behavior, as is known for cities . The exponents, however, are radically different in the two cases ($`\mu 1`$ for nations, $`\mu 2`$ for cities). We interpret these findings by developing growth models for cities and for nations related to basic properties of partition of the plane. These models allow one to understand the empirical findings without resort to the introduction of complex socio-economic factors.
The way in which urban aggregates are distributed was first investigated by Zipf who, half a century ago, observed that the population distribution of cities follows a power-law behavior with exponent $`\mu 2`$. This Zipf law has a โuniversalโ character since it holds at the world level as well as within a single nation, and the exponent is essentially independent of the area of the nation and its socio-economical conditions. More recently it has been observed that the area distribution of satellite cities, towns and villages around huge urban centers also obeys a power-law with exponent $`\mu 2`$.
The remarkable universal result $`\mu 2`$ has very recently attracted the attention of a number of physicists who model urban growth processes . Makse et al. , using a correlated percolation model in the presence of a density gradient, reproduced the observed morphology of cities and the area distribution of sub-clusters in an urban system. Zanette and Manrubia proposed a stochastic model which generates intermittent spatiotemporal structures, and predicts a population distribution in agreement with that observed empirically. Very recently Marsili and Zhang proposed a model based on a master equation approach which is able to give a population distribution close to that obeyed by cities. The transition probabilities which enter the master equation were related to some estimate of the interactions among individuals living in the same city. Though the concept of interaction among human beings is not so clearly defined as that among particles, it is nevertheless true that individuals living in the same city are related to each other by a number of โlinksโ which in the end define the very concept of city.
People living within the same nation are also related to each otherโe.g., they share language and cultural heritage. Since the average โinteractionsโ among the inhabitants of a nation may differ from those among people living in the same city, one could ask if the distribution of nations obeys the same law as that of cities. In order to answer to this question we analysed the population and area distributions of the worldโs nations. The log-log plot of the population distribution $`f(P)`$ is shown in Fig. Hierarchical organization of cities and nationsa for all nations of the world. In the same figure we show also the population distribution for the 140 largest city agglomerates of the USA. Both distributions obey a power law dependence, $`f(P)P^{\mu _P}`$. Regression fits give $`\mu _P=0.97`$ for nations, and $`\mu _P=1.94`$ for cities. Figure Hierarchical organization of cities and nationsb shows a log-log plot of the area distribution $`f(A)`$ for all nations of the world as well as for nations belonging to a subset of all nations (Europe). In both cases $`f(A)`$ obeys a power law dependence, $`f(A)A^{\mu _A}`$ with exponent close to unity.
The strikingly different values of the exponent for cities and nations suggest fundamental differences in the historical and social processes that lead to such distributions. Here we propose a simple explanation of the distributions of nations and cities based only on geometrical considerations. We observe that, at least in principle, there is no restriction to the land accessible to a nation except that of the total existing land: a sufficiently powerful nation could expand to absorb all other nations. Cities, being the result of spontaneous aggregation of individuals around sites having attractive features, can form away from existing ones. This separates the plane region into land basins, and each new city spans a single basin. The resulting distribution of areas is not destroyed when a city expands to absorb nearby cities and gives rise to a compact urban aggregate. In fact, unlike nations, cities usually do not lose their land to neighbours: small towns and villages retain their identity and usually become administrative districts of the bigger aggregate (obeying, as shown in Ref. , the same area distribution that holds for separate cities). Thus a major difference in the way according to which land is occupied by cities and nations is that for cities the accessible land is fragmented into basins while it is not for nations.
Next we model the land occupation processes of nations and cities. We first note that the above considerations suggest that geometric properties may be behind the differences between city and nation distributions. Hence we analyse these processes as random partitionings of the plane. Since different nations (or cities) do not all form at the same time, we consider partition processes in which the different portions are sequentially selected .
One of the simplest ways of partitioning a plane is to divide it using straight lines that are randomly oriented and positioned. Each line divides the region into two portions, of which the smaller is selected and the larger is further partitioned. This process is close in spirit to the way in which land is occupied by nations. Because this partition model resembles the positioning of fences, we refer to it as fence model.
This new model can be solved analytically. After $`n`$ partitions, the land โavailableโ for further division is $`A_n`$, and $`A_n=r_nA_{n1}`$, where $`r_n`$ is a random factor uniformly distributed between 1/2 and 1. The area $`T_n`$ of the โtakenโ portion is $`T_n=(1r_n)A_{n1}`$, and its logarithm can be written as $`\mathrm{ln}T_n=_{i=1}^{n1}\mathrm{ln}r_i+\mathrm{ln}(1r_n)+\mathrm{ln}A_0`$. If we plot on the $`x`$ axis the values of $`\mathrm{ln}T_n`$, we get random points, the average distance between which is equal to $`\mathrm{ln}r_i=1\mathrm{ln}20.3`$. The largest nation corresponds to almost an entire continent $`A_0/2\mathrm{}`$, and the smallest one corresponds to $`A_0(1/2\mathrm{})e^{0.3n}`$, where $`n`$ is the number of nations on the continent, and the factor $`2\mathrm{}`$ corresponds to the average $`\mathrm{ln}(1r_4)1.7`$. Thus, the distribution of the logarithm of a nationโs area in each continent is a flat distribution between $`\mathrm{ln}A_00.3n^{1.7}`$ and $`\mathrm{ln}A_0^{1.7}`$ described by the probability density $`P(\mathrm{ln}S)=1/0.3n`$.
The world has 5 continents; some, such as Australia and North America, have very few nations and some, such as Europe and Africa, have many. Because there are, on average, approximately $`n=50`$ nations per continent, we can expect the approximate distribution of the logarithms of the number of nations to be uniformly distributed between the average largest country $`10^3`$km<sup>2</sup> and $`10^7e^{15}3`$km<sup>2</sup>, which is consistent with a spread observed in the distribution of real nations. The flat distribution of the logarithms $`const`$ $`d(\mathrm{log}S)`$ corresponds to the distribution $`(const/S)`$ $`dS`$ of the areas, which is close to what we observe in Fig. 1b. Note, however, that, due to the few nations in each bin, a particular realization of the partition process described above may significantly deviate from the flat distribution we would expect.
The above model could be, however, oversimplified. This concern can be alleviated by incorporating into it the possibility that nations can evolve, growing or shrinking. Consider, e. g., the variant of the fence model (called, in the following, evolution model) in which with probability $`1/2`$ a nation can grow or shrink by some given amount. The $`P(A)`$ histogram, remarkably, is not affected. To see this, we first note that such change in area corresponds to the variable $`\mathrm{log}(A)`$ increasing or decreasing by a constant number, i.e., a simple random walk in the variable $`\mathrm{log}(A)`$. As we mention above, $`P(A)1/A`$ is equivalent to $`P[\mathrm{log}(A)]const`$.
In the present case, the random walk is confined by reflecting boundaries: $`A_{\text{max}}=A_0`$ is some minimal size $`A_{\text{min}}`$ below which the nation is not stable. The distribution of a random walk confined in an interval with reflecting boundaries converges to a uniform distribution . Thus $`P[\mathrm{log}(A)]`$ converges to a constant and hence $`P(A)1/A`$. The the $`P(A)`$ distribution is immune to the โnoiseโ of growth and shrinking. The distribution of nation areas, simulated using the evolution model, is shown in Fig. 1c. The distribution of the logarithms of the nation areas are shown in Fig. 1d.
To model city distributions, we shall consider a radically different way of partitioning the plane. As suggested by the way urban geographers have thought about โcentral placeโ theory and the hierarchy of towns , we assume that cities have on the average a circular shape and thus can be approximately represented as circles. Land occupation by cities can thus be modeled through the partition of the plane in nonoverlapping circles. Each new portion is a circle (with radius chosen from a uniform distribution) at a randomly-chosen position, but outside of previously-selected circles. The maximal area of each new circle is limited by the distance from the closest existing circle. The resulting fragmentation of accessible land reduces the space available for the next circle more rapidly with respect to what happens for the portions generated through the fence model, thus making small portions much more probable. Hence we expect $`\mu _A`$ to be larger for the present model.
We simulate the circle model and find that the distribution of the circle areas follows a power-law behavior which is close to the empirical data of Fig.1a. A linear fit of the distribution (with the exclusion of the region affected by finite size effects) gives an exponent $`\mu 1.94`$ (Fig.1e). One advantage of the circle model, compared to the previous models , is that it is based only on the geometrical features of the land occupation process.
Another advantage of this model is that it can be solved analytically in the limit when the area of the newly-formed circle is a small number proportional to the area of the unoccupied land closest to the center of this circle, with proportionality coefficient $`k1`$. Suppose that there are $`N`$ circles in a region of total area $`A_0`$. Then the probability that a randomly chosen point is surrounded by an empty space of area larger than $`X`$ is governed by a Poisson distribution $`F(X,N)=\mathrm{exp}(XN/A_0)`$. The number of circles with area between $`A`$ and $`A+dA`$ is $`p(A,N)dA=N/(A_0k)\mathrm{exp}(AN/A_0k)dA`$, the derivative of the Poisson distribution function. The distribution of circles, $`P(A)`$, at some given instant when there are $`N_c`$ circles is given by the integral $`_0^{N_c}p(A,N)๐N=(kA_0/A^2)[1\mathrm{exp}(AN_c/A_0k)]`$. For large $`N_c`$, we recover $`P(A)A^2`$, so $`\mu _A=2`$, in agreement with empirical findings. It is interesting to note that the exponent $`\mu _A=2`$ is robust since, as our simulations show, it holds even in the case when $`k`$ is not a small number, but is any random number constrained only by the fact that two circles cannot overlap (as for urban agglomerates that consist of tightly-packed villages and towns).
In summary, we show, through the analysis of accurate geographical and demographical data, that the nation population and nation area distributions obey power-laws. The exponent is, however, surprising and completely unexpected. Moreover, its value (1) is quite different from the known exponent for cities (2). In addition to our empirical discovery, we propose an explanation for why the population distribution exponent of cities differs so strikingly from that of nations.
We thank J. S. Andrade Jr., M. Batty, R. S. Dokholyan, I. Grosse, P. L. Krapivsky, H. A. Makse, M. Morrissey, and S. Redner for interesting and stimulating discussions. We thank NSF for support. N. V. D. is supported by NIH postdoctoral fellowship (GM20251-01).
FIG.1. (a) Double-logarithmic plot of the histogram of the population, $`P`$, of the 255 world nations and 140 largest urban agglomerates in the USA in mid-1997. For nations, the slope gives $`\mu _P=0.97`$ and the linear regression coefficient is $`R=0.97`$; for cities $`\mu _P=1.94`$ and $`R=0.99`$. For visual convenience, city data are multiplied by $`10^2`$. (b) Double-logarithmic plot of the histogram of areas, $`A`$, of the 255 nations of the world and 49 European nations; the corresponding coefficients are $`R=0.99`$ and 0.99 respectively. European data are divided by $`10^2`$. The source for both plots is http://www.stats.demon.nl. (c) Double-logarithmic plot of the histogram of areas for the evolution model. (d) Distributions of logarithms of nations areas produced by (i) 255 nations and (ii) by the evolution model. (e) Double-logarithmic plot of the histogram of areas for the circle model; the linear regression coefficient is $`R=0.96`$.
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# Metastable States of a Coupled Pair on a Repulsive Barrier
## Abstract
Resonance penetration of two coupled particles through a repulsive barrier is considered. It is shown that a local minimum of the total potential generates metastable bound states, and their spectrum determines the position of resonances in the penetration probability. It is pointed out that the probabilities of tunneling of two interacting particles from the false vacuum can be essentially higher than it has been assumed earlier.
PACS number(s): 03.65.Nk, 11.10.Jj, 21.45.+v
In paper by N.Saito and Y.Kayanuma , it was pointed out that there exists a new quantum phenomenon โ resonance transparency of a single repulsive barrier for a coupled pair of particles. To consider this effect, a one-dimensional rectangular repulsive barrier and an infinite one-dimensional rectangular potential well coupling the pair were chosen. Since the interactions were simple, it was possible to solve the initial two-dimensional Schroedinger equation by reducing it to a system of one-dimensional equations by means of projection onto 7 eigenfunctions. However, there still remained the question of how this effect manifests itself in other systems.
This note continues the study of the effect of resonance transparency of two type one-dimensional barriers for a pair of identical particles coupled by the oscillatory interaction. This pair interaction allows us to reduce the problem of three-dimensional scattering of a three-dimensional oscillator to the solution of a two-dimensional equation analogous to the equation derived in ref. . Besides, it is just this sort of pair interaction that is used in literature devoted to the probability of induced decay of the false vacuum in collisions of high-energy particles (see, for instance, ). It was pointed out there that it is possible to describe the processes of transition from the false vacuum on the basis of quantum-mechanical tunneling of a pair of particles through the barrier; but the study was performed for a system where only one of the oscillator particles interacts with the barrier. Here, we show that, when the two particles interact with the barrier, there arises the same effect of resonance transparency as in ref. .
The first potential barrier we study is taken in the Gauss form from ref. in order to show that it is possible to drastically increase the probability of induced decay of the false vacuum. The second potential barrier of the Coulomb form is investigated in order to draw attention to the fact that the resonance tunneling of the barrier is feasible in the problems of fusion of heavy nuclei. The method of investigation is based on the numerical solution of the two-dimensional Schroedinger equation without any further simplifications.
Consider the penetration of a pair of identical particles with masses $`m_1=m_2=m`$, and coordinates $`๐ซ_1`$ and $`๐ซ_\mathrm{๐}`$ coupled by an oscillatory interaction through the potential barrier $`V_0(x_1)+V_0(x_2)`$. The Hamiltonian of this system ($`\mathrm{}`$ = 1)
$$\frac{1}{4m}\mathrm{}_R\frac{1}{m}\mathrm{}_r+\frac{m\omega ^2}{4}r^2+V_0(๐๐ซ/\mathrm{๐})+V_0(๐+๐ซ/\mathrm{๐}),$$
written in coordinates of the center of inertia of the pair $`๐=(๐ซ_1+๐ซ_2)/2`$ and in an internal coordinate of the relative motion $`๐ซ=๐ซ_1๐ซ_2`$ describes the three-dimensional motion of a three-dimensional oscillator with the frequency of vibrations $`\omega `$. Since the potential barrier depends only on one variable, and the oscillatory interaction is additive in $`๐ซ`$ projections, the wave function is factorized, and its nontrivial part describing the process of scattering depends only on two variables. It is convenient to represent these variables in the dimensionless form
$$x=\sqrt{\frac{m\omega }{2}}(x_1x_2),y=\sqrt{\frac{m\omega }{2}}(x_1+x_2).$$
The Schroedinger equation in these variables is of the form
$$\left(_x^2_y^2+x^2+V(xy)+V(x+y)E\right)\mathrm{\Psi }=0,$$
(1)
where the energy $`E`$ is written in units $`\omega /2`$, and the potential barrier $`V(x\pm y)=\frac{2}{\omega }V_0((x\pm y)/\sqrt{2m\omega })`$ in what follows will be written in a form convenient for us. Equation (1) should be supplemented with boundary conditions. Let the process of scattering proceed from left to right, and the initial state of the oscillator be the state $`n`$. Then the boundary conditions are written in the form
$$\begin{array}{cc}\mathrm{\Psi }|_y\mathrm{}=\hfill & \hfill \mathrm{exp}(ik_ny)\phi _n(x)\underset{jN}{}S_{nj}\mathrm{exp}(ik_jy)\phi _j(x);\\ \mathrm{\Psi }|_{y+\mathrm{}}=\hfill & \hfill \underset{jN}{}R_{nj}\mathrm{exp}(ik_jy)\phi _j(x);\\ \mathrm{\Psi }|_{x\pm \mathrm{}}=\hfill & \hfill 0.\end{array}$$
(2)
The wave functions of the oscillator $`\phi _j(x)`$ obey the Schroedinger equation
$$\left(_x^2+x^2\epsilon _i\right)\phi _i=0$$
(3)
with the energy $`\epsilon _j=2j+1`$ ($`j=0,1,2,\mathrm{}`$), momenta $`k_j=\sqrt{E\epsilon _j}`$, and the number $`N`$ of the last open channel ($`E\epsilon _{N+1}<0`$). In what follows, we consider an oscillator composed of bosons, whose spectrum is convenient to number from 1. So, hereafter, $`\epsilon _j=4j3`$ ($`j=1,2,\mathrm{}`$).
The probabilities of penetration $`W_{ij}`$ and reflection $`D_{ij}`$ are defined as the ratio of the density of a penetrated or a reflected flux to the density of the flux of incident particles:
$$W_{ij}=|R_{ij}|^2\frac{k_j}{k_i},D_{ij}=|S_{ij}|^2\frac{k_j}{k_i}.$$
It is obvious that $`_{jN}\left(W_{ij}+D_{ij}\right)=1`$.
To determine the probabilities of penetration (reflection) in the above way, it is necessary to solve the two-dimensional differential equation (1); its numerical solutions will be given below. The numerical solution was based on a three-diagonal approximation of second derivatives with constant steps in $`x`$ and $`y`$: $`h_x=0.025`$, $`h_y=0.005`$, respectively. Finite dimensions $`|y_{\mathrm{max}}|=12`$ and $`|x_{\mathrm{max}}|=7`$ of the range of numerical calculations, at a given degree of discretization, provided accuracy to within the third decimal point in all the presented calculations.
As it was indicated above, in papers devoted to the induced decay of false vacuum, use was made of the model of quantum-mechanical tunnelling of a pair coupled by the oscillatory interaction through a barrier . The case when only one particle interacts with the barrier was considered. In the framework of the resonance tunnelling, we can expect that the picture of tunnelling would essentially change when the interaction of both the incident particles with the barrier is switched on. The barrier used in ref. , upon being made dimensionless, is of the form
$$V(X)=\frac{2}{g^2\omega }\mathrm{exp}(g^2X^2/\omega );X=x\pm y,$$
(4)
where $`\omega `$ is the oscillator frequency, and $`g^21`$ is the model parameter of false vacuum. In ref. , the dependence of the tunneling probability was calculated in the range of $`g^2`$ from 0.09 till 0.01 and at fixed frequency $`\omega =1/2`$. In the present calculations, the same value of $`\omega `$ but greater values of $`g^2`$ are accepted. The reason is that there arise extremely (in the framework of numerical calculations) narrow resonances in the energy dependence of the tunneling probability of a coupled pair through the barrier. Therefore, in Fig.1, we present the results of numerical calculations of equation(1) with barrier (4) at $`g^2=`$ 0.5, 0.3, 0.2 denoted by letters A, B, C, respectively.
In this Fig., we draw the probabilities of penetration of a coupled pair from the ground state into all the possible states, i.e. $`W=_{jN}W_{1j}`$ at different values of $`g^2`$. A prominent resonance dependence of the tunneling probability shows that there do exist barrier resonances under discussion. Note that the first resonance with decreasing $`g^2`$ shifts towards higher energies, but the quantity $`g^2E_r`$ diminishes.
When $`g^2=0.2`$, the probability of resonance penetration is many times ($`10^8`$) as large as the probability of penetration in the nonresonance region (background). Therefore, in Fig.1, we present the results of calculations on the logarithmic scale. They demonstrate not only the indicated exceeding but also the coincidence of the background part of the curve with the probability of penetration of a structureless particle, i.e. with the solution of equation (1) for the barrier potential $`2V(y)`$. To estimate the contribution of narrow resonances to the probability of penetration of the particle flux distributed over energy, in Fig.1, we plot the penetrated flux
$$I(E)=\frac{1}{EE_0}_{E_0}^EW(E^{})๐E^{},$$
in the case when the incident flux is distributed uniformly from $`E_0`$ till $`E`$; the quantity $`E_0=5`$. It is seen that the main contribution to the probability of penetration comes from resonances. At $`E=23`$, the difference from the background penetration amounts to 4 orders.
Now, we describe a simple scheme of arising barrier metastable states that make the barrier transparent. It is not difficult to verify that the potential energy $`U(x,y)=V(x+y)+V(xy)+x^2`$ possesses a local minimum at $`y=0`$ (the center of mass in the middle of the barrier) and at some values of $`x=\pm x_0`$. A maximum is at $`x=0`$.
This, there exist two potential โwellsโ separated by the barrier. Bound states of that system split into even and odd states. The magnitude of splitting is determined by the probability of penetration through the internal barrier. When $`2V(x=0,y=0)1`$, this shift can be very small, and the spectrum of even states is determined by the spectrum of an isolated โwellโ. In the first approximation, the position of resonances can be described by the oscillator spectrum of bound states at $`y=0`$ and $`x=x_0`$
$$E_{n_xn_y}=E_0+2\omega _x(1/2+n_x)+2\omega _y(1/2+n_y),$$
(5)
where $`n_x`$ and $`n_y`$ are oscillator quantum numbers; $`E_0=2V(x_0)`$; and frequencies $`\omega _x`$ and $`\omega _y`$ are determined by the second derivatives at the point of local minimum: $`\omega _x=\sqrt{_x^2U(x,y)/2}`$, $`\omega _y=\sqrt{_y^2U(x,y)/2}`$.
For potential (4), it is not difficult to obtain the oscillator-model parameters
$$\begin{array}{c}E_0=\omega (1+2\mathrm{ln}(2/\omega ))/g^2;\\ \omega _x^2=4\mathrm{ln}(2/\omega );\\ \omega _y^2=\omega _x^21.\end{array}$$
(6)
Let us compare the positions of resonances drawn in Fig.1 at $`g^2=0.2`$ with the results of calculation by formulae (5) and (6) from which it is clear that there exist small but clearly seen satellite resonances. The position of the first resonance $`E_r=13.65`$ is well described by the oscillator energy in the ground state $`E_{00}=13.92`$. The second group of resonances is generated by a single excitation of oscillators either along $`y`$ or along $`x`$: $`E_{01}=18.18`$, $`E_{10}=18.63`$. They are associated with the resonances at energies $`17.24`$ and $`17.74`$. The third group of resonances is generated by a double excitation $`E_{02}=22.45`$, $`E_{11}=22.90`$, $`E_{20}=23.34`$, and respectively, resonance energies are $`20.58`$, $`20.88`$, $`21.72`$. So, the simple oscillator model of a metastable barrier state gives a correct qualitative picture of the origin of resonances. Comparison with Fig.1 shows that the largest values of the tunneling probability correspond to metastable states with a minimal excitation along the coordinate of the center of inertia.
Decomposition around the point of equilibrium does not exhaust all possibilities of the oscillator model. Agreement between the resonance energy and the energy of a metastable state can be improved by a simple variational procedure. To this end, we consider the position of a minimum $`x_0`$ and frequencies $`\omega _x`$ and $`\omega _y`$ to be unknown quantities that are determined by the minimum of the average of the total Hamiltonian
$$\overline{H}=\varphi _x\varphi _y|_x^2_y^2+x^2+V(xy)+V(x+y)|\varphi _y\varphi _x$$
over normalized eigenfunctions of the oscillator in the ground state:
$$\varphi _x=\left(\frac{\omega _x}{2\pi }\right)^{1/4}\mathrm{exp}\left(\frac{\omega _x}{4}(xx_0)^2\right);$$
$$\varphi _y=\left(\frac{\omega _y}{2\pi }\right)^{1/4}\mathrm{exp}\left(\frac{\omega _y}{4}y^2\right).$$
Varying $`\overline{H}`$ over $`x_0`$, $`\omega _x`$ and $`\omega _y`$, we can derive a system of three nonlinear equations to be not presented here in view of their being cumbersome. The of the three equations are solved analytically: $`x_0=x_0(\omega _x,\omega _y)`$, $`\omega _y=\sqrt{\omega _x^21}`$. In this way, the function of one variable $`\overline{H}=\overline{H}(\omega _y)`$ is to be estimated numerically; its minimum determines the variational estimate $`E_{\mathrm{var}}`$ for the first resonance. Note that the variational connection of $`\omega _x`$ and $`\omega _y`$ is the same as in the case of decomposition (6). In Table I, we present the comparison of $`E_{\mathrm{var}}`$ with positions of the first resonance $`E_r`$ drawn in Fig.1. Agreement can be considered good in the framework of the above-indicated accuracy. TABLE I.: Comparison of positions of the first resonance with the variational estimate $`g^2`$ $`E_{\mathrm{var}}`$ $`E_r`$ $`g^2E_r\omega /2`$ 0.5 7.649 7.62 1.30 0.3 10.416 10.38 0.779 0.2 13.680 13.65 0.683 When $`g^21`$, the variational expressions get simplified and allow the following decomposition:
$$E_{\mathrm{var}}^{\mathrm{as}}E_{00}+O(g^2).$$
It coincides, with an accuracy up to $`O(g^2)`$, with the energy derived by a simple decomposition around the minimum of $`U(x,y)`$. So, the estimate of the resonance spectrum made by formulae (5) and (6) is asymptotic as $`g^20`$. In particular, when $`g^20`$, we can indicate the limiting position of the first resonance in units of $`g^2`$, i.e. the quantity $`g^2E\omega /2`$ used in ref. :
$$g^2E\omega /2\omega ^2(1+2\mathrm{ln}(2/\omega ))/2.$$
At $`\omega =1/2`$, this energy tends to $`0.472`$; this is shown in the fourth column of Table I. In ref. where the study was made of the penetration of a pair through the barrier, a smooth curve was obtained for the probability of penetration in the energy interval from $`1.2`$ to $`2`$. From the presented calculations it follows that, if the interaction of both particles with the potential barrier is taken into account, this curve becomes essentially nonmonotone.
The barriers considered above are of the form of the Gauss function. For completeness, below we present the calculations for a barrier of the Coulomb shape cutoff both at short and long distances:
$$V(X)=\{\begin{array}{cc}Q/X_{\mathrm{min}}& :|X|<X_{\mathrm{min}}\hfill \\ Q/|X|& :X_{\mathrm{min}}|X|X_{\mathrm{max}}\hfill \\ Q/X_{\mathrm{max}}& :|X|>X_{\mathrm{max}}\hfill \end{array};X=x\pm y.$$
(7)
The cutoff at short distances was introduced for modelling the nuclear Coulomb barrier in the framework of constraints imposed by the one-dimensional scattering. With this cutoff, the notion of โbarrier heightโ is meaningful for the one-dimensional model of scattering. For a greater analogy, the barrier width at $`|X|=X_{\mathrm{min}}`$ should be small in spatial units of the problem, i.e., as compared with the mean-square dimension of the oscillator. The cutoff at long distances was introduced to make it possible to use an asymptotics of the type (2). The quantity $`X_{\mathrm{max}}`$ should be larger than 1 for imitating the barrier of small transparency. Here we took $`X_{\mathrm{min}}=0.1`$ $`X_{\mathrm{max}}=5`$. The quantity $`Q`$ determines the energy height of the barrier. In Fig.2, we show the results of calculations for $`Q=`$ 2, 4, 10 denoted by letters A, B, and C, respectively.
In this case, a clear picture of the resonance tunneling of a coupled pair is also observed. We do not present the analysis of the oscillator model for the position of resonances, because the chosen potential is essentially of the model character. We only mention that satellite resonances manifest themselves clearly, and energies of principal resonances are equidistant.
The considered mechanism of the transparency of barriers for a coupled pair manifests itself for all the potential barriers chosen for the investigation. As the resonance transparency was first observed for barriers of the rectangular shape, and the coupling in a pair was of nonoscillator type, it could be assumed that the resonance transparency of barriers for composite particles could be observed for a wide class of interactions. Therefore, the effects of quantum transparency could occur in various fields of physics. In particular, when the interaction of two particles with a barrier is taken into account, the picture of induced decay of the false vacuum changes essentially.
The author expresses his deep gratitude to A.K.Motovilov for the fruitful idea of realization of the numerical scheme and to Yu.M.Chuvilโskii for pointing the importance of the study of the resonance transparency of Coulomb barriers.
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# On the transition to self-gravity in low mass AGN and YSO accretion discs
## 1 Introduction
Viscous discs can acquire large dimensions under the effect of angular momentum redistribution. This is corroborated by the observation of Young Stellar Objects (YSO) which reveals the emission of wide gaseous discs with an outer radius reaching $`1001000`$ AU (Beckwith et al., 1990; Pudritz et al., 1996; Duvert et al., 1998; Guilloteau & Dutrey, 1998; Shepherd & Kurtz, 1999). Also, discs hosted by Active Galactic Nuclei (AGN), although yet unresolved โexcept in the active galaxy NGC4258 (Herrnstein et al., 1999)โ could probably be of much larger size. The existence of strong similarities between YSOs and AGNs, like jets and outflows (Falcke, 1998; Frank, 1998), some excesses in the spectral energy distribution at infrared wavelengths (Adams & Shu, 1986; Zdziarski, 1986; Bertout, 1989; Sanders et al., 1989; Voit 1991; Kenyon, Yi & Hartmann, 1996), a surrounding (dusty) torus (Gรผsten, Chini & Neckel, 1984; Thรฉ & Molster, 1994; Drinkwater, Combes & Wiklind, 1996; Sandqvist, 1999) indicates that the same physical mechanisms should take place in discs, despite the difference in the central mass scale. This constitutes an interesting challenge for the disc theory, especially to understand the transport of angular momentum.
Many models are based on the $`\alpha `$-theory of thin discs (Shakura & Sunyaev, 1973; Pringle, 1981) and assume a vertically averaged structure (e.g., Collin & Dumont 1990; Ruden & Pollack, 1991; Cannizzo & Reiff, 1992; Hurรฉ et al. 1994a; Artemova et al. 1996; Burderi, King & Szuszkiewicz, 1998; Drouart et al. 1999; Aikawa et al., 1999). A physically more satisfactory approach is the investigation of the vertical structure in detail, as done for stars. Such a problem has been discussed by several groups already, in different contexts, with various goals and degrees of sophistication: at the scale of AGNs (Cannizzo, 1992; Wehrse, Stรถrzer & Shaviv, 1993; Dรถrrer et al., 1996; Siemiginowska, Czerny & Kostyunin, 1996; Hubeny & Hubeny, 1998 and references therein; Sincell & Krolik, 1997; Rรณลผaลska et al., 1999), in Cataclysmic Variables (CVs) (Smak, 1984; Mineshige & Osaki, 1983; Meyer & Meyer-Hofmeister, 1982; Pojmaลski, 1986; see Cannizzo, 1993 for a review; Milsom, Chen & Taam; 1994; Dubus et al., 1999) and in YSOs (Malbet & Bertout, 1991; Bell et al., 1997; DโAlessio et al., 1998) including the Primitive Solar Nebula (PSN) (Lin & Papaloizou, 1980; Papaloizou & Terquem, 1999). Basic considerations show that the outermost regions of (low mass) discs are expected to be regulated by self-gravity (Goldreich & Lynden-Bell, 1965; Shlosman & Begelman, 1987). This may be the case of AGN and YSO discs. A few constraints on the disc thickness, mass and rotation motion at large radii are set by observations (Guilloteau & Dutrey, 1998; Mundy, Looney & Welch, 2000; Herrnstein et al., 1999). So, it appears essential to construct models for as accurate and realistic as possible, despite the lack of knowledge regarding turbulent viscosity which causes accretion. To our knowledge, no self-consistent 2D-model accounting for vertical self-gravity in the outermost parts of discs has been published yet. This is the aim of this paper. The present model includes simultaneously vertical convection, self-gravitation, turbulent pressure, and realistic equation of state (EOS) and opacities. Of particular interest here is the transition from the โclassicalโ disc to the self-gravitating disc which is predicted to be as close as a few $`10^310^4R_{}`$ ($`R_{}`$ is the radius of the central object) by vertically averaged models (Ruden & Pollack, 1991; Dโalessio, Calvet & Hartmann, 1997; Hurรฉ, 1998). The hypothesis of the model and relevant equations are developed in Sect. 2. Two versions of the $`\alpha `$-prescription are tested. The ingredients (EOS and opacities) as well as the numerical method are presented in Sect. 3. We discuss in Sect. 4 a few properties of the model, namely the effect of turbulent pressure and depth dependent viscosity, and the position of the inner edge of the self-gravitating disc, for prototypal circumstellar and AGN discs, for various accretion rates. An Appendix contains a note concerning the treatment of convection and a high precision formula fitting the EOS.
## 2 Model for the vertical structure: hypothesis and relevant equations
### 2.1 General considerations. Accounting for self-gravity
A distinctive feature of steady state keplerian accretion discs is the absence of coupling between the vertical structure and the radial structure (e.g. Frank, King & Raine, 1992). This attractive property is undoubtedly an oversimplification and probably does not match reality. However, it is particularly advantageous from a modeling point of view because any annulus can be treated individually, whatever the state of neighbors, unlike in thick discs where pressure gradients and energy transport in the radial direction gain in importance with respect to their vertical counterparts (Maraschi, Reina & Treves, 1976; Abramowicz, Calvani & Nobili, 1980; Abramowicz et al., 1988; Narayan, Madevan & Quataert, 1998). Note that the keplerian assumption which fixes the rotation law to the value
$$\mathrm{\Omega }=\sqrt{\frac{GM}{R^3}}$$
(1)
where $`M`$ is the central mass and $`R`$ the polar radius, requires that the gas remains confined at altitudes $`z`$ such as $`z^2R^2`$, and small radial pressure gradients too.
Self-gravity may influence and even dominate the equilibrium structure and dynamical evolution of almost any kind of disc, not only in massive or thick discs or tori where effects are global (Bodo & Curir, 1992; Hashimoto, Eriguchi, Mรผller, 1995; Boss, 1996; Laughlin & Rรณลผyczka, 1996, Masuda, Nishida & Eriguchi, 1998), but also in low mass keplerian discs (Paczyลski, 1978; Shlosman & Begelman, 1987) as soon as the mass density of the accreted gas exceeds $`\mathrm{\Omega }^2/4\pi G`$ locally, $`\mathrm{\Omega }`$ being the rotation frequency. In that latter case which is of interest here, self-gravity increases vertical pressure gradients and gathers matter closer to the midplane (Sakimoto & Coroniti, 1981; Shore & White, 1982; Cannizzo & Reiff, 1992; Hurรฉ et al., 1994; Hurรฉ, 1998). So, only the outermost regions of keplerian discs where $`\mathrm{\Omega }`$ reaches low values can be affected by vertical self-gravity, except in very special situations (e.g. Sincell & Krolik, 1997).
Accounting correctly for the disc gravity requires the resolution of the Poisson equation (Hunter, 1963; Stรถrzer, 1993; Bertin & Lodato, 1999)
$$\mathrm{\Delta }\mathrm{\Phi }^{\mathrm{disc}}(R,z)=4\pi G\rho (R,z)$$
(2)
where $`\rho `$ is the gas mass density and $`\mathrm{\Phi }^{\mathrm{disc}}`$ is the gravitational potential due to the bare disc (the $`\varphi `$-invariance is assumed). This is a very difficult task, in particular because the disc surface has a non trivial form which is not known a priori and the deviation from sphericity is extreme. As Eq.(2) connects all annuli together, the method of โindependent ringsโ no longer applies, unless an iterative scheme in which the potential would be step by step improved from the actual density field until convergence (Stahler, 1983). There is however no guarantee that such a scheme effectively converge and the computational time might be prohibitive (Eriguchi & Mรผller, 1991). Here, we follow another, more simple approach: we adopt the infinite and $`R`$-homogeneous slab approximation (Paczyวนski, 1978; Sakimoto & Coroniti, 1981; Shore & White, 1982; Cannizzo & Reiff, 1992; Liu, Xie & Ji, 1994; Hurรฉ, 1998) which yields the gravity due to the disc
$$g_z^{\mathrm{disc}}(R,z)=\frac{\mathrm{\Phi }^{\mathrm{disc}}}{z}4\pi G\mathrm{\Sigma }(R,z),$$
(3)
where the surface density $`\mathrm{\Sigma }`$ is defined by
$$\mathrm{\Sigma }(R,z)=_0^z\rho (R,z^{})๐z^{},$$
(4)
but other assumptions are possible (Mineshige & Umemura, 1997). This enables again to investigate the disc structure annulus by annulus, but introduces a bias: it tends to overestimate (underestimate) self-gravity in regions of high (respectively low) densities. We expect that the present approximation gives acceptable results, at least in regions where the surface density radial gradients remain low with respect to density. Anyway, the error made with respect to a self-consistent model is not known and would be interesting to estimate.
### 2.2 The equations for the disc interior
When the upward transport of heat is treated in the diffusion approximation, the vertical structure of a keplerian disc, possibly irradiated, can be determined from the resolution of a system of four first order coupled ordinary differential equations (ODEs) in between the midplane and the top of the disc (Pojmaลski, 1986; Tuchman, Mineshige & Wheeler, 1990; Cannizzo, 1992; Meyer & Meyer-Hofmeister, 1982; Milsom, Chen & Taam, 1994; Dubus et al., 1999). The complexity of the problem rises when a multi-frequency radiative transfer is performed, as required to compare theoretical spectra with observations and make key predictions (Ross, Fabian & Mineshige, 1992; Wehrse, Stรถrzer & Shaviv, 1993; Dรถrrer et al., 1996; Sincell & Krolik, 1997; Dโalessio et al. 1998; El-Khoury & Wickramasinghe, 1998; Hubeny & Hubeny, 1998; De Kool & Wickramasinghe, 1999). The four equations specify respectively (Frank, King & Raine, 1992):
* the pressure gradient $`_zP`$ which describes the hydrostatic equilibrium of each fictitious slab
$$\frac{1}{\rho }\frac{dP}{dz}=\mathrm{\Omega }^2z4\pi G\mathrm{\Sigma }g_z,$$
(5)
* the heat flux gradient $`_zF`$ due to viscous heating
$$\frac{dF}{dz}=\frac{9}{4}\rho \nu \mathrm{\Omega }^2,$$
(6)
where $`\nu `$ is the $`z`$-dependent viscosity law (see Sect. 2.4),
* the temperature gradient $`_zT`$ which is determined by the heat net flux transported upwards through radiation and convection
$$\frac{dT}{dz}=T\frac{}{\lambda _p},$$
(7)
where $`\lambda _p\frac{dz}{d\mathrm{ln}P}`$ is the pressure height scale and $`\frac{d\mathrm{ln}T}{d\mathrm{ln}P}`$ is the actual gradient (see the Appendix, Sect. A),
* the surface density gradient $`_z\mathrm{\Sigma }`$
$$\frac{d\mathrm{\Sigma }}{dz}=\rho $$
(8)
Note that this last equation is not relevant when self-gravity is left apart (the total surface density of the disc + atmospheres $`\mathrm{\Sigma }_\mathrm{t}=2\times \mathrm{\Sigma }(R,\mathrm{})`$ can be easily computed a posteriori from Eq.(4), once the mass density distribution is known). The above ODEs must be supplemented by a closure relation, an equation of state. For a mixture of radiation and perfect gas undergoing atomic ionization and molecular dissociation at LTE, the total pressure is linked to the density and temperature of matter by
$$P=\frac{\rho kT}{\mu m_\mathrm{H}}+\frac{4\sigma }{3c}T^4$$
(9)
where $`\mu m_\mathrm{H}`$ is the mean mass per particle which depends on the density and temperature (see Sect. 3.1 and the Appendix, Sect. B). The above expression for radiation pressure is compatible with an optically thick disc only.
### 2.3 Accounting for turbulent pressure in the framework of the $`\alpha `$-prescription
Turbulence is the main mechanism driving accretion in discs. It is an extra source of pressure. According to the standard theory of $`\alpha `$-discs (Shakura & Sunyaev, 1973; Pringle, 1981), the typical velocity of turbulent eddies can be taken as $`\sqrt{\alpha }c_\mathrm{s}`$ if one assumes the equipartition between length and velocity turbulent scales ($`c_\mathrm{s}`$ is the adiabatic sound speed). So, the turbulent pressure is
$$p_\mathrm{t}=\alpha \mathrm{\Gamma }_1P$$
(10)
where $`\mathrm{\Gamma }_1=(\frac{d\mathrm{ln}P}{d\mathrm{ln}\rho })_{\mathrm{ad}}`$ is the first adiabatic exponent in the total (gas plus radiation) pressure. The effect of turbulent pressure on the hydrostatic equilibrium is therefore expected if $`\alpha `$ is not too small, as simulations shall confirm. If we include $`p_\mathrm{t}`$ into Eq.(5) and rewrite it in terms of a density gradient equation, we find
$$\frac{d\rho }{dz}=\frac{\rho ^2}{P}\left(\mathrm{\Omega }^2z+4\pi G\mathrm{\Sigma }\right)\frac{1\chi _T}{\chi _\rho \left(1+\alpha \mathrm{\Gamma }_1\right)}$$
(11)
where $`\chi _T`$ and $`\chi _\rho `$ are respectively the temperature and density exponents of the total pressure, and we have assumed $`_z\mathrm{\Gamma }_1=0`$ (this is justified since $`0.9\mathrm{\Gamma }_11.7`$, see for instance the Appendix, Fig.(11)). This expression clearly shows the existence of a density inversion (that is, a zone where $`_z\rho >0`$) each time the actual gradient satisfies $`\chi _T>1`$ (Rรณลผaลska et al., 1999). Such an inversion is likely to occur in radiative pressure dominated layers where $`\chi _T`$ is the largest or/and in convectively unstable zones where $``$ may be large. Note that the inversion still remains if turbulent pressure plays a role, but with a weaker amplitude.
It is likely that turbulent pressure plays a role, not only on the hydrostatic equilibrium as considered here, but also on advection of matter and energy both radially and vertically. We ignore these effects.
### 2.4 On depth-dependent viscosity laws
In vertically averaged disc models, the anomalous viscosity takes locally a mean value defined from midplane quantities, namely $`\nu \nu ^{\mathrm{midplane}}(=\alpha c_\mathrm{s}^2(0)/\mathrm{\Omega }`$). Here, we need to specify how $`\nu `$ varies with the altitude. This point is critical since our knowledge of turbulent viscosity remains extremely limited, if not null. There are however two standard ways to do this within the framework of the $`\alpha `$-prescription. The most common one is to assume that the shear stress is proportional to some pressure (the so-called โ$`\alpha ๐ซ`$-formalismโ; e.g. Cannizzo, 1992). It leads to a viscosity law
$$\nu _1=\frac{2\alpha ๐ซ}{3\mathrm{\Omega }\rho }\nu _1(R,z)$$
(12)
where $`๐ซ`$ can be either gas pressure, or total pressure or some combination of the two (Siemiginowska, Czerny & Kostyunin, 1996; Artemova et al., 1996; Hameury et al., 1998; Papaloizou & Terquem, 1999), with different consequences on the disc stability (Camenzind, Demole & Straumann 1986; Clarke, 1988). The other way uses the local version of the $`\alpha `$-prescription (Meyer & Meyer-Hofmeister, 1982)
$$\nu _2=\alpha c_\mathrm{s}\overline{\lambda }_p\nu _2(R,z)$$
(13)
where $`\overline{\lambda }_p`$ is the reduced pressure height scale (see the Appendix, Sect. A). Whatever the expression the authors adopt, the $`\alpha `$-parameter is most often assigned to a fixed value in a disc. However, theoretical arguments, numerical simulations as well as observational constraints indicate that $`\alpha `$ should vary, not only with the radius (Shakura & Sunyaev, 1973; Cannizzo, 1993; Lasota & Hameury, 1998), but also with the altitude via the temperature or density, or other quantities (Brandenburg, 1998). Strictly speaking, the need for strong variations of the $`\alpha `$-parameter means the failure of the $`\alpha `$-viscosity model.
Since $`\nu _1`$ and $`\nu _2`$ are formally different, even with the same $`\alpha `$-parameter, they lead to different vertical structures and consequently to different discs, as we shall see below. Besides, the correspondence between the two laws (a multiplying factor $`\frac{3}{2}\mathrm{\Omega }\lambda _p\mathrm{\Gamma }_1/c_\mathrm{s}`$ from $`\nu _1`$ to $`\nu _2`$, if $`๐ซP`$) is much more subtile than shifting $`\alpha `$: in general, $`\mathrm{\Omega }\lambda _pc_\mathrm{s}`$ at any altitude (this is specially true at the equatorial plane where $`\lambda _\mathrm{p}\mathrm{}`$; see Sect. 4.3). The volumic energy production associated to $`\nu _2`$ is
$$\frac{dF}{dz}=\frac{9}{4}\mathrm{\Omega }^2\lambda _p\alpha \sqrt{\mathrm{\Gamma }_1P\rho },$$
(14)
Note that both Eqs.(12) and (13) satisfy $`\frac{d\nu }{dz}<0`$ above the equatorial plane, meaning that the gas is accreted much faster at the equatorial plane than at the top of the disc. So, these laws implicitly suggest the existence a plan parallel shear which might be able to trigger turbulence and to expand turbulent eddies in the radial direction.
There are still no physical arguments to decide if $`\nu _2`$ is better than $`\nu _1`$. That is why we use both expressions in the following (in particular, for $`\nu _1`$ we take $`๐ซP`$). In fact, a wide class of functions $`\nu (R,z)`$ should be considered and their effects compared. It is possible that turbulent viscosity shows a weaker dependence with $`z`$ than considered so far and even does not depend on the altitude at all. This is precisely the case with the $`\beta `$-viscosity prescription which is suggested by laboratory experiments (Pringle & Rees, 1972; Lynden-Bell & Pringle, 1974; Richard & Zahn, 1999; Duschl, Biermann & Strittmatter, 2000; Hurรฉ, Richard & Zahn, 2000).
### 2.5 Equations for the Eddington atmosphere
In the Eddington approximation, the structure of the atmosphere is governed by three (only two in the absence of self-gravity) first orders ODEs. Two of these (namely Eqs.(5) and (8)) are basically the same as for the disc interior, and the third one controls the optical depth $`\tau `$ in the atmosphere
$$\frac{d\tau }{dz}=\kappa \rho $$
(15)
where $`\kappa `$ is a grey absorption coefficient. In order to prevents any temperature runaway which would lead to the formation of a hot corona (Shaviv & Wehrse, 1991; Rรณลผaลska et al., 1999), we assume that there is no viscous energy generation in this layer and no active turbulence. It means that the atmosphere is not accreted. We are aware that the validity of all these approximations, including the Eddington approximation, is easily open to criticism. The temperature within the atmosphere follows the law (Mihalas, 1978)
$$T^4=\frac{3}{4}T_{\mathrm{eff}}^4\left(\tau +\frac{2}{3}\right)$$
(16)
where $`\tau \frac{2}{3}`$ and $`T_{\mathrm{eff}}`$ is the effective temperature which is fixed by the outgoing flux due to internal and external heating sources, at the photosphereโs base. Actually, the global effect of the disc irradiation can be taken into account at this level by a suitable definition of $`T_{\mathrm{eff}}`$ (Kenyon & Hartmann, 1987; Hubeny, 1990; Robinson, Marsh & Smak, 1993; Dubus et al., 1999). Wee see from Eqs.(5), (15) and (16) that the density gradient is
$$\frac{d\rho }{dz}=\frac{\rho ^2}{\chi _\rho P}\left(\mathrm{\Omega }^2z+4\pi G\mathrm{\Sigma }\right)+\frac{\chi _T\kappa \rho }{4\chi _\rho \left(\tau +\frac{2}{3}\right)}$$
(17)
and may also become positive. As suggested above, we do not take into account turbulent pressure in this layer. This induces a discontinuity in the density (or gas pressure) gradient at the altitude where the disc and its atmosphere join together, but not in the density itself.
### 2.6 Boundary conditions and matching conditions at the disc/atmosphere interface
There are two problems to solve simultaneously: the structure of the disc interior from Eqs.(6)(7)(8) and (11) and the structure of its atmosphere from Eqs.(8)(15) and (17). These differential equations are subject to Dirichlet boundary conditions. At the midplane, the symmetry imposes $`\mathrm{\Sigma }=0`$ and $`F=0`$ like for a star. At the top atmosphere located at $`z=H`$ (to be determined), conditions are $`P=P_{\mathrm{amb}}`$ the ambient pressure (which may include the pressure due to external illumination), $`\tau =0`$ and $`\mathrm{\Sigma }=\frac{1}{2}\mathrm{\Sigma }_\mathrm{t}`$, half the total surface density of the disc and atmosphere. Some authors specify the density (for instance $`\rho 0`$) instead of the pressure at the boundary (e.g. Pojmaลski, 1986; Dรถrrer et al. 1993; El-Khoury & Wickramasinghe, 1999). We choose a value typical of the interstellar medium, namely $`P_{\mathrm{amb}}/k=10^5`$ K.cm<sup>-3</sup> (Duley & Williams, 1988), significantly smaller than in Dโalessio et al. (1998). At the base of the atmosphere located at $`z=h`$ (to be determined too), the matching conditions are $`\tau =\frac{2}{3}`$ and $`F=\sigma T_{\mathrm{eff}}^4`$ (see Sect. 4.1).
It is worthy of note that boundary conditions at the disc surface have a very weak effect on midplane quantities as long as the disc is optically thick and has a large surface density, locally. Such an insensitivity is well known in stellar structure computations (Kippenhahn & Weigert, 1990). Let us quote also that, contrary to a general belief, models that do not consider the atmosphere and use $`P=\frac{2}{3}\frac{g_z}{\kappa }`$ as a boundary condition at $`z=h`$ always violate the Eddington approximation since this relation does not guarantee that the atmosphere has the right optical thickness (i.e. $`\frac{2}{3}`$). Further, the approximation
$$_{P(\mathrm{})}^{P(h)}\frac{\kappa }{g_z}๐P\left(\frac{\kappa }{g_z}P\right)_{z=h}$$
(18)
becomes bad in regions where the surface density of the atmosphere is comparable to that of the disc, when the mean absorption within the atmosphere is low, or when the disc is not optically very thick. For these reasons, we sustain the idea that the use of the Eddington approximation has sense only if the structure of the atmosphere is computed together with that of the disc interior.
## 3 Ingredients and computational method
### 3.1 Equation of state and opacities for a cosmic gas
The physical problem depends on a certain amount of thermodynamical and radiative data. For the present application, the gas has cosmic abundances and is subject to atomic ionizations and molecular dissociations. Local Thermodynamic Equilibrium (LTE) is assumed (Hurรฉ et al., 1994b). The mean mass per particle $`\mu m_\mathrm{H}`$, coefficients $`\chi _T`$ and $`\chi _\rho `$, adiabatic gradient $`_{\mathrm{ad}}`$, heat capacity $`c_p`$ (needed to treat convective transport) and adiabatic exponent $`\mathrm{\Gamma }_1`$ have been computed accurately from chemical abundances at thermal equilibrium (Hurรฉ, 1998). Grids of data with $`\rho `$ and $`T`$ as the inputs have been generated. Interpolations between mesh points are performed with bicubic splines. We have computed a high precision analytical expression fitting $`\mu `$ (and subsequently coefficients $`\chi _T`$ and $`\chi _\rho `$) as functions of the temperature and density with an accuracy less than 4$`\%`$ relative to raw data in the case of a zero metallicity gas (see the Appendix, Sect. B1). This fit can be used for a cosmic gas as well without producing important errors (since metals in a cosmic mix have very low abundances relative to hydrogen, with almost no influence on the EOS and related quantities).
Opacity is probably the most important ingredient of the model and the choice for the grey absorption coefficient $`\kappa `$ is critical, specially near the surface (Mihalas, 1978). Here, we take (Hameury et al., 1998)
$$\kappa =\theta \kappa _\mathrm{R}+(1\theta )\kappa _\mathrm{P}$$
(19)
where $`\kappa _\mathrm{R}`$ and $`\kappa _\mathrm{P}`$ are the Rosseland and Planck means respectively, and $`\theta `$ is a function which varies continuously from 0 to 1 as the optical depth increases from 0 to infinity. A suitable choice is
$$\theta _m(\tau )=\frac{1}{1+\tau ^m}$$
(20)
where the index $`m`$ sets the stiffness of the transition from $`\kappa _\mathrm{P}`$ to $`\kappa _\mathrm{R}`$ ($`m=1`$ is taken in the following applications). The $`\theta `$-function is arbitrary and should strongly influence physical quantities at the disc surface. Also, we are aware that the best grey coefficient is probably neither given by $`\kappa _\mathrm{R}`$, nor by $`\kappa _\mathrm{P}`$, nor by Eq.(19). A flux weighted opacity mean could be better (Mihalas, 1978; Hubeny, 1990).
Tables of Rosseland and Planck means have been taken from various sources (Pollack et al., 1994; Hurรฉ et al., 1994b; Alexander & Fergusson, 1994; Henning & Stognienko, 1996, Seaton et al., 1994). But Planck absorption means published by the Opacity Project (OP) (Seaton et al., 1996) have not be selected due to apparently large errors of unknown origin (Alexander, 1998; Zeippen, 1999). Two grids of data have been generated after passing through filters (mainly running averages) to smooth out discontinuities. Finally, opacity grids cover the temperature range $`1010^8`$ K and extend from extremely high to extremely low densities where the ideal gas and LTE assumptions are expected to fail. An example of Rosseland and Planck means versus the temperature is displayed in Fig.(1). As done for the equation of state and related quantities, opacity values between mesh points are interpolated using bicubic splines. Working with analytical opacities may present some advantages (Burgers & Lamers, 1989; Collin & Dumont, 1990; Bell & Lin, 1994), but they can also produce numerical instabilities in the integration of the disc structure equations since they mostly consist in continuous but non derivable piecewise functions.
### 3.2 Computational method
The altitude of the top atmosphere being not known a priori, many numerical methods for solving Two Boundary Value Problems (TBVPs) are almost unusable in the actual situation. It is however possible to apply variable changes, namely $`_z_\tau `$ in the atmosphere and $`_z_F`$ in the disc interior (Cannizzo & Cameron, 1988) in order to recover a common TBVP with fixed boundaries (e.g. $`\frac{F}{\sigma T_{\mathrm{eff}}^4}[0,1]`$), but this drastically increases vertical gradients near the surface and make the problem more unstable from a numerical point of view. Some groups work with relaxation algorithms (Cannizzo, 1992; Milsom, Chen & Taam, 1994; Dโalessio et al., 1998; Hameury et al., 1998) which require (good) vertical profiles as starting guesses and are known to be rapidly converging methods. However, zones with steep gradients generally need a local mesh refinement. Other authors prefer straight integrators and algorithms based on shooting methods (Lin & Papaloizou, 1980; Meyer & Meyer-Hofmeister, 1982; Mineshige & Osaki, 1983; Smak, 1984; Mineshige, Tuchman & Wheeler, 1990; Rรณลผaลska et al., 1998; Papaloizou & Terquem, 1999) typical of Initial Values Problems (IVPs), with only a few quantities to guess but with possible troubles regarding the precision.
In the present case, we have written a Fortran code named VS<sup>3</sup>KAD which performs the integration of the disc and atmosphere equations using single side shooting and numerical routines specially designed for stiff IVPs. Integration steps are variable, internally to the routines and the accuracy can be as small as the machine precision. In its present version, the code works with the following variables: $`\mathrm{ln}T`$, $`\mathrm{ln}\rho `$, $`\frac{F}{\sigma T_{\mathrm{eff}}^4}`$, $`\frac{2\mathrm{\Sigma }}{\mathrm{\Sigma }_\mathrm{t}}`$ and $`\tau `$, but more clever choices are possible, specially to reduce artificially the stiffness of the equations. We have noticed that the computational time is considerably smaller than with Runge-Kutta algorithms, even with a variable step. Like in many stiff problems, the numerical integration may turn out to be non conservative: integrating from the top atmosphere down to the midplane can give a solution which, in the surface neighborhood, can be very different than that obtained when integrating in the opposite direction (Mineshige, Tuchman & Wheeler, 1990; Press et al., 1992; De Kool & Wickramasinghe, 1999). This has been observed in the present case. It reminds the influence of both boundary conditions and underlying methodology. The best stability and reliability of the results are obtained by starting from the boundary layer where almost all the stiffness of the problem is concentrated.
In practical, given a central mass $`M`$, accretion rate $`\dot{M}`$, $`\alpha `$-parameter, radius $`R`$ and total energy deposition $`\sigma T_{\mathrm{eff}}^4`$ at the bottom atmosphere (see Sect. 4.1), the computation starts at an arbitrary altitude $`z=H^{(0)}`$ (the top atmosphere) where the surface density $`\mathrm{\Sigma }^{(0)}(H^{(0)})`$ is guessed. The integration then proceeds downwards. Once the bottom atmosphere is reached, at an altitude $`z=h^{(0)}`$, the equations for the disc interior are integrated down to the midplane where the net flux and surface density generally differ from zero. By successive iterations on $`H^{(i)}`$ and $`\mathrm{\Sigma }(H^{(i)})`$ โperformed with a Newton-Raphson methodโ, quantities $`F^{(i)}(0)`$ and $`\mathrm{\Sigma }^{(i)}(0)`$ can be driven to very small values. The problem has converged after $`n`$ iterations, when simultaneously $`F^{(n)}(0)0`$ and $`\mathrm{\Sigma }^{(n)}(0)0`$ with the requested precisions
$$ฯต_{}^F\frac{F^{(n)}(0)}{\sigma T_{\mathrm{eff}}^4}ฯต_+^F,$$
(21)
$$ฯต_{}^\mathrm{\Sigma }\frac{2\mathrm{\Sigma }^{(n)}(0)}{\mathrm{\Sigma }_\mathrm{t}}ฯต_+^\mathrm{\Sigma }$$
(22)
and
$$ฯต_{}^H1\frac{H^{(n)}}{H^{(n1)}}ฯต_+^H,$$
(23)
where $`ฯต_+^F`$, $`ฯต_+^F`$, $`ฯต_+^\mathrm{\Sigma }`$, $`ฯต_{}^\mathrm{\Sigma }`$, $`ฯต_{}^H`$ and $`ฯต_+^H`$ are very small, positive values. Since our treatment of self-gravity is very approximate, it is legitimate to allow a much lower (but sufficient) precision on the surface density (this accelerates convergence and lowers the computational time). On average, it takes a few milliseconds CPU on a single user personal computer and $`n25`$ with $`ฯต_0^F=ฯต_0^H=(ฯต_0^\mathrm{\Sigma })^2=10^{10}`$ where $`ฯต_0^F=\mathrm{Sup}(ฯต_{}^\mathrm{F},ฯต_+^\mathrm{F})`$, $`ฯต_0^\mathrm{\Sigma }=\mathrm{Sup}(ฯต_{}^\mathrm{\Sigma },ฯต_+^\mathrm{\Sigma })`$ and $`ฯต_0^H=\mathrm{Sup}(ฯต_{}^\mathrm{H},ฯต_+^\mathrm{H})`$. This precision is sufficient in most cases. When self-gravity is neglected, Eq.(8) becomes obsolete. The iteration is on the flux only and the code then runs much faster (by a factor $`5`$), with $`n6`$.
The code has been tested and compared in many situations (e.g. without self-gravity, without turbulent pressure, with and without external irradiation) with reference models (Rรณลผaลska, 1998; Hameury et al., 1998) and appears to behave very well<sup>1</sup><sup>1</sup>1The author plans to make the executable available to the community. In the meanwhile, vertical structure computations may be performed on special request.. It automatically scales to the central mass, and is able to model any kind of disc: AGN discs, CV discs, YSO discs and sub-nebulae as well. A simulation of D/H enrichments in the PSN with this code is currently under way (Hersant, Gautier, Hurรฉ, 2000).
## 4 A few applications of the model
### 4.1 The hypothesis of a non irradiated disc
Vertically averaged disc models show that the non self-gravitating and gas pressure dominated parts of $`\alpha `$-discs have an aspect ratio $`\frac{h}{R}0.0010.1`$, depending on the central mass and accretion rate mainly. There is a slight flaring (i.e. $`\frac{d\mathrm{ln}h}{d\mathrm{ln}R}1`$) which eventually may vanish in the outermost regions (e.g. Ruden & Pollack, 1991; Bell et al., 1997). The existence of a flaring inner disc is supported by observations: the spectrum of T-Tauri stars and AGN contains a prominent infrared component which is currently interpreted as light re-processing in the superficial layers of the disc (Adams & Shu, 1986; Voit, 1991; Chiang & Goldreich, 1997). It is possible that self-irradiation also plays a role (Fukue, 1992). For circumstellar discs, the flaring angle of bare $`\alpha `$-disc is too low to enhance the infrared spectral component to the observed level (Kenyon & Hartmann, 1987). And the heating by the central star should not produce a significant increase of the flaring (e.g. Dโalessio et al., 1999). But the long wavelength emission can efficiently be boosted if the disc gets into a warped configuration (Terquem & Bertout, 1993; Miyoshi et al., 1995; Bachev, 1999).
Beyond a certain radius, self-gravity becomes important and reduces the disc thickness such that $`\frac{d\mathrm{ln}h}{d\mathrm{ln}R}<0`$ (Sakimoto & Coroniti, 1981; Shore & White, 1982; Cannizzo & Reiff, 1992; Liu, Xie & Ji, 1994; Hurรฉ et al., 1994a; Hurรฉ, 1998). As we shall see below, the same effect is observed from vertical structure computations. It means that outer parts should not receive directly light emitted at the center, contrary to the inner parts. Note that a disc is generally not isolated but embedded into a โwarmโ environment which may in turn heat it up. Discs in YSO are surrounded by an envelope of gas and dust resulting from the cloud core collapse (Mundy, Looney & Welch, 2000; Chick & Cassen, 1997; DโAlessio, Calvet & Hartmann, 1997). In the case of AGN, clouds moving above the disc (BLR clouds) can scatter light back onto it (Shields 1977; Collin-Souffrin, 1987; Osterbrock, 1993). It follows that, even outer regions which are not directly exposed to the luminous central source can be substantially irradiated. The disc response to irradiation is a complex problem to solve self-consistently since it depends on many parameters (size and location of the ionizing source(s), shape of the ionizing spectrum, intensity of irradiation, disc flaring angle, surface albedo, disc optical thickness, gas metalicity, etc.) which are neither well known, nor well constrained by current observations. For discs having a large total surface density, mostly the superficial layers are heated up and ionized, and the midplane is essentially not affected (Sincell & Krolik, 1997; Collin & Hurรฉ, 1999; Nayakshin, Kazanas & Kallman, 1999; Igea & Glassgold, 1999). Deep structural changes (for instance, a temperature inversion or an iso-thermalization) may however occur in some cases (Dโalessio et al. 1998). Although discs should experience some external heating, even in their outermost (self-gravitating) parts, we do not consider irradiation in this paper, for simplicity. So, considering internal viscous heating only, the disc effective temperature, far from the central object, is given by
$$T_{\mathrm{eff}}^4=\frac{3}{8\pi \sigma }\mathrm{\Omega }^2\dot{M}$$
(24)
where $`\dot{M}`$ is the mass accretion rate.
In the following, we apply the code to the computation of the vertical structure of two different systems subject to vertical self-gravity: circumstellar discs and AGN discs. We restrict to central masses of $`1`$ M and $`10^8`$ M respectively, and we show on a few examples how self-gravity and turbulent pressure which are the main specificities of this model affect the disc structure. Computations are performed discarding any kind of instabilities the disc could undergo. These are stopped when the temperature at the top atmosphere (i.e. $`T(H)=\left(\frac{1}{2}\right)^{1/4}T_{\mathrm{eff}}`$) attains 10 K, for both physical and practical reasons.
### 4.2 Effect of turbulent pressure on the vertical temperature and density stratification: an example
To see clearly the influence of turbulent pressure (see Eq.(10)), we have chosen a simulation with $`\alpha =1`$ which corresponds to the very upper limit regarding supersonic turbulence. In addition, convection and self-gravitation are turned off. Under these circumstances, we have computed the vertical structure with viscosity law $`\nu _1`$ for an disc around a massive black hole with $`M=10^8`$ M, at two different radii: $`R=500R_\mathrm{S}`$ ($`R_\mathrm{S}=2GM/c^2`$ is the Schwarzchild radius of the black hole) and $`R=1000R_\mathrm{S}`$. The accretion rate is $`\dot{M}=10^1`$ M/yr. The temperature and density profiles so obtained are plotted in Fig.(2). We see that turbulent pressure makes the disc thicker, as expected intuitively. Actually, according to Eq.(11), the density gradient is lowered with turbulent pressure and the transition from the ambient medium to the disc interior is softer. The flux gradient depending on the density to a positive power โin the absence of convection at leastโ, a larger integration range is needed to obtain a zero net flux at the midplane, hence a thicker disc. The effect is specially visible at $`R=1000R_\mathrm{S}`$, with a $`41\%`$ disc thickening (measured on $`h`$; we have $`33\%`$ on $`H`$). The disc flaring therefore increases by the same relative quantity. Since we have not explored the whole parameter space, it is probably easy to find cases where the effect is more important. The total surface density is almost unchanged (a $`2\%`$ increase only). At $`500R_\mathrm{S}`$, the disc thickening is less important, about $`15\%`$. Note in this example the presence of a density inversion which occurs well below the base of the atmosphere. It tends to be washed out by turbulent pressure, as argued in Sect.2.3.
The way turbulent pressure acts on both the temperature and density at the midplane temperature is however less straightforward and depends intimately on the vertical stratification. For a given effective temperature, a disc with a larger surface density would theoretically have a hotter core. But this happens only in the calculation at $`500R_\mathrm{S}`$, and the effect is very minor.
We have checked that the conclusions derived here-above are qualitatively unchanged using $`\nu _2`$ (see Sect. 4.5 for another effect of turbulent pressure). We will not discuss the case of circumstellar discs because values of the $`\alpha `$-parameter currently adopted for these systems are rather very low (typically $`\alpha 10^310^2`$; DโAlessio, Calvet & Hartmann, 1997) and so, no effect of turbulent pressure is indeed observable. This has been verified.
### 4.3 Effect of the viscosity law: $`\nu _1`$ versus $`\nu _2`$
In this paragraph, we show a few differences between $`\nu _1`$ and $`\nu _2`$. Convection is taken into account and turbulent pressure is left aside. The midplane temperature, total surface density, disc thickness and disc mass are plotted versus radius in Fig.(3) for $`M=1`$ M, $`\dot{M}=10^7`$ M/yr and $`\alpha =10^3`$. These parameter values are typical of discs around T-Tauri stars (DโAlessio, Calvet & Hartmann, 1997). The disc mass $`M^{\mathrm{disc}}`$ is
$$M^{\mathrm{disc}}(R)=2\pi _{R_{\mathrm{in}}}^RR^{}\mathrm{\Sigma }_\mathrm{t}(R^{})๐R^{}\pi \mathrm{\Sigma }_\mathrm{t}R^2$$
(25)
where $`R_{\mathrm{in}}`$ is the inner edge of the disc (unimportant as long as $`RR_{\mathrm{in}}`$). Results computed without self-gravity are also shown in comparison. We see that, in the non self-gravitating region, differences on the geometrical thickness and midplane temperature are small, the disc being slightly thicker and hotter with $`\nu _1`$ than with $`\nu _2`$. The most important effect is on the surface density and consequently on the disc mass: the $`\nu _1`$-prescription leads to slightly more massive disc than the $`\nu _2`$-prescription (by a factor $`2`$ here).
We see that the disc is definitely affected by self-gravity from approximately $`7`$ AU whatever the viscosity law (in fact, a little bit less with $`\nu _1`$ which agrees with the fact that that prescription leads to more massive discs, as long as $`M^{\mathrm{disc}}M`$). Let us remind that the importance of self-gravity can be measured by the quantity (see Eq.(5))
$$\zeta (R,z)=\frac{4\pi G\mathrm{\Sigma }(R,z)}{\mathrm{\Omega }^2z}$$
(26)
which remains finite at $`z=0`$ as a Taylor expansion shows. It follows that a โnaturalโ definition for the limit between the standard disc and the self-gravitating disc can be $`\zeta =1`$ at the midplane (if $`\zeta >1`$ then the disc contribution to gravity exceeds that due to the central object, vertically).
On the example, the disc thickness goes through a maximum near 10 AU for both viscosity laws. Beyond this radius, the disc becomes thinner and thinner. Regions located farther away can therefore not intercept photons emitted at the center. Interestingly, the temperature and the surface density show very different behaviors. With law $`\nu _1`$, the temperature decreases monotonically as $`R`$ increases, almost as in the absence of self-gravity. So does the total surface density. The midplane density, not shown on graphs, gently increases with the radius. Note that $`\mathrm{\Sigma }_\mathrm{t}`$ (and $`M^{\mathrm{disc}}`$) are surprisingly not affected by self-gravity. The origin of this insensitivity has not been identified at the time being. Conversely, with law $`\nu _2`$, the surface density and the density violently increase with the radius. The disc thickness falls in also very rapidly. The temperature reaches a plateau. In fact, this somewhat โsingularโ behavior is predicted by the vertically averaged model (Shlosman & Belgelman, 1987; Hurรฉ et al. 1994a; Hurรฉ, 1998; see also Duschl, Biermann & Strittmatter, 2000). In particular, the temperature $`T_{\mathrm{sg}}`$ at the plateau is mainly fixed by the ratio $`\dot{M}/\alpha `$, namely
$`T_{\mathrm{sg}}`$ $`=`$ $`{\displaystyle \frac{\mu m_\mathrm{H}}{k}}\left({\displaystyle \frac{4G^2\dot{M}^2}{9\alpha ^2}}\right)^{1/3}`$ (27)
$``$ $`24100\mu \left({\displaystyle \frac{\dot{M}}{1\mathrm{M}_{}/\mathrm{yr}}}\right)^{2/3}\alpha ^{2/3}\mathrm{K}`$
For $`\dot{M}/\alpha =10^4`$ M/yr as in Fig.(3), Eq.(27) yields $`T_{\mathrm{sg}}120`$ K (assuming $`\mu =2.33`$). Although this value is derived from the one zone model, it is in good agreement with the vertical structure computation which gives a plateau at about $`70`$ K. Note that the drastic increase of the surface density produces an important increase of the disc mass which even exceeds the value obtained with law $`\nu _1`$. When $`M^{\mathrm{disc}}M`$ (this occurs at $`2328`$ AU on the example), the disc becomes very self-gravitating, probably gravitationally unstable (Goldreich & Lynden-Bell, 1965; Shu et al., 1990) and should therefore not be well described with current viscosities (e.g. Lin & Pringle, 1987) nor by steady state solutions (Goldreich & Lynden-Bell, 1965). One reaches here the limit of the model. It is however interesting to note that, asymptotically, there is either a hot very dense solution or a cold more diffuse solution (see Paczyลski, 1978).
The same quantities obtained under the same conditions but for an AGN disc with $`M=10^8`$ M, $`\dot{M}=10^2`$ M/yr and $`\alpha =0.1`$ are displayed in Fig.(4). Globally, we find similar trends. The singular behavior observed in the previous example with $`\nu _2`$ inside the self-gravitating regime seems even stronger. Self-gravitation becomes important from about $`500R_\mathrm{S}`$. The midplane temperature stabilizes radially at $`3900`$ K. This value again compares quite well with Eq.(27) which gives $`6600`$ K (assuming $`\mu =1.27`$). The presence of extremely steep surface density and density radial gradients probably means that the assumption made on the disc potential (see Sect. 2.1) is no longer valid. This would be worthwhile to check. The physical picture there is then that of a disc surrounded by a very dense ring ($`\mathrm{\Sigma }`$ increases by more than one order of magnitude over a relative length $`\frac{\mathrm{\Delta }R}{R}15\%`$) which contains almost all the disc mass and angular momentum. When $`M^{\mathrm{disc}}M`$ (beyond a few $`10^3R_\mathrm{S}`$ for law $`\nu _2`$), effects of self-gravity are expected to be global, with a change in the rotation law.
All these results require a few comments. First, it is important to note that the effect of self-gravity is slightly underestimated when it is neglected in the computations. This is specially true with law $`\nu _2`$. Besides, self-gravity becomes important before $`\zeta `$ reaches unity, like in the one zone model (Hurรฉ, 1998). At the inner edge of the self-gravitating disc (that is at $`\zeta =1`$ following our definition), the disc mass is much less than the central mass: we find $`\frac{M^{\mathrm{disc}}}{M}10^210^1`$ (depending on the viscosity law) in the YSO case, and $`\frac{M^{\mathrm{disc}}}{M}10^3`$ in the AGN case. This suggests that low mass discs must not be automatically classified as non self-gravitating discs as often asserted. In particular, it is not excluded that T-Tauri discs we observe be subject to vertical self-gravity, despite their relatively low mass (Beckwith et al., 1990).
It has been stated in Sect. 2.4 that the relation between $`\nu _1`$ and $`\nu _2`$ is not trivial. In the non self-gravitating parts however, midplane temperatures $`T(0)`$ show a striking quasi-parallel variation with the radius in logarithmic scales (this is also true for quantities $`\mathrm{\Sigma }_\mathrm{t}`$ and $`h`$), which could be attributed to the fact that $`\nu _1`$ is indeed proportional to $`\nu _2`$ (or one goes from one solution to the other by a change of $`\alpha `$; see power law solutions for $`\alpha `$-discs). To demonstrate that these two viscosity prescriptions are intrinsically different (as suggested by what happens in the self-gravitating regions), we have computed the ratio $`\frac{\nu _1}{\nu _2}=\frac{3}{2}\mathrm{\Omega }\lambda _p\mathrm{\Gamma }_1/c_\mathrm{s}`$ with $`\nu _1`$, for the examples discussed above. We have chosen two radii, the one lies in the classical disc part, the other is in the self-gravitating part. The results are shown in Fig.(5).
### 4.4 Example of internal structure: 2D-density maps
We show in Fig.(6) the density field within a circumstellar disc for the same parameter values as in Fig.(3). This simulation has been carried out using $`\nu _1`$ and $`\nu _2`$ as well, including convection, self-gravity and turbulent pressure. We also indicate iso-contours of the temperature, the limit between gas pressure dominated zones and radiative pressure dominated ones, and lines where $`\zeta (R,z)=\frac{1}{10},1`$ and $`10`$. We notice that the density and the temperature show strong variations in between the midplane to the photosphereโs base and we are far from a vertically isothermal, homogenous disc as often considered. Note, specially for viscosity law $`\nu _1`$, the re-increase of the midplane density in the radial direction (another density inversion) as self-gravity gains in importance. It is remarkable that self-gravity does not appear suddenly (i.e. on a small radial range) but installs gently and steadily over a very extended domain. For instance, with law $`\nu _2`$, it contibutes by $`10\%`$ in the hydrostatic equilibrium near $`24`$ AU and by $`50\%`$ at about 7 AU, and all disc quantities are significantly modified over this domain. We see clearly that the disc density is slightly smaller with $`\nu _2`$ than with $`\nu _1`$ (same color code for both maps).
We notice that the whole disc is optically very thick, due to its large surface density. It is then likely that a mean external heating should not change noticeably the temperature and density at the midplane. From this point of view, the neglect of irradiation is entirely justified. It is interesting to see that, as long as an $`\alpha `$-model can be applied to the Solar Nebula (e.g. Ruden & Pollack, 1991; Papaloizou & Terquem, 1999; Drouart et al, 1999), our giant planets (displayed on graphs) lie inside the self-gravitating part of the disc (this is true whatever the accretion rate, but depends on the $`\alpha `$-parameter; see the next Section). As already noticed (Ruden & Pollack, 1991), this coincidence is somewhat striking and we do not know to what extent self-gravity could have played a role in the process of the formation of outer planet and other objects (Lissauer, 1993).
Similar 2D-maps are displayed in Fig.(7) for an AGN disc with $`M=10^8`$ M, $`\dot{M}=10^2`$ M/yr and $`\alpha =10^1`$, using the same color table. Note the existence of an inner radiative pressure dominated region at $`R100R_\mathrm{S}`$, common for AGN discs as well as the great steepness of density gradients near the surface which may cause numerical difficulties. Also visible is the density inversion at $`R160200R_\mathrm{S}`$ which has less amplitude and is less extended radially in the presence of convection and turbulent pressure.
### 4.5 A criterion to check the importance of vertical self-gravity
We see on both Figs.(6) and (7) that $`\zeta (R,\mathrm{})<\zeta (R,0)`$ and iso-values of $`\zeta (R,z)`$ form very curved lines (โdiaboloโ-shape surfaces in 3D). This is important if one wishes to estimate correctly the position of the self-gravitating regime. For instance, in the case of the law $`\nu _1`$ discussed in Fig.(6), $`\zeta =1`$ in the midplane at 7 AU (we have $`450R_\mathrm{S}`$ in the AGN case depicted in Fig.(7)) whereas this occurs in the disc atmosphere at about twice the distance.
For many purposes, it is interesting to know the position of the (vertically) self-gravitating disc. That is why we have performed a systematic computation to find the radius $`R_{\mathrm{sg}}`$ satisfying $`\zeta (R_{\mathrm{sg}},0)=1`$ as a function of the accretion rate, for 4 values of the $`\alpha `$-parameter in the range $`10^31`$, and for $`M=1`$ M and $`M=10^8`$ M. These parameter values should cover a wide variety of circumstellar discs and AGN discs. The calculations have been performed with $`\nu _2`$ including convection, self-gravitation, with and without turbulent pressure (the use of $`\nu _1`$ would systematically give a slightly lower value of $`R_{\mathrm{sg}}`$). The results are plotted in Fig.(8) for parameter values typical of circumstellar discs. We see that the location of self-gravitating regime is very sensitive to the viscosity parameter. The larger the value of $`\alpha `$, the further away the self-gravitating region. Regarding the sensitivity to the accretion rate, there are three different trends. For $`\alpha 10^2`$, the lower $`\dot{M}`$, the larger $`R_{\mathrm{sg}}`$. For $`\alpha 0.1`$, the dependence is rather weak: $`R_{\mathrm{sg}}27\pm 3`$ AU. For higher values of the viscosity parameter (generally not appropriate to fit properties of observed discs), $`R_{\mathrm{sg}}`$ is an increasing function of $`\dot{M}`$. Note also that turbulent pressure, when important, slightly increases $`R_{\mathrm{sg}}`$ ($`12\%`$ in the example for $`\alpha =1`$ and the highest accretion rate).
Results for AGN discs are shown in Fig.(9). As above, large viscosity parameters imply wider non self-gravitating inner disc (because discs are less dense with large values of the $`\alpha `$-parameter), and turbulent pressure pushes further away the self-gravitating regime. For $`\alpha 10^2`$, we have $`R_{\mathrm{sg}}520\times 10^{15}`$ cm. The sensitivity to the accretion rate remains weak. For higher values of the viscosity parameter and for very low to moderate accretion rates ($`\dot{M}10^210^1`$ M/yr, depending on $`\alpha `$), $`R_{\mathrm{sg}}`$ globally decreases as $`\dot{M}`$ increases. For higher accretion rates, $`R_{\mathrm{sg}}`$ increases as $`\dot{M}`$ increases and self-gravity is important within the radiative pressure dominated region.
It is important to notify that the โToomreโs criterionโ ($`Q_\mathrm{T}\frac{c_\mathrm{s}(0)\mathrm{\Omega }}{\pi G\mathrm{\Sigma }_\mathrm{t}}1`$) is an helpful tool to trace regions where some gravitational instabilities occur in a stellar disc (Toomre, 1965), but it differs significantly from the criterion derived by Goldreich & Lynden-Bell (1965) for gaseous discs ($`Q_{\mathrm{GLB}}\frac{4\mathrm{\Omega }^2}{\pi G\rho }1`$). Besides, $`Q_\mathrm{T}`$ is commonly written in different forms which are not strictly equivalent (e.g., Ruden & Pollack, 1991; Sincell & Krolik, 1997; Dโalessio, Calvet & Hartmann, 1997, Papaloizou & Terquem, 1999), specially in a 2D-model. Althought these $`Q`$-parameters are undisputably related to the $`\zeta `$-parameter in some ways (for instance, $`Q_\mathrm{T}\times \zeta (R,\mathrm{})=2\frac{c_\mathrm{s}(0)}{\mathrm{\Omega }H}>1`$ and $`Q_{\mathrm{GLB}}\times \zeta (R,0)16`$), they must be used with care if one wishes to check the importance of self-gravity. They may lead to noticeable over-estimates of $`R_{\mathrm{sg}}`$, specially if, as it is always the case, the $`Q_\mathrm{T}`$-parameter is computed from models that discards self-gravity. This is illustrated in Fig.(10) where we have plotted values of $`\zeta `$ at three key altitudes and $`Q_\mathrm{T}`$ versus the radius. It seems therefore preferable to make the check on $`\zeta (R,0)`$ rather than on any other quantities, otherwise the importance of self-gravity is under-estimated.
## 5 Conclusion
In this article, we have presented the equations accounting for the vertical structure of steady state keplerian accretion discs, including simultaneously turbulent pressure, convective transport in the framework of the Mixing Length Theory, and the disc self-gravity within the infinite slab approximation. The emphasis has been placed on outer regions where vertical self-gravity becomes increasingly influencal. We have shown that turbulent pressure makes the disc thicker, but it can be neglected if $`\alpha 0.1`$. Also, self-gravity may be important even if the disc mass is very low compared to the central object, contrary to usual assertions. Further, the transition from the classical disc to the self-gravitating disc takes place gradually and so, concerns a large radial domain.
Another important conclusion of this work is that the model of mechanical energy deposition towards the vertical direction is crucial. This is not surprising since viscosity is the source of heating and accretion. Two different behaviors are predicted inside the self-gravitating (and gravitationally unstable) region. On the one hand, the $`\alpha ๐ซ`$-formalism yields a solution which is asymptotically cold and diffuse, whereas the local version of the standard prescription gives a more โsingularโ solution, asymptotically hot and dense, which is also predicted by the vertically averaged model. We can imagine that a suitable choice for the function $`\nu (z)`$ should provide any intermediate solution, meaning that this model for outer regions has almost no predictive power. Although the present stationary keplerian disc model is probably irrelevant to describe the disc structure in its strongly self-gravitating parts, the existence of two asymptotically distinct solutions might indicate that the outer disc can get into, either a flat and cold configuration where fragmentation and formation of indiviual clouds (Kumar, 1999) and compact objects like planets and stars (Collin & Zahn, 1999) could occur, or into a thick diffuse configuration (a torus) (Paczynski, 1978). Time dependent simulations should shed light on this question.
This study confirms that discs around super-massive black holes are self-gravitating close to the center, beyond a few hundreds Schwarzchild radii, depending on the central mass accretion rate and $`\alpha `$-parameter. We have found that the disc surface density remains high when self-gravity dominates and that the disc mass definitely rises outwards. Let us remind that a current problem with the fueling of active nuclei is that the accretion time scale $`M^{\mathrm{disc}}/\dot{M}`$ is usually much too long at large radii, that is why other more efficient mechanisms are invoked (Frank, 1990). However, given the sensitivity of the disc structure to viscosity, it is not excluded that depth dependent viscosity laws other than considered here would lead to smaller surface density distributions, and consequently to shorter accretion time scales. It is therefore important to test other models for the energy deposition along the $`z`$-axis as well as other kinds of viscosity prescriptions (Lin & Pringle, 1987; Cannizzo & Cameron, 1988; Richard & Zahn, 1999; see Hurรฉ & Richard, 2000).
With accretion rates and values of the viscosity parameter usually considered to model discs in T-Tauri stars, we conclude that our giant planets (if not Jupiter, the other ones) were probably formed within the self-gravitating part of the Solar Nebula (Drouart et al. 1999). This is a fortiori true if these objects experienced any inward migration. More generally, as long as the $`\alpha `$-theory may be applied to describe the inner parts of circumstellar discs around forming stars, self-gravity is expected to play a role at about ten AU from the center. This is not uncompatible with the discs we observed (Beckwith et al., 1990).
###### Acknowledgements.
I specially thank A. Rรณลผaลska for many helpful discussions and for providing some reference vertical profiles during the writing and test of the computational code, S. Collin for reading the manuscript. I thank also D. Gautier for highlights on circumstellar discs and the Primitive Solar Nebular, J.-P. Zahn for pointing out the role of turbulent pressure. I am personally grateful to J.M. Hameury for his initiation, some years ago, to the physics of stellar interiors and vertical structure computations. I also thank D. Richard for interesting comments on numerical aspects as well as F. Hersant for widely testing the code. Also, I thank for their hospitality, people at the ITA-Heidelberg where this work was completed, and specially W.J. Duschl.
## Appendix A Treatment for convection in the presence of self-gravitation
Convection is treated with in the framework of the Mixing Length Theory (MLT) (Cox & Giuili, 1968). It is modified in order to include turbulent pressure in the determination of the pressure scale height, and self-gravititation. Following the original version of the MLT, we neglect the variation of gravity with the altitude $`z`$. It is likely that turbulent pressure modify significantly the velocity of rising elements within convectively unstable zones, specially if the $`\alpha `$-parameter is close to unity, but this effect is not considered here. The total pressure height scale writes
$$\lambda _\mathrm{p}\frac{dz}{d\mathrm{ln}\left(P+p_\mathrm{t}\right)}>0$$
(28)
where $`P`$ is the total (gas plus radiation) pressure and $`p_\mathrm{t}`$ is the turbulent pressure. This expression being singular at the midplane ($`\lambda _p\mathrm{}`$) (see Eq.(5)), the pressure scale height is usually limited to the disc thickness
$$\overline{\lambda }_\mathrm{p}=\mathrm{Inf}[h,\frac{P+p_\mathrm{t}}{\rho \left(\mathrm{\Omega }^2z+4\pi G\mathrm{\Sigma }\right)}]$$
(29)
where $`\rho `$ is the gas density, $`\mathrm{\Omega }`$ is the rotation frequency, $`\mathrm{\Sigma }`$ is the surface density and $`h`$ is the altitude of the bottom atmosphere.
According to Eq.(7), the temperature gradient is
$$\frac{d\mathrm{ln}T}{dz}=\frac{}{\overline{\lambda }_p}$$
(30)
where the actual gradient $`\frac{d\mathrm{ln}T}{d\mathrm{ln}P}`$ depends on the adiabatic gradient $`_{\mathrm{ad}}=\left(\frac{d\mathrm{ln}T}{d\mathrm{ln}P}\right)_{\mathrm{ad}}`$ (computed from the EOS; see the next Section) with respect to the fictitious radiative gradient $`_\mathrm{r}`$ defined as
$$_\mathrm{r}=\frac{3\rho \kappa F\overline{\lambda }_p}{16\sigma T^4}$$
(31)
where $`F`$ is the flux to be transported upwards and $`\kappa `$ is a grey absorption coefficient which must approach the Rosseland mean $`\kappa _\mathrm{R}`$ at great optical depth (see Sect. 3.1).
When $`_\mathrm{r}_{\mathrm{ad}}`$ (i.e. in convectively stable zones), heat is transported through radiation only and we have $`=_\mathrm{r}`$. Conversely, when $`_\mathrm{r}>_{\mathrm{ad}}`$, $``$ is computed from
$$=(1x^3)_\mathrm{r}+x^3_{\mathrm{ad}}$$
(32)
where $`x`$ is the real root of the third degree equation
$$\frac{9}{4}^2x^3+x^2+x\frac{9}{4}^2=0$$
(33)
with
$$=\left[\frac{4}{9}๐^2\left(_\mathrm{r}_{\mathrm{ad}}\right)\right]^{1/3}$$
(34)
and
$$๐=\frac{1}{48\sqrt{2}\sigma }\frac{c_p\kappa _\mathrm{R}\alpha _{\mathrm{MLT}}^2\lambda _p^2}{T^3}\sqrt{\frac{\rho ^5}{P}}\left(\mathrm{\Omega }^2z+4\pi G\mathrm{\Sigma }\right)$$
(35)
where $`\alpha _{\mathrm{MLT}}`$ is the mixing length parameter ($`\alpha _{\mathrm{MLT}}=1.5`$ here) and $`c_p`$ is the constant pressure specific heat capacity.
## Appendix B Notes on the EOS and related quantities
### B.1 High precision fitting formula for $`\mu `$ as functions of the temperature and density โ $`\chi `$-coefficients
Within the assumption of LTE, the pressure of a mixture of radiation and ideal gas undergoing atomic ionization molecular dissociation is given by Eq.(5). In particular, the mean mass per particle $`\mu `$, in units of the proton mass, is defined as
$$\mu =\frac{_im_ip_i}{m_\mathrm{H}_ip_i}$$
(36)
where $`m_i`$ and $`p_i`$ are respectively individual mass and partial pressure of the chemical compounds, and the sum extends over the total number of compounds, including electrons. We have computed equilibrium abundances for a mixture of hydrogen and helium (H:He=1:$`\frac{1}{10}`$) only as function of $`\rho `$ and $`T`$ (see Hurรฉ, 1998). Data relative to the EOS so obtained are plotted versus the temperature in Fig.(11) for three values of the gas density. In particular, the resulting function $`\mu (\rho ,T)`$ has then been fitted with great accuracy by the following expression
$$\mu ^{\mathrm{fit}}=\delta _0+\underset{i=1,4}{}\delta _i\mathrm{tanh}\mathrm{\Phi }_i,$$
(37)
with
$$\mathrm{\Phi }_i=\frac{\mathrm{log}(T)\theta _i}{\mathrm{\Delta }_i}$$
(38)
where coefficients $`\delta _i`$ and $`\tau _i`$ and $`\mathrm{\Delta }_i`$ are listed in Tab.(1). Hyperbolic functions account successively for the transitions H<sub>2</sub>/HI, HI/HII, He/HeI and HeI/HeII. Relative errors on raw data $`\frac{\mathrm{\Delta }\mu }{\mu }=\frac{\mu ^{\mathrm{fit}}\mu }{\mu }`$ are displayed in Fig.(12) and never exceed 4 % over the whole domain of temperature and density. Besides, the standard deviations with respect to raw data are less than 1.3 %. We have tried to fit the residuals with a few Gaussians, but no major improvement has been obtained. Note that, by construction, the fitting formula shows no singular behavior at high of low densities and temperature. The fit can be used for a cosmic gas without producing large errors (the effect of heavy elements on the EOS is very weak as long as $`ZX`$).
The temperature and density exponents of the total pressure are repectively given by
$$\chi _T=\left(\frac{\mathrm{ln}P}{\mathrm{ln}T}\right)_\rho =4(1\beta )+\beta \left[1\left(\frac{\mathrm{ln}\mu }{\mathrm{ln}T}\right)_\rho \right]$$
(39)
and
$$\chi _\rho =\left(\frac{\mathrm{ln}P}{\mathrm{ln}\rho }\right)_T=\beta \left[1\left(\frac{\mathrm{ln}\mu }{\mathrm{ln}\rho }\right)_T\right]$$
(40)
where $`\beta `$ is the ratio of gas pressure to total (gas plus radiation) pressure. It follows from Eq.(37)
$`\chi _T`$ $`=`$ $`{\displaystyle \frac{\beta }{\mu \mathrm{log}10}}{\displaystyle \underset{i=1,4}{}}{\displaystyle \frac{\delta _i}{\mathrm{\Delta }_i}}\left(1\mathrm{tanh}^2\mathrm{\Phi }_i\right)`$ (41)
$`+`$ $`\beta \left(1\delta _0\right)+4(1\beta )`$
and
$`\chi _\rho `$ $`=`$ $`{\displaystyle \frac{\beta }{\mu \mathrm{log}10}}{\displaystyle \underset{i=1,4}{}}{\displaystyle \frac{1\mathrm{tanh}^2\mathrm{\Phi }_i}{\delta _i}}\left({\displaystyle \frac{d\theta _i}{d\mathrm{log}\rho }}\right)_T`$ (42)
$`+`$ $`\beta \left(1\delta _0\right)`$
and the precision with respect to raw data is also of the order of a few percents.
### B.2 Low precision expressions for coefiicients $`\chi _T`$, $`\chi _\rho `$, $`_{\mathrm{ad}}`$, $`\mathrm{\Gamma }_1`$ and $`c_p`$
The adiabatic gradient $`_{\mathrm{ad}}`$, heat capacity at constant pressure $`c_p`$ and adiabatic exponent $`\mathrm{\Gamma }_1`$ are respectively defined by
$$_{\mathrm{ad}}=\left[\chi _T+\frac{\chi _\rho }{\chi _T}\left(\frac{U}{\mathrm{ln}\rho }\right)_T\frac{\rho }{P}\right]^1$$
(43)
where $`U`$ is the specific internal energy of gas and radiation,
$$c_p=\frac{\chi _TP}{\chi _\rho \rho T_{\mathrm{ad}}}$$
(44)
and
$$\mathrm{\Gamma }_1=\frac{\chi _\rho }{1\chi _T_{\mathrm{ad}}}$$
(45)
Since disc models are still very uncertrain, one may use, as a first approximation, low precision expressions for all coefficients related to the EOS discarding changes in the internal energy due to ionizations and dissociations, for instance
$$\chi _\rho \beta $$
(46)
$$\chi _T43\beta $$
(47)
and then
$$_{\mathrm{ad}}\frac{1}{43\beta +\frac{12\beta }{43\beta }\left(1\frac{21}{24}\beta \right)},$$
(48)
$$c_p\left(\frac{43\beta }{\beta ^2}\right)\frac{1}{\mu _{\mathrm{ad}}}\frac{k}{m_\mathrm{H}}$$
(49)
and
$$\mathrm{\Gamma }_1\frac{\beta }{1(43\beta )_{\mathrm{ad}}}$$
(50)
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# Isovector and Isoscalar superfluid phases in rotating nuclei
## Abstract
The subtle interplay between the two nuclear superfluids, isovector T=1 and isoscalar T=0 phases, are investigated in an exactly soluble model. It is shown that T=1 and T=0 pair-modes decouple in the exact calculations with the T=1 pair-energy being independent of the T=0 pair-strength and vice-versa. In the rotating-field, the isoscalar correlations remain constant in contrast to the well known quenching of isovector pairing. An increase of the isoscalar (J=1, T=0) pair-field results in a delay of the bandcrossing frequency. This behaviour is shown to be present only near the N=Z line and its experimental confirmation would imply a strong signature for isoscalar pairing collectivity. The solutions of the exact model are also discussed in the Hartree-Fock-Bogoliubov approximation.
PACS numbers : 21.60.Cs, 21.10.Hw, 21.10.Ky, 27.50.+e
There is overwhelming evidence that the isovector, T=1 pairing field among identical nucleons is an essential component of the nuclear mean-field potential. The bulk of nuclear ground-state properties, like the odd-even mass differences and the moments of inertia of deformed nuclei can be accounted for by considering nucleons to be in a superfluid (T=1, J=0) paired-phase. These effects have been studied mostly in heavier nuclei with N$`>`$Z, where the Fermi surfaces of protons and neutrons lie in different major shells.
In recent years, however, due to a substantial progress achieved in the sensitivity of the detecting systems it has been possible to study nuclei near the N$`=`$Z line in the mass A$`=`$70 and 80 regions. Furthermore, with the availability of radioactive beams these studies are expected to reach even heavier N$`=`$Z nuclei. For these nuclei, one expects the pairing between protons and neutrons to become important, since the Fermi surfaces of both protons and neutrons lie in the same major shell.
The role of the isovector T=1 pairing between protons and neutrons in the low-spin regime has been discussed in recent studies by . The importance of the isoscalar T=0 pairing can be inferred from masses and studies of high-spin states. However, most of these studies are based on the mean-field approximation which predict a transitional behaviour for rotating nuclei for the T=1 and T=0 pair-fields as a function of the rotational frequency and the strength of the T=0 interaction.
The purpose of the present study is to examine properties of the isoscalar and isovector correlations within an exactly soluble model of a deformed sinlge-j shell and to compare to the predictions of the mean-field HFB approximation. The observable consequences of the T=0 pair-field which have remained illusive are also discussed in the present study.
The model hamiltonian consists of a cranked deformed one-body term and a scalar two-body interaction
$$H^{}=h^{}+V_2,$$
(1)
where,
$$h^{}=h_{def}\omega J_x,$$
(2)
with
$$h_{def}=4\kappa \sqrt{\frac{4\pi }{5}}\underset{ij}{}<j|Y_{20}|i>\delta _{\tau _i\tau _j}\delta _{m_im_j}c_j^{}c_i.$$
(3)
The labels $`i,j,\mathrm{}`$ denote the magnetic quantum-number ($`m`$) of the j- shell and the isospin projection quantum-number $`\tau `$ \[$`\tau `$=1/2 (neutron) and -1/2(proton)\]. The deformation energy $`\kappa `$ is equal to the usual deformation parameter $`\beta `$ by $`\kappa =0.16\mathrm{}\omega (N+3/2)\beta `$ in units of G (ref. ). The two-body interaction in Eq. 1 is given by
$$V_2=\frac{1}{2}\underset{JMTT_z}{}E_{JT}A_{JM;TT_z}^{}A_{JM;TT_z},$$
(4)
with $`A_{JM;TT_z}^{}=(c_{j\frac{1}{2}}^{}c_{j\frac{1}{2}}^{})_{JM;TT_z}`$ and $`A_{JM;TT_z}=(A_{JM;TT_z}^{})^{}`$. For the antisymmetric-normalized two-body matrix-element ($`E_{JT}`$), we use the delta-interaction which for a single j-shell is given by
$$E_{JT}=G\frac{(2j+1)^2}{2(2L+1)}\left\{\left[\begin{array}{ccc}j& j& J\\ \frac{1}{2}& \frac{1}{2}& 0\end{array}\right]^2+\frac{1}{2}\{1+(1)^T\}\left[\begin{array}{ccc}j& j& J\\ \frac{1}{2}& \frac{1}{2}& 1\end{array}\right]^2\right\}$$
(5)
where the bracket $`[]`$ denotes the Clebsch-Gordon coefficient.
As mentioned in the introduction, one of the objectives of the present work is to investigate the HFB approximation. In the following, we present some basic HFB formulae, for details see for instance ref.. The HFB equations are given by
$$^{}\left(\begin{array}{c}U\\ V\end{array}\right)=E_i^{}\left(\begin{array}{c}U\\ V\end{array}\right),$$
(6)
where
$`^{}=`$ $`\left(\begin{array}{cc}h^{}& \mathrm{\Delta }\\ \mathrm{\Delta }^{}& (h^{})^{}\end{array}\right).`$ (9)
with
$`h_{ij}^{}`$ $`=`$ $`ฯต_{ij}^{}+\mathrm{\Gamma }_{ij},`$ (10)
$`ฯต_{ij}^{}`$ $`=`$ $`<i|h_{def}|j>(\lambda _pZ+\lambda _nN+\omega m_i)\delta _{ij},`$ (11)
$`\mathrm{\Gamma }_{ij}`$ $`=`$ $`{\displaystyle \underset{kl}{}}<ik|v_a|jl>\rho _{lk},`$ (12)
$`\mathrm{\Delta }_{ij}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{kl}{}}<ij|v_a|kl>\kappa _{kl}.`$ (13)
$$\rho =V^{}V^T,\kappa =V^{}U^T=UV^{}.$$
(14)
In order to evaluate the angular-momentum dependence of the pair-energy, we define the coupled pair-field through
$`\mathrm{\Delta }_{ij}={\displaystyle \underset{JMTT_z}{}}\left[\begin{array}{ccc}j& j& J\\ m_i& m_j& M\end{array}\right]\left[\begin{array}{ccc}\frac{1}{2}& \frac{1}{2}& T\\ \tau _i& \tau _j& T_z\end{array}\right]\mathrm{\Delta }_{JT},`$ (19)
with
$`\mathrm{\Delta }_{JT}=E_{JT}{\displaystyle \underset{ij}{}}`$ (20)
$`\left[\begin{array}{ccc}j& j& J\\ m_i& m_j& m_i+m_j\end{array}\right]\left[\begin{array}{ccc}\frac{1}{2}& \frac{1}{2}& T\\ \tau _i& \tau _j& \tau _i+\tau _j\end{array}\right]\kappa _{ij}.`$ (25)
The pair-energy can now be expressed in terms of the coupled pair-fields as
$$E_{pair}=\frac{1}{2}\underset{JT}{}\frac{\mathrm{\Delta }_{JT}\mathrm{\Delta }_{JT}^{}}{E_{JT}}.$$
(26)
The above expression is quite useful since in the exact calculations there is no gap parameter $`\mathrm{\Delta }`$, but one may associate โ$`E_{pair}`$โ with the expectation value of the two-body residual interaction, $`V_2`$, Eq.4 To obtain the $`\mathrm{\Delta }`$-value from the exact analysis, Eq.26 is then simply inverted.
The HFB solutions have been obtained by solving the Eqs. (6-14) self-consistently. In order to treat both the T=0 and the T=1 pair-fields simultaneously, it is necessary to define complex HFB potentials, since the symmetries of the T=1 and T=0 n-p fields are different. The initial complex HFB wavefunctions have been constructed by using the expressions for real and imaginary $`V`$โs and $`U`$โs of the HFB transformation in terms of the pair-gaps. We would like to mention that no symmetry restrictions have been imposed on the HFB wavefunction since it is known that symmetries lead to exclusion of particular correlations. For more details concerning the HFB-transformation in the presence of both T=1 and T=0 pairing, we refere the reader to refs. .
Several mean-field studies show that the T=0 and T=1 pairing-modes are exclusive in the BCS-approximation. The system is always choosing the mode that generates the lowest energy, which in the case of equal weight for each pair results in either T=0 and T=1 pairing. In the presence of approximate particle-number projection, the two modes coexist, but only above a critical strength. Using a more complex model space also results in the possibility of mixed solutions. The question, therefore, arises whether the exclusiveness is persistent in an exact model. Fig. 1 shows the size of the T=1 correlations as a function of increasing T=0 strength in the exact analysis. The figure clearly shows that the two modes are essentially independent. There is no critical strength for either pairing mode and therefore one expects to have both modes present in nuclei. From this we can conclude that the exclusion between the two modes is a mean-field effect. It also implies that atomic nuclei exhibit the unique possibility of exhibiting two different pairing condensates simultanously.
In order to explore further the mutual exclusiveness of the T=1 and T=0 pair-fields obtained in earlier studies, we have studied the HFB solution as a function of the strength of the T=0 interaction. The results are presented in Fig. 2. For the normal strength $`G_{T=0}=1`$, the solution corresponds to a T=1 pair-field. With increasing $`G_{T=0}`$ the HFB energy remains constant which is obvious since the solution has only the T=1 component and there is no T=0 component. The T=0 solution shown in Fig. 2 has been obtained by solving the HFB equations for a very large value of $`G_{T=0}`$ ($`G_{T=0}=2.8`$) and then using this solution for lower values of $`G_{T=0}`$. In this manner, it was possible to obtain a T=0 solution also below the critical point, see Fig. 2. We note from Fig. 2 that the two solutions coexist for most of the $`G_{T=0}`$ values. They represent two different solutions of the HFB equations.
The exact solution, presented in Figs. 1 contains both the T=0 and T=1 pair-modes, whereas HFB gives two separate solutions, corresponding to either T=0 or T=1 pair-fields. The difference between the two models resides in the fact that in the exact model, the two-body interaction always is a scalar whereas in the HFB-aproximation, the pairing potential is either a T=0 or T=1 field, with the corresponding symmetry. Our analysis shows that starting from a certain solution, with a given symmetry, this symmetry propagates to the next solution (with different $`G_{T=0}`$), analoguous to other self-consistent symmetries of the HFB hamiltonian, see e.g. the discussion in . The different pair-fields appear as independent of each other. Our results further indicate, that for a certain strength of the $`G_{T=0}`$ pair field, energy can be gained. This conforms with earlier results to associate the Wigner energy with T=0 pair correlations.
As a next step, we consider the response of the nuclear pair-potential to the rotating fields. In Fig. 3, we show the total pair-field (Eq. 15) as well as selected individual $`(J,T)`$ contribution as a function of the rotational frequency ( $`\mathrm{}\omega `$ ) for 4 particles (2 protons and 2 neutrons) in the $`f_{7/2}`$ shell. First of all, we may note the distinct difference between the T=1 and T=0 pairing fields. Whereas the T=1 field is dominated by one component with J=0, the T=0 mode is dominated by the J=1 and $`J=2j`$ part of the interaction, also the intermediate spins $`J=3,5`$ play a role. This already indicates that a discussion of a pairing force restricted to $`L=0`$ may be appropriate for the T=1 part of the interaction, but not for T=0, see also ref. .
As we increase the rotational frequency, the T=1 pairing-correlations (solid line) reveal the well known drop due to particle-alignment from the $`f_{7/2}`$-shell at around $`\mathrm{}\omega =0.7G`$. At this crossing point, the yrast band changes character from the paired $`(J=0)`$ configuration to the aligned $`(J=M_x=6+6)`$ state.
Similar calculations were performed also for the case of the (4+2) and (4+4) systems. Qualitatively, they all show the same trend, where of course the size of the drop in correlation energy depends on the number of particles present in the single-$`j`$ shell. For (4+2) system, the correlations of the J=0 component only for one-pair disappear whereas the drop for the 4+4 particles is less pronounced. This is due to the fact that only one-proton and one-neutron pair have aligned at the first crossing. Hence, the J=0 correlations are still active for the remaining two-pairs. For higher frequencies, the next pair will align, and then the J=0 (and in consequence) the T=1 correlations will drop in a similar fashion as for the system with one-proton and one-neutron pair only. The important message remains, as is evident from Fig. 1, that the T=1 field is largely built up from the J=0 pair-correlations, that are diminished in the process of particle alignment. Although, the components with higher-$`J`$ contribute at higher values of angular-momentum, the T=1 correlations are strongly reduced by the rotational motion.
In contrast, the T=0 correlations evolve quite differently with rotational frequency. The contribution of the coupling to low-$`J`$, like the J=1 pairs, behave similar to the coupling to J=0. This is quite natural, since they are built up by pairs of $`L=0`$ and $`L=2`$. However, although the contribution of the J=1 to the T=0 correlations drop in a similar fashion as the J=0, the value of the total T=0 correlations remain essentially unchanged. Apparently, the part that is lost by J=1 and $`J=3`$ is gained by $`J=7`$ and $`J=5`$. This implies, that the high-$`J`$ components of the T=0 correlations compensate the loss of the low-$`J`$. This feature appears to be independent of the number of particles in the system. It means that for a given interaction in the pp-channel, the total T=0 correlations remain almost unaffected by rotation. The presence of increasing $`L`$-values in the pairing field will affect deformation properties. This is what one expects in a fully self-consistent approach, which of course is beyond our present model analysis. Note that a recent analysis within the Monte Carlo Shell model shows that at high angular momenta, the T=0 correlations with $`2j`$ increase.
From the above analysis, one may conclude that the T=0 correlations are not able to affect rotational properties, since the increase in the stretched $`J=2j`$ component is exactly nullified by the decrease of the J=1 part, see also the discussion in Ref. . Indeed, these are the results e.g. for the $`f_{7/2}`$ shell where one is dealing with a โsingle-jโ shell. However, in heavier nuclei, when $`Z>28`$, the active shell is composed of, e.g., $`p{}_{3}{}^{}/2,f_{5/2},p_{1/2}`$ and $`g_{9/2}`$. For those cases, the J=1 part of the T=0 interaction becomes more coherent, since every subshell can contribute. In contrast, the $`J=2j`$ components become fragmented, since they have a different value for each subshell. Therefore, one may expect a different response of the T=0 pair field to rotation in heavier nuclei. Since we are dealing in our model with a single-j shell it is not possible to deal with such a case. One may, however, mimick in an adhoc way the region beyond $`Z=28`$ by increasing the strength of the J=1 part of the interaction.
The effect of a redistributed strength of the T=0 correlations, where the J=1 part has been increased by a factor of two, is shown in Fig. 5. Indeed, the crossing frequency is shifted. In other words, a coherence of J=1, T=0 pairs results in a change of the crossing frequency. What is even more striking is that this effect is suppressed when $`NZ`$. In Fig. 5, we show the case of (2+4) nucleons in the $`f_{7/2}`$ shell and, indeed, the first crossing frequency remains essentially unchanged. This feature persists also in the HFB approximation. Although our model is highly simplistic, one can certainly conclude that T=0, J=1 collectivity results in a shift of the crossing frequency to higher values and that this property is expected to be present also in more realistic calculations. Of course, as discussed above, there are other factors that affect the crossing frequency, like the deformation which in turn can be influenced by the T=0 pairing field.
A shift of the crossing frequency has been reported for the case of the N=Z nucleus <sup>72</sup>Kr. There have been efforts to explain this shift of the crossing in terms of T=1 np-pairing. Since to a very good approximation, the nuclear force is charge independent, only the total isospin $`T`$ matters for the interaction, not the projection of isospin ($`T_z`$). This is analoguous to the assumption that the nuclear force does not depend on the angular-momentum projection $`J_z`$, only on total $`J`$. This basic assumption implies that the T=1 pair-gaps are not affected by rotation in isospace, i.e. the total $`T=1`$ pair gap ($`\mathrm{\Delta }_{nn}^2+\mathrm{\Delta }_{pp}^2+\mathrm{\Delta }_{np}^2`$) is an invariant quantity. In an attempt of Ref. to account for the shift of the crossing frequency, the T=1 $`\mathrm{\Delta }_{np}`$ pair-gap was simply increased from 0 to a value of 2.5 MeV. Such an increase strongly violates charge independence. Following the arguments given above, one could as well increase the nn- or pp-pairing gaps. Of course, any increase of the T=1 pairing energy will result in a shift of the crossing frequency but this has nothing to do with np-pairing.
In summary, we have studied the competition between the T=0 and T=1 pair-fields in an exactly soluble deformed single-j shell model. It is shown that the HFB approach gives rise to two decoupled solutions corresponding to T=1 and T=0 modes. Although, in the exact shell model analysis, the solution contains both T=0 and T=1 modes, the two modes are independent with T=1 pair-energy independent of the strength of the T=0 correlations and vice-versa. The T=0 correlations in a single-j shell have a complex structure where the total amount is not affected by rotation. For realistic cases in heavy nuclei (Z$`>`$28), with several j-shells, the J=1 part will effectively acquire a larger strength. It has been demonstrated that increasing the value of the (T=0, J=1) pair-strength results in a shift of the bandcrossing frequency. Such a shift of the crossing frequency in heavy N=Z nuclei, therefore, is an indication of the collective (T=0, J=1) correlations.
This work has been supported by the Gรถran Gustafsson Foundation and the Swedish Natural Research Council (NFR).
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# Algebra of Principal Fibre Bundles, and Connections.
## 1 Principal Fibre Bundles
Let us consider a groupoid object $`\mathrm{\Phi }\text{}C`$ in a left exact category $`\underset{ยฏ}{E}`$. Let us also consider a subobject $`AC`$ and a global section $`:1C`$. We shall talk about $`\underset{ยฏ}{E}`$ as if it were the category of sets, so we may say โsubsetโ instead of โsubobjectโ. We assume that the domain- and codomain formation maps are effective descent maps in $`\underset{ยฏ}{E}`$, and that the groupoid is transitive, meaning that โthe anchorโ map $`<d_0,d_1>:\mathrm{\Phi }C\times C`$ is also an effective descent map.
Then the set $`P`$ of those arrows of $`\mathrm{\Phi }`$, whose codomain is in $`A`$ and whose domain is $``$, carries the structure of a principal fibre bundle over $`A`$, with group $`G=\mathrm{\Phi }(,)`$. Any principal fibre bundle in $`\underset{ยฏ}{E}`$ comes about this way from a groupoid (see the remarks below). The algebraic structure of $`P`$ comes about from that of $`\mathrm{\Phi }`$, and may be made explicit as follows. First, codomain formation $`d_1:\mathrm{\Phi }C`$ resticts to a map $`\pi :PA`$, which is the structural map of the bundle. The group $`G=\mathrm{\Phi }(,)`$ acts from the right on $`P`$, by precomposition (we compose from right to left). Clearly, this action is free, and transitive on the fibres $`\pi ^1(a)`$ ($`aA`$).
Any element $`g`$ of $`G=\mathrm{\Phi }(,)`$ may, for any $`aA`$, be written in the form $`x^1z`$ for a pair of elements in $`\pi ^1(a)`$. This representation of elements in $`G`$ by โfractionsโ $`x^1z`$ prompts us to use the (Ehresmann) notation $`P^1P`$ for $`G`$. Then clearly $`xx^1z=z`$ (where $``$ denotes the $`G`$-action). Any choice of $`x\pi ^1(a)`$ provides us with an explicit bijection $`\pi ^1(a)G`$, given by $`zx^1z`$.
Let us also consider the (transitive) subgroupoid $`\mathrm{\Phi }_A`$ of $`\mathrm{\Phi }`$ consisting of those arrows whose domain and codomain both belong to $`A`$. The groupoid $`\mathrm{\Phi }_A`$ acts on the left on $`PA`$ by postcomposition; any arrow $`ab`$ in it may be presented in the form $`yx^1`$, for some $`x\pi ^1(a)`$ and $`y\pi ^1(b)`$. Then clearly $`yx^1x=y`$ (where $``$ denotes the action). The representation of arrows in $`\mathrm{\Phi }_A`$ by โfractionsโ $`yx^1`$ prompts us to use the (Ehresmann) notation $`PP^1`$ for $`\mathrm{\Phi }_A`$.
Remark. The set $`P`$ itself carries a partially defined ternary operation, given by the composite $`yx^1z`$ in $`\mathrm{\Phi }`$ (defined subject to the book-keeping condition that $`\pi (x)=\pi (z)`$), and this operation satisfies a couple of equations and book-keeping conditions, making it into a โpregroupoidโ on $`A`$, in the sense of . Out of such pregroupoid, a transitive groupoid $`\mathrm{\Phi }`$ on $`A+1`$ may be constructed, which in turn gives rise to $`P`$ by the procedure described above (provided $`\pi :PA`$ is an effective descent map); this is in essence demonstrated in . Principal fibre bundles (in the classical sense) $`PA`$, in the category of smooth manifolds, say, may, in a rather evident way, be provided with pregroupoid structure. So our โgroupoid theoreticโ way of describing the notion of principal fibre bundle subsumes the classical notion, and it is essentially Ehresmannโs conception. โ A (non-transitive) generalization where $`:1C`$ is replaced by a subset $`BC`$, is considered in ; this generalization is relevant for foliation theory, cf. loc.cit. and .
We shall henceforth be interested in the case where the โbaseโ $`A`$ of the bundle $`PA`$ is to be thought of as a manifold, so we denote it by $`M`$ rather than by $`A`$.
Remark on fibre bundles in general. A principal fibre bundle $`PM`$ with group $`G`$, may by the above be identified with a groupoid $`\mathrm{\Phi }`$ with set of objects $`M+1`$, (and with $`G=\mathrm{\Phi }(,)`$, where $``$ is the isolated point of $`M+1`$). Similarly, a fibre bundle $`\pi :EM`$, with associated principal bundle $`P`$ and with fibre a left $`G`$-set $`F`$, becomes identified with a discrete opfibration over $`\mathrm{\Phi }`$ (in the algebraic sense, i.e. an action by $`\mathrm{\Phi }`$), with $`F=\pi ^1()`$, and $`E=\pi ^1(M)`$. Such fibre bundle is determined up to isomorphism by $`P`$ and $`F`$ (with its left $`G`$-action). In the present general context, this is the upshot of . We shall not explicitly be using this correspondence for general fibre bundles here. But let us remark that $`P`$ itself is a fibre bundle, with fibre $`G`$ (with $`G`$-action by left multiplication). Likewise, if we let $`G`$ act on $`G`$ by conjugation, $`gh:=ghg^1`$, we get a group bundle, namely what calls the gauge group bundle of $`PP^1`$. (It is also known under the name $`Ad(P)`$.) We shall utilize this latter bundle, but shall recall it without reference to this general fibre bundle theory. For a groupoid $`\mathrm{\Psi }\text{}M`$, the gauge group bundle $`\text{gauge}(\mathrm{\Psi })`$ is a bundle over $`M`$, which for its fibre over $`aM`$ simply has the group
$$(\text{gauge}(\mathrm{\Psi }))_a=\mathrm{\Psi }(a,a).$$
It carries a left action by $`\mathrm{\Psi }`$, given by conjugation: if $`f:ab`$ in $`\mathrm{\Psi }`$ and $`h\mathrm{\Psi }(a,a)`$, then $`fhf^1\mathrm{\Psi }(b,b)`$.
The existence, for any principal fibre bundle $`P`$, of an embedding of $`P`$ into a groupoid $`\mathrm{\Phi }`$, implies a โmetatheoremโ, namely that we may calculate freely with expressions, like $`vu^1`$, as if we were dealing with actual compositions in a groupoid. The โactionโ dots, like in $`yx^1x`$ are then superfluous, and the same applies to many parentheses; so they are mainly kept for readability. The message (which I also tried to get through in and in several other places) is that a fair amount of calculations in geometry can be performed on this very basic โmultiplicativeโ level.
Since an arrow $`f:ab`$ in the groupoid $`PP^1`$ may be represented as a โfractionโ $`yx^1`$ (with $`yP_b`$ and $`xP_a`$), it follows that an element $`h`$ over $`a`$ in the gauge group bundle $`\text{gauge}(PP^1)`$ may be represented by a fraction $`yx^1`$ with $`y`$ and $`x`$ both $`P_a`$. For the case where the group $`G`$ is commutative, it is well known, and easy to see, that we have an isormorhism of group bundles
$$\text{gauge}(PP^1)M\times G,$$
(1)
given by sending $`h=yx^1PP^1`$ to $`x^1yP^1P`$. This cannot be done for non-commutative $`G`$: for any $`gG`$, the same $`h`$ may also be represented by the fraction $`yg(xg)^1`$, but $`(xg)^1(yg)=g^1(x^1y)g`$ which is not equal to $`x^1y`$ in general.
## 2 Connections versus connection forms
Consider a principal bundle $`\pi :PM`$, with group $`G`$, as above. We shall assume that $`M`$ and $`P`$ are equipped with reflexive symmetric relations $``$, called the neighbour relation. The set of pairs $`(x,y)M\times M`$ with $`xy`$ is a subset $`M_{(1)}M\times M`$, called the first neighbourhood of the diagonal, and similarly for $`P_{(1)}P\times P`$. We assume that $`\pi :PM`$ preserves the relation $``$, and also that it is an โopen submersionโ in the sense that if $`ab`$ in $`M`$, and $`\pi (x)=a`$, then there exists a $`yx`$ in $`P`$ with $`\pi (y)=b`$. In fact, we assume that for any โinfinitesimal $`k`$-simplexโ $`a_0,\mathrm{},a_k`$ in $`M`$ (meaning a $`k+1`$-tuple of mutual neighbours), and for any $`x_0P`$ above $`a_0`$, there exists an infinitesimal $`k`$-simplex $`x_0,\mathrm{},x_k`$ in $`P`$ (with the given first vertex $`x_0`$) which by $`\pi `$ maps to $`a_0,\mathrm{},a_k`$. Finally. the action of any $`gG`$ on $`P`$ is assumed to preserve the relation $``$ on $`P`$.
This is motivated by Synthetic Differential Geometry (SDG), cf. , and more recently , where the notion of connection (infinitesimal parallel transport) and differential form is elaborated in these terms.
The groupoid viewpoint for connections is also in essence due to Ehresmann. In SDG, this connection notion becomes paraphrased (see , or , Section 8): for a groupoid $`\mathrm{\Phi }\text{}M`$, a connection in it is just a map $`:M_{(1)}\mathrm{\Phi }`$ of reflexive symmetric graphs over $`M`$.
Let $`\pi :PM`$ be a principal fibre bundle. To any connection $``$ in the groupoid $`PP^1`$, one may associate a 1-form $`\omega `$ on $`P`$ with values in the group $`P^1P`$, as follows. For $`u`$ and $`v`$ neighbours in $`P`$, with $`\pi (u)=a`$, $`\pi (v)=b`$, put
$$\omega (u,v):=u^1((a,b)v).$$
(2)
Note that both $`u`$ and $`(a,b)v`$ are in the $`\pi `$-fibre over $`a`$, so that the โfractionโ $`u^1((a,b)v)`$ makes sense as an element of $`P^1P`$.
The defining equation is equivalent to
$$u\underset{P^1P}{\underset{}{\omega (u,v)}}=\underset{PP^1}{\underset{}{(\pi (u),\pi (v))}}v.$$
(3)
If we agree that (for $`u,v`$ in $`P`$ a pair of neighbours in $`P`$) $`(u,v)`$ denotes $`(\pi (u),\pi (v))`$, this equation may be written more succinctly
$$u\omega (u,v)=(u,v)v.$$
(4)
It is possible to represent the relationship between $``$ and the associated $`\omega `$ by means of a simple figure:
The figure reflects something geometric, namely that $`\omega (u,v)`$ acts inside the fibre (vertically), whereas $``$ defines a notion of horizontality.
We have the following two equations for $`\omega `$. First, let $`xy`$ in $`P`$, and assume that $`g`$ has the property that also $`xgy`$. Then
$$\omega (xg,y)=g^1\omega (x,y).$$
(5)
Also, for $`xy`$ and any $`gG`$
$$\omega (xg,yg)=g^1\omega (x,y)g.$$
(6)
To prove (5), let us denote $`\pi (x)=\pi (xg)`$ by $`a`$ and $`\pi (y)`$ by $`b`$. Then we have, using the defining equation (3) for $`\omega `$ twice,
$$xg\omega (xg,y)=(a,b)y=x\omega (x,y),$$
and now we may calculate as in a groupoid: first cancel the $`x`$ on the left, then multiply the equation by $`g^1`$ on the left. To prove (6), we have, with $`a`$ and $`b`$ as above,
$$xg\omega (xg,yg)=(a,b)yg=x\omega (x,y)g,$$
by the defining equation (3) for $`\omega (xg,yg)`$, and by (3) for $`\omega (x,y)`$, multiplied on the right by $`g`$, respectively. From this, we get the result by first cancelling $`x`$ and then multiplying the equation by $`g^1`$ on the left.
The following Proposition is now the rendering, in our context, of the relationship between a connection $``$ and its connection 1-form $`\omega `$:
###### Proposition 1
The process $`\omega `$ just described, establishes a bijective corresondence between 1-forms $`\omega `$ on $`P`$, with values in the group $`P^1P`$ and satisfying (5) and (6), and connections $``$ in the groupoid $`PP^1`$.
Proof. Given a 1-form $`\omega `$ satisfying (5) and (6), we construct a connection $``$ as follows. Let $`ab`$ in $`M`$. To define the arrow $`(a,b)`$ in $`PP^1`$, pick $`uv`$ above $`ab`$, and put
$$(a,b)=u(v\omega (v,u))^1.$$
We first argue that this is independent of the choice of $`v`$, once $`u`$ is chosen. Replacing $`v`$ by $`vgu`$, we are in the situation where (5) may be applied; we get
$$u(vg\omega (vg,u)^)1=u(vgg^1\omega (v,u))^1=u(v\omega (v,u))^1;$$
the left hand side is $`(a,b)`$ defined using $`u,vg`$, the right hand side is using $`u,v`$.
To prove independence of choice of $`u`$: any other choice is of form $`ug`$ for some $`gG`$. For our new $`v`$, we now chose $`vg`$ (the result will not depend on the choice, by the argument just given). Again we calculate. By (6), we have the first equality sign in
$$ug(vg\omega (vg,ug))^1=ug(vgg^1\omega (u,v)g)^1=ug(v\omega (u,v)g)^1=u(v\omega (u,v))^1,$$
and the two expressions here are $`(a,b)`$ defined using, respectively, $`ug,vg`$ and $`u,v`$.
The calculation that the two processes are inverse of each other is trivial (using $`\omega (u,v)=\omega (v,u)^1`$ and $`(a,b)=(b,a)^1`$).
## 3 Gauge forms versus horizontal equivariant forms
We consider a principal fibre bundle $`\pi :PM`$ as in the previous section. The horizontal $`k`$-forms that we now consider, are $`k`$-forms on $`P`$ with values in the group $`G=P^1P`$. Horizontality means for a $`k`$-form $`\theta `$ that
$$\theta (u_0,u_1,\mathrm{},u_k)=\theta (u_0,u_1g_1,\mathrm{},u_kg_k)$$
(7)
for any infinitesimal $`k`$-simplex $`(u_0,u_1,\mathrm{},u_k)`$ in $`P`$, and any $`g_1,\mathrm{}g_kP^1P`$ with the property that $`(u_0,u_1g_1,\mathrm{},u_kg_k)`$ is still an infinitesimal simplex (which is a strong โsmallnessโ requirement on the $`g_i`$โs).
Note that the connection form $`\omega `$ for a connection $``$ is not a horizontal 1-form, since $`\omega (x,yg)=\omega (x,y)g`$, not $`=\omega (x,y)`$.
We say that a $`k`$ form $`\theta `$, as above, is equivariant if for any infinitesimal $`k`$-simplex $`(u_0,\mathrm{},u_k)`$, and any $`gP^1P`$, we have
$$\theta (u_0g,u_1g,\mathrm{},u_kg)=g^1\theta (u_0,u_1,\mathrm{},u_k)g.$$
(8)
Note that connection forms are equivariant in this sense, by (6).
###### Proposition 2
Assume that the group $`G=P^1P`$ is commutative. Then any horizontal equivariant $`k`$-form $`\theta `$ on $`P`$ can be written $`\pi ^{}(\mathrm{\Theta })`$ for a unique $`G`$-valued $`k`$-form $`\mathrm{\Theta }`$ on the base space $`M`$.
Proof. It is evident that any form $`\pi ^{}(\mathrm{\Theta })`$ is horizontal and equivariant (which here is better called invariant, since the equivariance condition now reads $`\theta (u_0g,u_1g,\mathrm{},u_kg)=\theta (u_0,u_1,\mathrm{},u_k)`$). Conversely, given an equivariant (= invariant) $`k`$-form $`\theta `$ on $`P`$, and given an infinitesimal $`k`$-simplex $`a_0,\mathrm{},a_k`$ in $`M`$, define
$$\mathrm{\Theta }(a_0,\mathrm{},a_k):=\theta (x_0,\mathrm{},x_k)$$
where $`x_0,\mathrm{},x_k`$ is any infinitesimal $`k`$-simplex above $`a_0,\mathrm{},a_k`$. The proof that this value does not depend on the choice of the $`x_i`$โs proceeds much like the proof of the well-definedness of a connection given a connection-form, in Proposition 1 above: First we prove, for fixed $`x_0`$ above $`a_0`$, that the value is independent of the choice of the remaining $`x_i`$โs, and this is clear from the verticality assumption on $`\theta `$. Next we prove that changing $`x_0`$ to $`x_0g`$ (and picking $`x_1g,\mathrm{},x_kg`$ for the remaining vertices in the new $`k`$-simplex) does not change the value either, and this is clear from equivariance (= invariance).
Recall that a $`k`$-form with values in a group bundle $`EM`$ associates to an infinitesimal $`k`$-simplex $`a_0,\mathrm{},a_1`$ in $`M`$ an element in the fibre of $`E_{a_0}`$. We are interested in the case where $`E`$ is the gauge group bundle of a groupoid; such forms we call gauge forms, for brevity.
###### Proposition 3
There is a natural bijective correspondence between horizontal equivariant $`k`$-forms on $`P`$ with values in $`G=P^1P`$, and $`k`$-forms on $`M`$ with values in the gauge group bundle $`\text{gauge}(PP^1)`$.
Proof/Construction. Given a horizontal equivariant $`k`$-form $`\theta `$ on $`P`$ as above, we construct a gauge valued $`k`$-form $`\stackrel{ห}{\theta }`$ on $`M`$ by the formula
$$\stackrel{ห}{\theta }(a_0,\mathrm{},a_k):=(u_0\theta (u_0,\mathrm{},u_k))u_0^1,$$
(9)
where $`(u_0,\mathrm{},u_k)`$ is an arbitrary infinitesimal $`k`$-simplex mapping to the infinitesimal $`k`$-simplex $`(a_0,\mathrm{},a_k)`$ by $`\pi `$ (such exist, since $`\pi `$ is a surjective submersion). Note that the enumerator and the denominator in the fraction defining the value of $`\stackrel{ห}{\theta }`$ are both in the fibre over $`x_0`$, so that the value is an endo-map at $`a_0`$ in the groupoid $`PP^1`$, thus does belong to the gauge group bundle. โ We need to argue that this value does not depend on the choice of the infinitesimal simplex $`(u_0,\mathrm{}u_k)`$. We first argue that, once $`u_0`$ is chosen, the choice of the remaining $`u_i`$โs in their respective fibres does not change the value. This follows from (7). To see that the value does not depend on the choice of $`u_0`$: choosing another one amounts to choosing some $`u_0g`$, for some $`g`$. But then we just change $`u_1,\mathrm{},u_k`$ by the same $`g`$; this will give the arrow in $`PP^1`$
$$(u_0g\theta (u_0g,\mathrm{},u_kg))(u_0g)^1.$$
Now we calculate using the โmetatheoremโ, so we drop partentheses and multiplication dots; using the assumed equivariance (8), this expression then yields
$$u_0gg^1\theta (u_0,\mathrm{},u_k)gg^1u_0^1,$$
which clearly equals the expression in (9).
Conversely, given a gauge valued $`k`$-form $`\alpha `$ on $`M`$, we construct a $`P^1P`$-valued $`k`$-form $`\widehat{\alpha }`$ on $`P`$ by putting
$$\widehat{\alpha }(u_0,u_1,\mathrm{},u_k):=u_0^1(\alpha (a_0,a_1,\mathrm{},a_k)u_0)$$
(10)
where $`a_i`$ denotes $`\pi (u_i)`$. Since, for $`i1`$, this expression depends on $`u_i`$ only through $`\pi (u_i)=a_i`$, it is clear that (7) holds, so the form $`\widehat{\alpha }`$ is horizontal. Also,
$$\widehat{\alpha }(u_0g,\mathrm{},u_kg)=(u_0g)^1(\alpha (a_0,\mathrm{},a_k)(u_0g));$$
by the metatheorem, this immediately calculates to the expression in (10).
Finally, a calculation with the metatheorem again (cancelling $`u_0^1`$ with $`u_0`$) immediately gives that the two processes $`\theta \stackrel{ห}{\theta }`$ and $`\alpha \widehat{\alpha }`$ are inverse to each other.
We may summarize the bijection $`\alpha \widehat{\alpha }`$ from $`\text{gauge}(PP^1)`$-valued forms on $`M`$ to horizontal equivariant $`P^1P`$-valued forms on $`P`$ by the formula
$$u_0\widehat{\alpha }(u_0,,,,,u_k)=(\pi ^{}\alpha )(u_0,\mathrm{},u_k)u_0.$$
(11)
In the case that the group $`G=P^1P`$ is commutative, we may cancel the โexternalโ $`u_0`$โs, and get
$$\widehat{\alpha }(u_0,\mathrm{},u_k)=(\pi ^{}\alpha )(u_0,\mathrm{},u_k),$$
for all infinitesimal $`k`$-simplices $`u_0,\mathrm{},u_k`$. So under the identification of gauge forms with $`G`$-valued forms implied by (1), we have that
$$\widehat{\alpha }=\pi ^{}\alpha .$$
(12)
Recall that if $``$ and $`_1`$ are two connections in a groupoid $`\mathrm{\Phi }\text{}M`$, we may form a 1-form $`_1^1`$ with values in the gauge group bundle; it is given by
$$_1^1(a,b)=_1(a,b)(b,a).$$
For the case where the groupoid is $`PP^1`$, we have the following Proposition, which we shall not use in the sequel, but include for possible future reference:
###### Proposition 4
Let $`PM`$ be a principal bundle, and let $``$ and $`_1`$ be two connections in the groupoid $`PP^1`$. Then
$$(_1^1)\widehat{}=\omega _1\omega ^1$$
where $`\omega `$ and $`\omega _1`$ are the connection forms of $``$ and $`_1`$, respectively.
Proof. Let $`xy`$, over $`a`$ and $`bM`$, respectively. Then
$`(_1^1)\widehat{}(x,y)`$ $`=`$ $`x^1(_1(a,b)(b,a)x)`$
$`=`$ $`x^1_1(a,b)y\omega (y,x)`$
$`=`$ $`x^1x\omega _1(x,y)\omega (y,x)`$
$`=`$ $`\omega _1(x,y)\omega (y,x)`$
$`=`$ $`(\omega _1\omega ^1)(x,y),`$
using the defining relation (11) for $`()\widehat{}`$, and the relation (3) for $``$ and $`_1`$, respectively.
## 4 Curvature versus coboundary
Recall that the curvature of a connection in a groupoid $`\mathrm{\Phi }\text{}M`$ is the $`\text{Gauge}(\mathrm{\Phi })`$-valued 2-form $`R=R_{}`$ given by
$$R(a_0,a_1,a_2)=(a_0,a_1)(a_1,a_2)(a_2,a_0),$$
and recall that if $`\omega `$ is a 1-form with values in a group $`G`$, them $`d\omega `$ is the $`G`$-valued 2-form given by
$$d\omega (x_0,x_1,x_2)=\omega (x_0,x_1)\omega (x_1,x_2)\omega (x_2,x_0).$$
We apply this to the case where $`\mathrm{\Phi }=PP^1`$ and $`G=P^1P`$, for a principal fibre bundle $`\pi :PM`$. Then the curvature $`R`$, which is a $`\text{gauge}(PP^1)`$ -valued 2-form on $`M`$, gives, by Proposition 3, rise to a (horizontal and equivariant) $`P^1P`$-valued 2-form $`\widehat{R}`$ on $`P`$.
We then have the following:
###### Theorem 1
Let $`\pi :PM`$ be a principal fibre bundle, and let $``$ be a connection in the groupoid $`PP^1`$ with connection form $`\omega `$. Then we have an equality of $`P^1P`$-valued 2-forms on $`P`$:
$$\widehat{R}=d\omega .$$
In particular, $`d\omega `$ is horizontal and equivariant.
The form $`\widehat{R}=d\omega `$ is the curvature form of the connection. See the remark below for comparison with the classical formulation.
Proof. Let $`x,y,z`$ be an infinitesimal 2-simplex in $`P`$, and let $`a=\pi (x)`$, $`b=\pi (y)`$, and $`c=\pi (z)`$. We calculate the effect of the (left) action of the arrow $`R(a,b,c)`$ on $`x`$ (note that $`R(a,b,c)`$ is an endo-arrow at $`a`$ in the groupoid):
$`R(a,b,c)x`$ $`=`$ $`(a,b)(b,c)(c,a)x`$
$`=`$ $`(a,b)(b,c)z\omega (z,x)`$
$`=`$ $`(a,b)y\omega (y,z)\omega (z,x)`$
$`=`$ $`x\omega (x,y)\omega (y,z)\omega (z,x)`$
$`=`$ $`xd\omega (x,y,z),`$
using the defining equations for $`R`$ and for $`dw`$ for the two outer equality signs, and using (3) three times for the middle three ones. This proves the Theorem.
Remark. By I.18, or in more detail, ), there is a bijective correspondence between $`G`$-valued $`k`$-forms $`\theta `$ on a manifold $`P`$ (where $`G`$ is a Lie group, say $`P^1P`$), and differential $`k`$-forms $`\overline{\theta }`$, in the classical sense, with values in the Lie algebra $`\underset{ยฏ}{g}`$ of $`G`$ (i.e. multilinear alternating maps $`TP\times _P\mathrm{}\times _PTP\underset{ยฏ}{g}`$. Under this correspondence, the horizontal equivariant 2-form $`d\omega `$ considered in the Theorem corresponds to the classically considered โcurvature 2-formโ $`\mathrm{\Omega }`$ on $`P`$, as in II.4, 5.3, or V bis 4, (perhaps modulo a factor $`\pm 2`$, depending on the conventions chosen). This is not completely obvious, since $`\mathrm{\Omega }`$ differs from the exterior derivative $`d\overline{\omega }`$ of the classical connection form $`\overline{\omega }`$ by a โcorrection termโ $`1/2[\overline{\omega },\overline{\omega }]`$ involving the Lie Bracket of $`\underset{ยฏ}{g}`$; or, alternatively, the curvature form comes about by modifying $`d\overline{\omega }`$ by a โhorizontaliztion operatorโ (this โmodificationโ also occurs in the treatment in ). The fact that this โcorrection termโ (or the โmodificationโ) does not come up in our context can be explained by Theorem 5.4 in (or see Theorem 18.5); here it is proved that the formula $`d\omega (x,y,z)=\omega (x,y)\omega (y,z)\omega (z,x)`$ already contains this correction term, when translated into โclassicalโ Lie algebra valued forms.
For the case where the group $`P^1P`$ is commutative, we may use the isomorphism (1) to identify $`\text{ gauge}(PP^1)`$-valued forms on $`M`$ with $`P^1P`$-valued forms on $`M`$. Also, by Proposition 2, and the horizontality and equivariance of $`d\omega `$, there is a unique $`P^1P`$-valued 2-form $`\mathrm{\Omega }=\stackrel{ห}{d\omega }`$ on $`M`$ with $`\pi ^{}(\mathrm{\Omega })=d\omega `$. We therefore have the following Corollary (notation as above), which is essentially what call the infinitesimal version of Gauss-Bonnet Theorem (for the case where $`G=SO(2)`$):
###### Corollary 1
Assume $`P^1P`$ is commutative, and let the connection $``$ in $`PP^1`$ have connection form $`\omega `$. Then the unique $`G`$-valued 2-form $`\mathrm{\Omega }`$ on $`M`$ with $`\pi ^{}\mathrm{\Omega }=d\omega `$ is $`R_{}`$.
$$R_{}=\mathrm{\Omega }.$$
Let us remark that also gives a version of the Corollary for the non-commutative case, their Proposition 6.4.1; this, however, seems not correct. In this sense, our Theorem 1 is partly meant as a correction to Prop. 6.4.1, partly a โtranslationโ of it into the pure multiplicative fibre bundle calculus, which is our main concern.
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# Higher-order QED corrections to single-๐ production in electron-positron collisions
## 1 Introduction
Four-fermion final states are of special interest for the physics programme of lep2 and future high-energy electron-positron colliders, being entangled with electroweak gauge boson production and decay . In particular, the process considered in the present paper, i.e.
$$e^+e^{}e^{}(e^+)\overline{\nu }_e(\nu _e)q^{}\overline{q}$$
(1)
is peculiar among all the possible four-fermion final states because the bulk of its cross section is due to two sub-processes, i.e. $`W`$-boson pair production and decay
$$e^+e^{}W^{}W^{}4\mathrm{fermions}$$
(2)
and the radiation of an almost on shell $`t`$-channel photon from the electron (positron), with subsequent production of a $`W`$-boson and a neutrino
$$e^+e^{}\gamma ^{}e^+(e^{})W^{}\overline{\nu }_e(e^{})4\mathrm{fermions}$$
(3)
Despite, strictly speaking, both sub-processes (2,3) always occur simultaneously and are indistinguishable, channel (2) dominates if the electron is emitted at large angle, whereas channel (3) dominates if the electron is emitted in the very forward region, because of the presence of a quasi-real $`t`$-channel photon.
In this paper the process (3) will be addressed, by restricting the analysis to the kinematical range of forwardly emitted electrons. This signal is usually referred to as single-$`W`$ production, since only the two final jets are detected .
The importance of this process has been emphasized since long time. In the lep2 energy range it is fundamental in order to study the self-interaction of the gauge bosons, together with the process (2), whereas in the energy regime of future colliders at the TeV scale it becomes the dominant electroweak process. In refs. cross sections and distributions were calculated in the approximation of real $`W`$-boson production, either by studying the reaction $`e^+e^{}e^{}\overline{\nu }_eW^+`$, or by employing the Weizsรคcker-Williams equivalent-photon approximation for the $`t`$-channel photon. In refs. it was pointed out the relevance of this process for the study of trilinear gauge boson couplings and some assessment of the sensitivity has been given. The first full four-fermion calculation, including the crucial effect of fermion masses, has been presented in ref. , where the lep2 sensitivity to anomalous gauge couplings has been studied. Since then, other complete four-fermion calculations of the single-$`W`$ process have appeared in the literature and implemented in computational tools for data analysis . In most of these calculations the effect of fermion masses is exactly accounted for in the dynamics and kinematics for the whole four-fermion phase-space , while in the approach of ref. the Weizsรคcker-Williams approximation is employed in the very forward, collinear region and matched with a massless four-fermion computation outside it. An up-to-date inventory of the present theoretical status is under preparation by the four-fermion working group of the lep2 mc workshop at cern .
Measurements of the single-$`W`$ cross section and the corresponding bounds on anomalous gauge couplings have been recently reported by the lep collaborations . Because the foreseen accuracy of final lep2 data is of the order of $`1`$-$`2\%`$ , accurate theoretical predictions for cross section and distributions are required.
The calculation of the cross section for single-$`W`$ processes poses several non-trivial theoretical problems . For a realistic account of gauge bosons properties it is mandatory to include the gauge boson width in the propagator. In general this mixes a fixed order calculation with an all order resummation of a class of Feynman diagrams and introduces a violation of the Ward identities of the theory. This issue is of special importance here since, due to the $`t`$-channel photon exchange, even a tiny violation of qed Ward identities is enhanced by a factor of $`s/m_e^2`$. This is indeed the case if a running width is used in the calculation. This problem has been extensively studied , and several options to address it have been explored. The most theoretically appealing procedure is the fermion loop scheme , which preserves both $`U(1)`$ and $`SU(2)`$ Ward identities. Recently, this scheme has been generalized to the case of massive external fermions, both in its minimal version, which considers the imaginary parts of the fermionic loops (ifl, and in its full realization with real and imaginary parts . In particular, in ref. a detailed numerical investigation has been performed, showing no significant difference between the ifl and the fixed width scheme, even in the region most sensible to $`U(1)`$ gauge invariance. For this reason, the fixed width scheme is adopted in the present calculation.
Another delicate issue is the so called resolved-photon component of the cross section. The quasi-real $`t`$-channel photon can split into a pair of almost massless quarks, leading to a situation where the partonic picture of hadrons is inadequate and both perturbative and non-perturbative qcd corrections become relevant. This issue is widely discussed in the literature , where the standard approach to this problem is also described and to which the reader is referred for details. However, for single-$`W`$ like events the resolved-photon component does not constitute a severe limitation: once a hard $`q\overline{q}`$ invariant mass cut is imposed, as done in realistic data analysis, the bulk of the signal is kept, whereas the resolved photon contribution becomes almost negligible .
A further relevant issue is given by radiative corrections due to photon radiation. Because exact $`O(\alpha )`$ electroweak corrections to single-$`W`$ production are still unknown, in most of the theoretical and experimental studies presented insofar only the large contribution of initial-state radiation (isr) has been taken into account, generally by using collinear structure function (hereafter denoted as sf) and assuming $`s=4E^2`$ as the proper scale for qed radiation. Due to the dominance of the quasi-real $`t`$-channel photon exchange, this can be expected not to be a suitable choice in the present case. On the other hand, it has been recently proposed to correct only the $`s`$-channel contributions to the single-$`W`$ signature, fixing the radiation scale in the usual manner, and to neglect the photonic corrections to the $`t`$-channel contributions . Following previous investigations of the pattern of photonic radiation in qed and electroweak processes, some theoretical arguments to determine the appropriate energy scale entering the sf are presented and compared with existing results. The analysis here described elucidates the theoretical details and provides further numerical results of a contribution by the authors to the activity of the four-fermion working group of the lep2 mc workshop at cern . Ideas similar to those adopted in the present work have recently appeared in ref. and there applied to the two-photon process $`e^+e^{}e^+e^{}\mu ^+\mu ^{}`$.
The paper is organized as follows. After a short review of the sf approach to leading log (LL) qed corrections in Sect. 2, Sect. 3 collects the analytical results valid for soft and collinear corrections to a generic scattering process. By comparing the results of Sect. 2 and Sect. 3, the radiation scales for single-$`W`$ production are determined in Sect. 4. Sect. 5 deals with the problem of taking into account the effect of the photon vacuum polarization in the single-$`W`$ process, while Sect. 6 shows the numerical results of the present study obtained with a Monte Carlo (hereafter mc) code for the single-$`W`$ signature, including also the effect of anomalous trilinear gauge couplings. Conclusions and prospects are given in Sect. 7.
## 2 Structure Function approach to photon radiation
Since in high-energy processes the corrections due to the emission of soft and collinear radiation are quite large, the LL contribution must be calculated at every perturbative order. A common technique to achieve this goal is the qed structure function approach , which consists in convoluting the hard-scattering cross section with appropriate โpartonโ densities. As well known, these convolution factors, i.e. the qed structure functions, include, by construction, both the real and virtual part of the photon correction, in order to ensure the cancellation of the infrared singularities. If a generic Born-level prediction $`d\sigma _0`$ is considered, the cross section $`d\sigma `$ including LL qed radiative corrections is obtained, by virtue of factorization theorems, according to the following general formula
$$d\sigma =\underset{i}{}๐x_iD(\mathrm{\Lambda }^2,x_i)๐\sigma _0$$
(4)
where $`1x_i`$ are the energy fractions carried away by the radiated photons from the $`i`$-th leg, $`\mathrm{\Lambda }`$ is the characteristic scale of the sf $`D(\mathrm{\Lambda }^2,x_i)`$, whose evolution is driven by the Dokshitzer-Gribov-Lipatov-Altarelli-Parisi (DGLAP) equation and is dependent on $`\mathrm{\Lambda }`$. It is worth noticing that the choice of the scale $`\mathrm{\Lambda }`$ is not dictated by general arguments and it is therefore rather arbitrary.
Equation (4) can be rewritten by stressing the possibility of different scales for each sf as follows
$$d\sigma =\underset{i}{}๐x_iD(\mathrm{\Lambda }_i^2,x_i)๐\sigma _0$$
(5)
In particular, if the integrated hard-scattering cross section is a smooth function of the centre of mass (c.m.) energy, once the integrations over the energy fractions $`x_i`$ are performed in the soft-photon approximation, the $`O(\alpha )`$ double-log expansion of eq. (5) can be written as follows
$$d\sigma =d\sigma _0\left(1+\underset{i}{}\frac{\alpha }{\pi }\mathrm{log}\frac{\mathrm{\Delta }E}{E}L(\mathrm{\Lambda }_i^2)\right)$$
(6)
where $`\mathrm{\Delta }E/E`$ is the maximum energy for undetected photons, to be identified with finite energy resolution of the photon detector, and $`L(\mathrm{\Lambda }_i^2)\mathrm{log}(\mathrm{\Lambda }_i^2/m^2)`$ is the collinear logarithm.
Since the functional form of the qed sf is accurately known , the main problem in evaluating eq. (5) is to fix the process scales $`\mathrm{\Lambda }_i`$. A generally adopted attitude is given by the comparison of eq. (6) with a perturbative calculation, which can be performed within any approximation, provided it reproduces the correct double-log contribution of the $`O(\alpha )`$ correction. This issue is addressed in the next Section.
## 3 Analytical results
The double-log contribution to photon radiation traces back to soft and collinear bremsstrahlung and its virtual counterpart , and, in the case of a calorimetric measurement of the energy of the final-state (fs) particles, to hard radiation collinear to the fs particles themselves . At each perturbative order, the leading contribution can be expanded in terms of infrared and collinear logs. For example, when the photons are emitted from the initial-state (is) particles only, such an expansion can be arranged in terms of double-log contributions of the form $`\alpha ^nl^nL_s^n`$, where $`l\mathrm{log}(\mathrm{\Delta }E/E)`$ is the infrared log and $`L_s\mathrm{log}(s/m^2)`$ is the collinear log. This is the reason why $`\mathrm{\Lambda }^2=s`$ is the โnaturalโ energy scale to be used for sf in the presence of isr only. When also fs radiation is considered, the collinear log, $`L`$, is in general modified by additional factors coming from the angular integration over the photon variables. A typical example is given by the radiation emitted from one leg in the $`t`$-channel qed contribution to Bhabha scattering . In the soft-photon approximation the radiation cones, one from the is electron and one from the fs electron, have a half-opening which is determined by the angle between the emitting particles, because of a destructive interference. As a consequence, the energy scale $`s`$, which appears in $`L_s`$, transforms into $`|t|s(1\mathrm{cos}\theta )`$, where $`\theta `$ is the electron scattering angle, and, therefore, $`L_sL_t`$. Hence the perturbative expansion contains collinear logs which are modified because of the angular ordering introduced by the radiation cones. In the presence of large scattering angles, for which $`|t|s`$, the above modification is numerically small, but it becomes more and more important in the forward angular range, which is the dynamically favourite region by $`t`$-channel Bhabha scattering and where $`ts`$. The net result is a numerically significant depletion of qed radiation effects just in the most important part of the hard-scattering $`t`$-channel dynamics. Actually, when using sf to evaluate qed LL corrections to small-angle Bhabha scattering, the energy scale $`\mathrm{\Lambda }^2=|t|`$ is employed in all phenomenological applications . More in general, in order to take into account dominant initial-final-state interference effects in addition to initial- and final-state leading terms, $`s`$ and $`t`$ qed contributions to Bhabha scattering can be corrected in terms of a unique combination of Mandelstam invariants given by $`st/u`$, as discussed in refs. . Therefore, the energy scale $`\mathrm{\Lambda }^2=st/u`$ turns out to be a suitable choice for the evaluation in terms of sf of LL corrections to qed Bhabha scattering, as demonstrated, in comparison with the exact $`O(\alpha )`$ calculation, in ref. . Similar arguments for an appropriate choice of the energy scale for QED radiation, based on the inspection of the soft and collinear limit of the $`O(\alpha )`$ correction, have been also advocated in ref. for the reaction $`e^+e^{}W^+W^{}`$ and, very recently, in ref. for the process $`e^+e^{}e^+e^{}\mu ^+\mu ^{}`$. Comparisons performed in refs. with available exact $`O(\alpha )`$ calculations explicitly exhibit the validity of such a strategy, which is therefore pursued in the present analysis. As already remarked, the result for LL corrections in the presence of a calorimetric detection of fs particles must include the contribution of photons which, regardless of their energy, can not be discriminated from closely collinear fermions, as a consequence of the finite angular resolution of the calorimeters. The role of such hard photons collinear to the fs particles becomes therefore unavoidable in the case of a calorimetric measurement of the energy of the fs particles, as discussed in the following.
### 3.1 Soft-photon contribution
In this Section the contribution of photons, too soft to be detected in the calorimeter, will be computed for a generic process with $`n`$ ingoing-legs. <sup>1</sup><sup>1</sup>1 This choice fixes our conventions. Outgoing particles will appear as ingoing ones with momentum and charge according to crossing symmetry. The following approximations are understood
$$\{\begin{array}{cc}q_ik\hfill & \\ s_{ij}m_i^2,m_j^2\hfill & \end{array}$$
(7)
where $`q_i`$ is the momentum of a particle of mass $`m_i`$ emitting a real photon of momentum $`k`$, and $`s_{ij}(q_i+q_j)^2`$ is the invariant mass of the pair $`ij`$.
By following the standard derivation of the eikonal factors due to soft bremsstrahlung and by generalizing it to particles with different masses and charges, the differential cross section, dressed by soft-photon emission, can be cast into the following factorized form
$$d\sigma _{\mathrm{soft}}=d\sigma _0\frac{d\omega }{\omega }\frac{2\alpha }{\pi }\underset{i>j}{\overset{n}{}}Q_iQ_j\mathrm{log}\frac{s_{ij}}{m_im_j}$$
(8)
where $`\omega `$ is the photon energy and $`Q_i`$ is the charge of the $`i`$-th particle.
It is worth noticing that in the limit $`s_{ij}m_i^2,m_j^2`$, provided the first inequality in eq. (7) still holds, the logarithmic behaviour present in eq. (8) disappears, leaving a power law which can be simply obtained by means of the substitution
$$\mathrm{log}\frac{s_{ij}}{m_im_j}\frac{1}{3}\frac{s_{ij}}{m_im_j}$$
(9)
Notice that, since the goal is to determine the scale entering the sf, only the contribution of real photons is explicitly calculated, because the virtual corrections, in order to preserve the cancellations of infrared singularities, must share the same leading collinear structure of the real part itself.
By including the virtual part needed to cancel the infrared singularity and integrating eq. (8) over the photon energy $`\omega `$ in the soft region $`0\omega \mathrm{\Delta }E`$, one gets
$$d\sigma _{\mathrm{S}+\mathrm{V}}=d\sigma _0\mathrm{log}\frac{\mathrm{\Delta }E}{E}\frac{2\alpha }{\pi }\underset{i>j}{\overset{n}{}}Q_iQ_j\mathrm{log}\frac{s_{ij}}{m_im_j}$$
(10)
### 3.2 Hard radiation collinear to the final-state particles
In the case of a calorimetric set-up, which is the realistic situation for single-$`W`$ production at lep, photons collinear to the detected fs particles can not be distinguished from the the emitting particles themselves, independently of the photon energy. Therefore, in order to obtain the correct structure of double-log corrections for such an event selection, the effect due to the emission of unresolved hard radiation collinear to the fs particles must be taken into account in addition to soft+virtual corrections.
To this end, let us re-consider the previous process with $`n`$ ingoing legs and think $`m`$ of them to be changed into outgoing legs at the end of the calculation (see the previous footnote). Then, the contribution of photons collinear to the fs particles can be cast into a gauge invariant form as follows
$$d\sigma _{\mathrm{hard}}=d\sigma _0\frac{d\omega }{\omega }\frac{2\alpha }{\pi }\underset{i}{\overset{m}{}}Q_i^2\mathrm{log}\frac{E_i\delta }{m_i}$$
(11)
where $`E_i`$ is the energy of the $`i`$-th particle, and $`\delta `$ is the half-opening angle of the calorimetric resolution.
By integrating eq. (11) over the photon energy $`\omega `$ in the range $`\mathrm{\Delta }E\omega E`$, the integrated correction due to hard photons collinear to the fs particles is given by
$$d\sigma _{\mathrm{hard}}=d\sigma _0\mathrm{log}\frac{E}{\mathrm{\Delta }E}\frac{2\alpha }{\pi }\underset{i}{\overset{m}{}}Q_i^2\mathrm{log}\frac{E_i\delta }{m_i}$$
(12)
### 3.3 The master formula
Equations (10) and (12) give the leading double-log contribution which must be compared to (6), the $`O(\alpha )`$ perturbative expansion of eq. (5), in order to fix the process scales $`\mathrm{\Lambda }_i`$. Summing the contributions of eq. (10) and eq. (12), the analytical cross section is in conclusion given by
$`d\sigma _{\mathrm{S}+\mathrm{V}}+d\sigma _{\mathrm{hard}}`$ $`=`$ $`d\sigma _0{\displaystyle \frac{2\alpha }{\pi }}\mathrm{log}{\displaystyle \frac{\mathrm{\Delta }E}{E}}\{{\displaystyle \underset{i=m+1}{\overset{n}{}}}Q_i^2\mathrm{log}{\displaystyle \frac{E_i}{m_i}}+`$
$`{\displaystyle \underset{i>j}{\overset{n}{}}}Q_iQ_j\mathrm{log}2(1c_{ij}){\displaystyle \underset{i}{\overset{m}{}}}Q_i^2\mathrm{log}\delta \}`$
where $`c_{ij}`$ is the cosine of the angle between particles $`i`$ and $`j`$.
Three different kinds of logarithms occur in eq. (3.3). The first term contains the mass and energy logarithms of the is particles only, since, as expected, the energies and the masses of the fs particles disappeared, in agreement with the kln theorem . The second term includes angular terms due to radiation interference, while the third one comes from the requirement of calorimetric measurement.
These terms must be compared with the collinear logarithms of eq. (6) in order to fix the scales $`\mathrm{\Lambda }_i`$ of the sf. In the following Section this task is accomplished in detail for the single-$`W`$ process.
## 4 Fixing the radiation scales in the single-$`W`$ process
Let us consider, for definiteness, the process $`e^+e^{}e^{}\overline{\nu }u\overline{d}`$ with the fs electron lost in the beam pipe (single-$`W`$ process). In this event selection (hereafter es) the leading contribution comes from $`\gamma ^{}e^+`$ scattering with the virtual photon emitted from the electron line. The leading dynamics is given by the $`t`$-channel Feynman diagrams shown in Fig. 1.
If a calorimetric measurement of the energies of the fs particles is performed, only the is legs need to be corrected by the sf. Furthermore, since the electron is scattered in the very forward region, the interference between the electron line and the rest of the process is very small. This allows a natural sharing of the logarithms coming from eq. (3.3) between the two is sf associated to the colliding electron and positron.
Hence the formula (3.3), when compared with eq. (6), translates into the two following scales ($`\mathrm{\Lambda }_{}`$ refers to sf attached to the electron line, while $`\mathrm{\Lambda }_+`$ to the sf attached to the positron line),
$$\mathrm{\Lambda }_{}^2=4E^2\frac{(1c_{})^2}{\delta ^2},\mathrm{\Lambda }_+^2=2^{\frac{14}{9}}E^2\frac{\left((1c_{\overline{d}})(1c_u)^2\right)^{\frac{2}{3}}}{\left((1c_{u\overline{d}})^2\delta ^5\right)^{\frac{2}{9}}}$$
(14)
where $`E`$ is the beam energy, $`c_{}`$ is the cosine of the electron scattering angle, $`c_u`$ and $`c_{\overline{d}}`$ are the cosine of the quark scattering angles with respect to the initial positron, and $`c_{u\overline{d}}`$ is the cosine of the relative angle between the quarks.
It is worth noticing that in the numerical implementation, whenever one of the two scales is less than a small cut-off ($`\mathrm{\Lambda }_{\mathrm{cut}\mathrm{off}}^2=4m_e^2`$, where $`m_e`$ is the electron mass), the radiation from the corresponding leg is switched off, in accordance with eq. (9). It was carefully tested that variations of the cut-off do not alter the numerical results.
Owing to the presence of a resonant $`W`$ boson, some modifications to the previous results may come from finite-width effects and from radiation decoherence . Finite-width corrections of the form of $`E_\gamma /\mathrm{\Gamma }_W`$ arise when the unstable particle propagator goes off its mass-shell, but this is not the present case, since the multi-fermion final state can accommodate a resonant $`W`$. Radiation decoherence is present whenever a resonance occurs and its effect is to cancel the angular dependence from the scale. As a consequence the scale $`\mathrm{\Lambda }_+`$ should be modified by dropping the angular interference factors in eq. (14) when the emitted photons have $`E_\gamma \mathrm{\Gamma }_W`$. Yet in the present case the effect is tiny, since the effects due to angular interference for the scale $`\mathrm{\Lambda }_+`$ are already small by themselves.
It is also possible to make a naive ansatz for the radiation scales without a detailed calculation, by thinking of the graphs of Fig. 1 in terms of the Weizsรคcker-Williams approximation , i.e. in terms of a convolution of the process $`e^+\gamma \nu _eW^{}`$ with an equivalent photon spectrum plus a real electron line. This leads to assigning two different scales to the single-$`W`$ process: one scale for the electron current and one for the positron current. The former scale is the proper one for a $`t`$-channel process, e.g. $`t`$-channel Bhabha scattering, so it is simply $`|q_\gamma ^{}^2|`$, where $`|q_\gamma ^{}^2|`$ is the squared momentum transfer in the $`ee\gamma ^{}`$ vertex. The latter is the sum of an $`s`$-channel electron exchange and a $`t`$-channel $`W`$ exchange (see Fig. 1). Assuming that the $`t`$-channel dominates, its natural cut-off is given by the $`W`$-boson mass, $`M_W`$. Hence, the following ansatz follows
$$\mathrm{\Lambda }_{,\mathrm{naive}}^2=|q_\gamma ^{}^2|,\mathrm{\Lambda }_{+,\mathrm{naive}}^2=M_W^2$$
(15)
where $`M_W`$ is the mass of the $`W`$ boson. The comparison between the scales given by eq. (14) and these naive scales, which will be performed numerically in the following Section, provides a useful cross-check of the analytical results derived by inspection with the soft/collinear limit of the $`O(\alpha )`$ correction.
A discussion of other possible approaches to the treatment of photonic corrections to single-$`W`$ production can be found in the four-fermion working group report of the lep2 mc Workshop .
## 5 The running of the electromagnetic coupling constant
Besides the higher-order qed corrections discussed in the previous Sections, other large logarithmic contributions to the single-$`W`$ cross section arise from the running of the electromagnetic coupling constant $`\alpha `$. Since in the case under study the dominant configurations come from the Feynman diagrams with an almost on-shell photon exchange, the appropriate scale for the evaluation of the electromagnetic coupling relative to the $`t`$-channel photon in the $`ee\gamma ^{}`$ vertex is the squared momentum transfer $`q_\gamma ^{}^2`$ defined above.
However, because $`G_F`$, $`M_W`$ and $`M_Z`$ are the typically adopted input parameters for electroweak processes at lep2, the electromagnetic coupling is fixed at tree-level to a high energy value as, for example,
$$\alpha _{G_F}=4\sqrt{2}\frac{G_FM_W^2s_W^2}{4\pi },\mathrm{with}s_W^2=1\frac{M_W^2}{M_Z^2}.$$
(16)
On the other hand, the single-$`W`$ production is a $`q_\gamma ^{}^20`$ dominated process and therefore the above high-energy evaluation of $`\alpha `$, $`\alpha _{G_F}`$, needs to be rescaled to its correct value at small momentum transfer. To this end, a gauge-invariant โreweightingโ procedure can be adopted, by rescaling the differential cross section $`d\sigma /dt`$ ($`tq_\gamma ^{}^2`$) in the following way
$$\frac{d\sigma }{dt}\frac{\alpha ^2(t)}{\alpha _{G_F}^2}\frac{d\sigma }{dt},$$
(17)
where $`\alpha (t)`$ is the running coupling constant computed at virtuality $`q_\gamma ^{}^2`$.
A detailed analysis of the effect of the running couplings in single-$`W`$ production has been recently performed within the massive fermion-loop scheme in ref. , where the couplings are automatically running in the calculation. As shown in ref. , the relative difference between the above reweighting prescription and the complete results of the fermion-loop scheme is at the 1%-2% level <sup>2</sup><sup>2</sup>2 Actually, for the single-$`W`$ final state under examination here and for realistic event selections, the differences between the two procedures are confined below the 1% level., and it is therefore in the expected range of theoretical uncertainty due to missing full one-loop electroweak corrections.
## 6 Numerical results
In this Section the mc code, developed to simulate the single-$`W`$ process, is described and a sample of numerical results obtained by means of it is shown and commented, with particular emphasis on the effects of higher-order qed corrections to single-$`W`$ production at lep2 energies.
### 6.1 The Monte Carlo code
A mc program, named SWAP, was developed to calculate cross sections and differential distributions for the single-$`W`$ signature.
As already emphasized, the main feature of this process is the fact that the $`t`$-channel photon of Fig. 1 becomes quasi-real. In the limit of massless fermions, the photon propagator becomes singular in the forward direction and the cross section develops a logarithmic singularity. Indeed, whenever the final electron is lost in the beam pipe, its mass becomes a natural cutoff for the very-forward singularities, compelling to build a massive matrix element and phase-space. The phase-space integration is performed in SWAP with the aid of a multi-channel importance sampling with stratification. The main peaking structures for the single-$`W`$ process are given by the dynamics depicted by the fusion and bremsstrahlung graphs of Fig. 1. They are the resonant $`W`$-boson invariant mass, treated with a Breit-Wigner weight, and the $`t`$-channel โsingularityโ of the quasi-real photon, treated with a $`1/|t|`$ weight. Moreover, the program can deal with the singularities of the sub-leading $`t`$-channel CC$`20`$ diagrams shown in Fig. 2, by means of the multi-channel approach.
The exact hard-scattering matrix element is computed by means of the ALPHA code for the automatic evaluation of Born scattering amplitudes. Fermion masses are exactly taken into account and the fixed-width scheme is adopted as gauge-restoring approach, by taking the massive gauge boson propagator as follows:
$$\mathrm{\Pi }^{\mu \nu }=\frac{i\left(g^{\mu \nu }\frac{k^\mu k^\nu }{M^2i\mathrm{\Gamma }M}\right)}{k^2M^2+i\mathrm{\Gamma }M},\mathrm{\Gamma }=\mathrm{cost}.$$
(18)
It is known that this scheme preserves $`U(1)`$ gauge invariance but still violates $`SU(2)`$ Ward identities. However, at least in the unitary gauge employed here, it is indistinguishable from other fully gauge-invariant schemes .
The contribution of anomalous gauge couplings is also accounted for in SWAP. The anomalous gauge boson couplings $`\mathrm{\Delta }k_\gamma `$, $`\lambda _\gamma `$, $`\delta _Z`$, $`\mathrm{\Delta }k_Z`$ and $`\lambda _\gamma `$ are implemented in the ALPHA code according to the parameterization of refs. . Photon radiation is implemented via sf formalism, according to the discussion of Sect. 4. It is worth noticing that, since the incoming electron/positron are required to be on-shell massive fermions, a naive four-momentum rescaling, due to photon emission, such as $`\widehat{p}_\pm =xp_\pm `$ leads to potentially dangerous gauge violations, according to what previously discussed. Therefore, the rescaled incoming four-momenta are implemented as $`\widehat{p}_\pm =(xE,0,0,\pm \sqrt{x^2E^2m_e^2})`$, by interpreting $`x`$ as the energy fraction after photon radiation, as motivated in ref. . If required, $`p_{}/p_L`$ effects can be provided in the treatment of ISR, by means of either $`p_{}`$-dependent SF or a QED Parton Shower algorithm . The effect of vacuum polarization is taken into account as described by eq. (17), by including the contribution of leptons, heavy quarks and light quarks, the latter according to the parameterization of ref. . The program supports realistic es and it can be employed either as a cross-section calculator or as a event generator, with both weighted and unweighted generation available.
The technical precision of the event generator SWAP has already been carefully proved in ref. , by means of detailed tuned comparisons between the predictions of independent codes. Perfect agreement was found, both at the level of integrated cross sections and distributions, also for purely leptonic final states.
### 6.2 Discussion of the numerical results
The numerical simulations are elaborated according to the es reviewed in Tab. 1, with the electroweak input parameters shown in Tab. 2.
In Figs. 3-4 the numerical impact of different choices of the $`\mathrm{\Lambda }^2`$-scale on the cross section of the single-$`W`$ process $`e^+e^{}e^{}\overline{\nu }u\overline{d}`$ in the lep2 energy range is shown. Since the energy scale $`\mathrm{\Lambda }_+`$ of eq. (14) depends on the quark scattering angles, two different quark angular acceptances are considered, namely no cut (Fig. 3) and $`|\mathrm{cos}\vartheta _{u,\overline{d}}|<0.95`$ (Fig. 4). The marker $``$ represents the Born cross section, $``$ represents the correction given by $`\mathrm{\Lambda }_\pm ^2=s`$ scale for both IS sf, $`\mathrm{}`$ represents the correction given by $`\mathrm{\Lambda }_\pm ^2=|q_\gamma ^{}^2|`$ scale for both IS sf, $`\mathrm{}`$ the correction given by the scales of eq. (14), $`\mathrm{}`$ the correction given by the naive scales of eq. (15). It can be seen that neither the $`s`$ scale, as implemented in some computational tools used for the analysis of the single-$`W`$ process, nor the $`|q_\gamma ^{}^2|`$ scale are able to reproduce the effects due to the scales of eq. (14) and eq. (15). These two scales are in good agreement and both predict a lowering of the Born cross section of about 8-9%, almost independent of the c.m. lep2 energy and quark angular acceptance. This fact can be understood by looking at Fig. 5, where it is shown the single-$`W`$ differential cross section with respect to the scales $`\mathrm{\Lambda }_\pm `$ of eq. (14). On the left, $`\mathrm{\Lambda }_+`$ exhibits a broad peak not far from $`M_W`$, while, on the right, the other scale $`\mathrm{\Lambda }_{}`$ peaks, as expected, at very small momentum transfer.
Figure 6 shows the effects of the reweighting procedure of eq. (17) for the evaluation of the qed running coupling constant. The marker $`\mathrm{}`$ represents the relative difference between the integrated cross section computed in terms of $`\alpha _{G_F}`$ and the cross section computed in terms of $`\alpha (0)`$, while the marker $`\mathrm{}`$ is the relative difference between the integrated cross section computed in terms of $`\alpha _{G_F}`$ and the cross section computed in terms of $`\alpha (t)`$. As can be seen, the rescaling from $`\alpha _{G_F}`$ to $`\alpha (t)`$ introduces a negative correction of about 5-6% in the lep2 energy range. The difference between $`\mathrm{}`$ and $`\mathrm{}`$, which is about 2-3%, is a measure of the running of $`\alpha _{QED}`$ from $`q_\gamma ^{}^2=0`$ to $`q_\gamma ^{}^2=t`$.
As an illustrative example of the effect of anomalous couplings on single-$`W`$ differential distributions, in Fig.7 the distribution of the $`q\overline{q}`$ invariant mass, around the peak of the $`W`$-boson resonance, and the distribution of the angle of the quarks with the line of flight of the reconstructed $`W`$-boson in the $`W`$-boson rest-frame are shown. The dashed lines correspond to the simulation as obtained by means of SWAP for the anomalous coupling $`\mathrm{\Delta }\kappa _\gamma =0.1`$, while the solid lines represent the Standard Model prediction. The effect of the anomalous coupling $`\mathrm{\Delta }\kappa _\gamma `$ at lep2 energies is just an overall rescaling of the total cross section. Therefore the lep2 sensitivity to $`\mathrm{\Delta }\kappa _\gamma `$ in single-$`W`$ events depends crucially on the accuracy of the theoretical evaluation of the total cross section.
## 7 Conclusions
The process of single-$`W`$ production in high-energy $`e^+e^{}`$ collisions is relevant at lep2 for the determination of the non-abelian self-couplings of the $`W`$ boson, and of primary importance at future Linear Colliders at the TeV scale, its cross section being dominant at very high energies with respect to other four-fermion processes.
In order to give a contribution to the reduction of the theoretical uncertainty presently associated to the calculation of the single-$`W`$ cross section, the issue of higher-order photonic corrections has been carefully investigated within the standard sf technique. Theoretical and phenomenological arguments for the choice of the energy scale entering the sf have been proposed. Two possible solutions for the scale of qed radiation have been obtained. The former has been derived by means of general arguments concerning the soft and collinear limit of the $`O(\alpha )`$ corrections coming from the radiation of external legs. The latter, which can be considered as a naive ansatz, has been driven by thinking of the single-$`W`$ process in terms of the Weizsรคcker-Williams approximation.
Numerical calculations show that the typically adopted choice of the center-of-mass energy of the reaction, as radiation scale for the process, can lead to over-estimate the radiative correction by a factor of $`1.5`$, implying an under-estimate of the cross section of about $`4\%`$. Also the choice of fixing the scale to the momentum transfer $`t`$ in the $`ee\gamma ^{}`$ vertex for both the is sf leads to an under-estimate of the photon correction of about $`4\%`$. The difference between the predictions given by the two set of scales of eq. (14) and eq. (15) is at the per mille level in the lep2 energy range. Therefore, the naive scales of eq. (15) provide a good ansatz for the energy scale of QED radiation in the single-$`W`$ process, which could be simply implemented in MC tools for data analysis and further corroborated by the comparison with the results of other approaches. The method here described for the energy scale determination in the sf can be simply generalized to other four-fermion process dominated by non-annihilation channels, such as single-$`Z`$ production.
In order to provide adequate phenomenological predictions for precision experiments, also the running of the electromagnetic coupling constant has been accounted for in an effective way, i.e. by rescaling the differential cross section for the ratio of the electromagnetic coupling constant, evaluated at the typical scale of the process, to the same coupling evaluated from the input parameters according to tree-level relations. The effect of such rescaling amounts to a negative correction of about $`5`$-$`6\%`$, in agreement with recent findings , as far as the effect of $`\alpha _{\mathrm{QED}}`$ is concerned.
In the light of the experimental precision for the single-$`W`$ process, the corrections considered in the present paper are phenomenologically relevant.
According to the theoretical approach described in the present paper, an original mc programme SWAP has been developed, including exact tree-level matrix elements with finite fermion masses effects, anomalous couplings, vacuum polarization and higher-order qed corrections. The code is available for experimental analysis.
The authors wish to thank the members of the four-fermion working group of the lep2 Monte Carlo Workshop (cern), in particular Y. Kurihara, G. Passarino, R. Pittau and M. Verzocchi, for useful discussions on the subject. A. Pallavicini is grateful to the infn, Sezione di Pavia, for having provided computer facilities.
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# Thick Domain Walls Intersecting a Black Hole
## I Introduction
Topological defects arise during spontaneous symmetry breaking associated with phase transitions, and cosmological evolution of them is considered to have played an important role in cosmology (see e.g., ). Topological defects produced in the early universe might give us some information on high energy phenomena which cannot be reached by accelerator experiments, and those produced at the late time phase transitions also have attracted attention as a potential source of the cosmic structures. To study the topological defects is therefore crucially important in cosmology and elementary particle physics.
In general relativity, topological defects have interesting features. Topological defects are extended and relativistic objects due to their large tension. In particular, the gravitational field produced by an infinitely thin domain wall shows repulsive nature . On the other hand, although there are many studies about the properties of thick domain walls in the flat and de Sitter spacetime backgrounds , little is known about the existence of thick wall configurations on inhomogeneous, strongly curved background such as a black hole spacetime. Thus it is intriguing to study the gravitational interaction between two extended relativistic objects: a topological defect and a black hole.
In most studies of defectโblack hole system, topological defects have been approximately treated as infinitely thin and non-gravitating objects whose dynamics is governed by Nambu-Goto action. Within this context, the scattering problem of a Nambu-Goto string by a background black hole has been studied in detail . In the domain wall case, recently M. Christensen, V. P. Frolov and A. L. Larsen considered Nambu-Goto walls embedded in the Schwarzschild black hole spacetime and found the static axisymmetric solutions. They showed that there exist a family of infinitely thin walls which intersect the black hole event horizon. However, little attention has been given to the system of a thick defect interacting with a black hole though a defect as a topologically stable configuration of a scalar field has a finite thickness.
Now we shall consider the validity of thin-wall approximation in the system of a topological defect and a black hole. In such a system, there are two characteristic scales: the thickness $`w`$ of the defect and the black hole radius $`R_g`$. In the case of the system of an astrophysical black hole with the mass $`M_{}`$ and a defect formed during a GUT phase transition, the thickness of the defect is much smaller than the black hole radius and therefore thin-wall approximation would be valid. However, it is not so hard to consider the situation that the thickness of defects cannot be negligible as compared with the size of a small black hole. Over the last few decades, many people have studied the formation of small black holes called primordial black holes (PBHs) and their cosmological implications. For example, studying the contribution of PBHs to cosmic rays enables one to place limits on the spectrum of density fluctuations in the early universe (see e.g. ). On the other hand, it has been discussed the possibility of thick defects and their roles in cosmology, e.g. as a source of large-scale structure in the universe , or as a candidate for some kind of dark matter . For example, it is thought that the typical mass of PBHs which evaporate at the present epoch is $`10^{15}\text{g}`$, so $`R_g10^{13}\text{cm}`$. When one considers the topological defects formed during a phase transition at $`100\text{MeV}`$, such defects become thicker than the size of PBHs and thin-wall approximation is no longer valid.
In this paper, we investigate the gravitational interaction between a domain wall and a black hole taking the thickness of the wall into account. We deal with scalar fields in the Schwarzschild black hole spacetime with $`\varphi ^4`$ and sine-Gordon potentials, which have a discrete set of degenerate minima. We explicitly show that static axi-symmetric thick domain walls intersecting the black hole do exist by numerical investigation. We consider a non-gravitating domain wall for simplicity. This test wall assumption might be valid when the symmetry braking scale of the scalar field is much lower than Planck scale as will be shown later by dimensional analysis.
This paper is organized as follows. In section 2, we derive the basic equation and discuss the boundary conditions which represent the situation we want to study. In section 3, we show the numerical result. We summarize our work in section 4. We also discuss the validity of the assumption that the effects of gravity of domain wall can be ignored near the black hole horizon. Throughout this paper, we use units such that $`c=\mathrm{}=G=1`$ unless otherwise stated.
## II The basic equation and the boundary conditions
We consider a static thick domain wall in a black hole spacetime. The domain wall is constructed by a scalar field with self-interaction in a given curved spacetime. In what follows, we neglect the self gravity of the scalar field, as will be justified later. As a background spacetime, we consider the Schwarzschild black hole
$$g=\left(1\frac{2M}{R}\right)\mathrm{d}t^2+\left(1\frac{2M}{R}\right)^1\mathrm{d}R^2+R^2(\mathrm{d}\vartheta ^2+\mathrm{sin}^2\vartheta \mathrm{d}\phi ^2).$$
(1)
For our purpose, we find that it is more convenient to work in the isotropic coordinates $`\{t,r,\vartheta ,\phi \}`$, where the new radial coordinate $`r`$ is defined by
$$R=r\left(1+\frac{M}{2r}\right)^2.$$
(2)
We mainly consider the region outside the event horizon in this paper, which corresponds to $`r>M/2`$. In this coordinate system, the metric has a spatially conformally flat form
$$g=\left(\frac{2rM}{2r+M}\right)^2\mathrm{d}t^2+\left(1+\frac{M}{2r}\right)^4\left[\mathrm{d}r^2+r^2(\mathrm{d}\vartheta ^2+\mathrm{sin}^2\vartheta \mathrm{d}\phi ^2)\right].$$
(3)
Let us consider a real scalar field $`\varphi `$ with a potential $`V[\varphi ]`$, of which Lagrangian is given by
$$=(detg)^{1/2}\left(\frac{1}{2}\varphi \varphi +V[\varphi ]\right).$$
(4)
The equation of motion for $`\varphi `$ is
$$^2\varphi V/\varphi =0.$$
(5)
In this paper, we consider following two familiar types of potentials which have a discrete set of degenerate minima; the $`\varphi ^4`$ potential
$$V_1[\varphi ]=\frac{\lambda }{4}(\varphi ^2\eta ^2)^2,$$
(6)
and the sine-Gordon potential
$$V_2[\varphi ]=\lambda \eta ^4[1+\mathrm{cos}(\varphi /\eta )].$$
(7)
Since the Schwarzschild spacetime is asymptotically flat, the asymptotic boundary condition for Eq. (5) would be given by the solution in the flat spacetime. The relevant solutions in the flat spacetime $`g=\mathrm{d}t^2+\mathrm{d}x^2+\mathrm{d}y^2+\mathrm{d}z^2`$ are the static and plane-symmetric solutions
$$\varphi _1(z)=\eta \mathrm{tanh}\sqrt{\lambda /2}\eta z$$
(8)
and
$$\varphi _2(z)=\eta [4\mathrm{arctan}\mathrm{exp}(\sqrt{\lambda }\eta z)\pi ],$$
(9)
for the potentials $`V_1`$ and $`V_2`$, respectively. These solutions represent domain walls in the flat spacetime characterized by the thickness of the wall
$$w=1/\sqrt{\lambda }\eta .$$
(10)
In the Schwarzschild background, the solution compatible with the above asymptotic boundary condition would have a static and axi-symmetric form $`\varphi =\varphi (r,\vartheta )`$. Then, the explicit form of the equation of motion (5) becomes
$$\left(\frac{2r}{2r+M}\right)^4\left[\frac{^2}{r^2}+\frac{8r}{(4r^2M^2)}\frac{}{r}+\frac{1}{r^2}\left(\frac{^2}{\vartheta ^2}+\mathrm{cot}\vartheta \frac{}{\vartheta }\right)\right]\varphi =\frac{V}{\varphi }.$$
(11)
This equation can be parameterized by a single dimensionless parameter
$$ฯต=M/2w,$$
(12)
by introducing dimensionless variables
$$\rho =2r/M,\mathrm{\Phi }(\rho ,\vartheta )=\varphi (r,\vartheta )/\eta .$$
(13)
In terms of these variables, Eq. (11) becomes
$$\left(\frac{\rho }{\rho +1}\right)^4\left[\frac{^2}{\rho ^2}+\frac{2\rho }{(\rho ^21)}\frac{}{\rho }+\frac{1}{\rho ^2}\left(\frac{^2}{\vartheta ^2}+\mathrm{cot}\vartheta \frac{}{\vartheta }\right)\right]\mathrm{\Phi }=ฯต^2\frac{U}{\mathrm{\Phi }},$$
(14)
where the dimensionless potential $`U[\mathrm{\Phi }]=V[\varphi ]/\lambda \eta ^4`$ is defined. $`U`$ has minima at $`\mathrm{\Phi }=\pm 1`$ for $`\varphi ^4`$ potential and at $`\mathrm{\Phi }=(2n+1)\pi `$ ($`n=0,\pm 1,\pm 2,\mathrm{}`$) for sine-Gordon potential. The parameter $`ฯต`$ is just a ratio of the horizon radius to the wall thickness measured in the asymptotic region, namely if $`ฯต`$ is smaller (larger) than unity, then the wall is said to be thick (thin) as compared to the size of the black hole.
We shall confine ourselves to the case that the core of wall is located at the equatorial plane $`\{\vartheta =\pi /2\}`$ of the black hole. The solutions without this assumption will be discussed in the separated paper . Accordingly, we impose the Dirichlet boundary condition at the equatorial plane
$$\mathrm{\Phi }|_{\vartheta =\pi /2}=0.$$
(15)
Now it is sufficient to consider the north region, namely the solution in the south region can be obtained via $`\mathrm{\Phi }(\rho ,\vartheta )=\mathrm{\Phi }(\rho ,\pi \vartheta )`$, $`\{\pi /2\vartheta \pi \}`$ of the spacetime. The regularity of the scalar field at the symmetric axis is give by the Neumann boundary condition
$$\frac{\mathrm{\Phi }}{\vartheta }|_{\vartheta =0}=0.$$
(16)
On the other hand, the boundary condition at the event horizon $`\{\rho =1\}`$ is given by the Neumann boundary condition
$$\frac{\mathrm{\Phi }}{\rho }|_{\rho =1}=0.$$
(17)
As is shown in Appendix, the condition (17) is the consequence of a natural requirement that the energy density observed by a freely falling observer remains finite at the event horizon. In practice, the region of numerical integration should be finite, so that we need an asymptotic boundary condition at $`\rho =\rho _{\mathrm{max}}`$ for $`\rho _{\mathrm{max}}1`$. Taking into account the flat background solutions (8) and (9), we impose the Dirichlet boundary condition
$$\mathrm{\Phi }_1|_{\rho =\rho _{\mathrm{max}}}=\mathrm{tanh}(2^{1/2}ฯต\rho _{\mathrm{max}}\mathrm{cos}\vartheta )$$
(18)
and
$$\mathrm{\Phi }_2|_{\rho =\rho _{\mathrm{max}}}=4\mathrm{arctan}\mathrm{exp}(ฯต\rho _{\mathrm{max}}\mathrm{cos}\vartheta )\pi ,$$
(19)
for the $`\varphi ^4`$ and the sine-Gordon potentials, respectively.
In the next section, we numerically integrate the field equation (14) using relaxation method under these boundary conditions at the equatorial plane (15), at the symmetric axis (16), at the event horizon (17) and in the asymptotic region (18), (19) for both the $`\varphi ^4`$ and the sine-Gordon potentials.
## III Numerical integration
The scalar field configurations $`\mathrm{\Phi }(x,z)`$ satisfying (14) and the boundary conditions are shown in Figs. 14, where $`x`$ and $`z`$ are the Cartesian coordinates $`x=\rho \mathrm{sin}\vartheta `$, $`z=\rho \mathrm{cos}\vartheta `$. Here we show the results in typical two cases; the $`ฯต=1`$ case in which the thickness of the kinks (18), (19) at $`\rho _{\mathrm{max}}`$ is comparable to the Schwarzschild radius (Fig. 1 for $`\varphi ^4`$ and Fig. 3 for sine-Gordon), and the $`ฯต=0.1`$ case in which the thickness of the kinks (18), (19) at $`\rho _{\mathrm{max}}`$ is one order of magnitude larger than the Schwarzschild radius (Fig. 2 for $`\varphi ^4`$ and Fig. 4 for sine-Gordon). In both the cases, we obtain a domain wall solution as a kink of the scalar field at the equatorial plane $`z=0`$. Particularly in the case $`ฯต=0.1`$, the black hole is enveloped in the core region of the wall.
We also show the energy density $`E`$ of the scalar field given by
$$E\frac{T_t^t}{\lambda \eta ^4}=\frac{1}{2ฯต^2}\left(\frac{\rho }{\rho +1}\right)^4\left[\left(\frac{\mathrm{\Phi }}{\rho }\right)^2+\frac{1}{\rho ^2}\left(\frac{\mathrm{\Phi }}{\vartheta }\right)^2\right]+U[\mathrm{\Phi }]$$
(20)
in Figs. 58, corresponding to Figs. 14, respectively. In all the cases, one can see that the configuration actually has a wall-like structure, namely the energy density is localized around the equatorial plane with a certain thickness corresponding to $`ฯต`$.
We shall comment on the computational domain and the grid spacing taken in our calculation. In order that the asymptotic boundary conditions (18) and (19) make sense, $`\rho _{\mathrm{max}}`$ must be large enough. In the above calculation, we take the computational domain which is fifty times as large as the horizon radius (i.e. $`\rho _{\mathrm{max}}=51`$), and we carry out the integration on $`500\times 90`$ grid (the grid spacing in $`\rho `$\- and $`\vartheta `$\- direction is $`0.1\times `$horizon-radius and $`1^{}`$, respectively). Then we clip the region $`\{|x|,|z|<20\}`$, where the above results are insensitive to the value of $`\rho _{\mathrm{max}}`$ and the number of grid points. In fact, the results do not change when we extend the computational domain to $`\rho _{\mathrm{max}}=101`$ with keeping the grid spacing, and the results differ from ones on the finer ($`1000\times 180`$) grid with keeping $`\rho _{\mathrm{max}}=51`$ at most 1%. We also comment that the reliability of the numerical code is checked in the flat space case and it reprodeces the exact solutions (8), (9) with accuracy of $`10^2`$.
## IV Summary and Discussions
In order to answer the question whether or not scalar fields can actually form a topological defect in the vicinity of a black hole, we have numerically solved the equation of motion for real scalar fields with $`\varphi ^4`$ and sine-Gordon potentials, which have a discrete set of degenerate minima, in the Schwarzschild black hole background. In both $`\varphi ^4`$ and sine-Gordon potential cases, we showed that there exist the static axi-symmetric field configurations which represent thick domain walls intersecting the black hole. In particular, we studied the specific case that the wallโs core is located at the equatorial plane of the Schwarzschild spacetime. We introduced the parameter $`ฯต`$ which characterizes the domain wall thickness compared to the black hole horizon radius; the smaller than unity $`ฯต`$ is, the larger the wall width is. We showed the domain wall solutions and their energy densities for $`ฯต=1`$ and $`ฯต=0.1`$ cases. In summary, we can say that a black hole is not an obstacle for scalar fields to form a domain wall configuration intersecting the black hole.
One might wonder about our present results; one naively expects that the scalar field could have no static distribution around a black hole and inevitably fall into the horizon as usual objects do. However, a domain wall is a relativistic object with a large negative pressure (or large tension) whose magnitude is comparable to that of energy density. Furthermore in our study we examined the domain walls which are extended infinitely in space. Then we can understand that the domain wall is suspended from the asymptotic region and supported against falling into the black hole by its tension, so that the static configuration is realized.
In our analysis, we assumed that the gravitational effect of the domain wall is negligible compared to that of the Schwarzschild black hole. We shall comment on the validity of this assumption. The energy density of a domain wall is given by
$$GT_t{}_{}{}^{t}G\lambda \eta ^4\frac{1}{w^2}\left(\frac{\eta }{m_{\mathrm{Pl}}}\right)^2,$$
(21)
where $`m_{\mathrm{Pl}}`$ is the Planck mass and $`G`$ is the Newtonian constant. On the other hand, the curvature strength, or gravitational tidal force, of Schwarzschild spacetime is estimated as
$$\frac{GM}{R^3}ฯต\frac{w}{R^3}$$
(22)
at areal radius $`R`$. Then a ratio of the gravity which will be produced by the domain wall to the gravity of the background black hole is given by
$$\omega \frac{GT_t^t}{GM/R^3}\frac{1}{ฯต}\left(\frac{R}{w}\right)^3\left(\frac{\eta }{m_{\mathrm{Pl}}}\right)^2.$$
(23)
When $`\omega `$ becomes much smaller than unity, the gravity of the domain wall is negligible compared to that of the black hole, and our test wall assumption becomes valid. From Eq. (23) we have $`\omega ฯต^2(\eta /m_{\mathrm{Pl}})^2`$ near the horizon ($`R2GM`$). Therefore, when we consider domain walls with the symmetry breaking scale being much lower than the Planck scale (i.e. $`\eta m_{\mathrm{Pl}}`$), we have $`\omega 1`$ and consequently our result gives a good description of shapes of gravitating thick domain walls near the black hole. We can also see from Eq.(23) that, in the thick wall ($`w>2GM`$) case, if $`\eta m_{\mathrm{Pl}}`$, the test wall assumption is still valid even at $`Rw`$.
In the asymptotic region ($`Rw`$), one may expect that the gravity of the domain wall is no longer negligible and changes the asymptotic geometry drastically. Bonjour, Charmousis and Gregory have recently investigated the spacetime of a thick gravitating domain wall with local planar symmetry and reflection symmetry around the wallโs core . They showed that the domain wall spacetime becomes spatially compact and has a cosmological horizon as de Sitter spacetime does. This suggests that, when a black hole exists and wallโs gravity is taken into account in the region far from the black hole, the whole spacetime has a cosmological horizon and an axi-symmetric domain wall intersects both the black hole and the cosmological horizons as an equatorial plane in a Schwarzschild-de Sitter spacetime. This motivates us to study further the interaction between thick domain walls and black holes in, for example, Schwarzschild-de Sitter background.
The domain wall solutions obtained here are thought to represent a possible final configuration of a gravitational capturing of a domain wall by a black hole. At present, it is far reaching for us to investigate a fully dynamical process such as the scattering and capture of thick domain walls by black holes. However, to get some insights into the problem, it is worth investigating the existence of the static axi-symmetric solutions which represent thick domain walls located away from a black hole.
A cosmic string is also an extended relativistic object with large tension and thought to play more important role in cosmology than a domain wall does. Study of a gravitationally interacting system of thick cosmic strings and black holes is an interesting problem as a generalization of the present analysis.
## ACKNOWLEDGMENTS
We would like to thank Profs. Hideki Ishihara, Hideo Kodama and Takashi Nakamura for many useful suggestions and comments. We also thank Susumu Higaki for stimulating discussions. D.I. was supported by JSPS Research Fellowships for Young Scientists, and this work was supported in part by the Grant-in-Aid for Scientific Research Fund (No. 4318). A.I. was supported by JSPS Research Fellowships for Young Scientists and this work was partially supported by Handai Yukawa Shogakukai (A.I.).
## A The Boundary Condition at the Event Horizon
The tangent of a freely falling observer parameterized by its proper time $`\tau `$ is $`u^\mu =(\mathrm{d}t/\mathrm{d}\tau ,\mathrm{d}r/\mathrm{d}\tau ,0,0)`$, and we have
$$1=g_{\mu \nu }u^\mu u^\nu =\left(\frac{2rM}{2r+M}\right)^2\left(\frac{\mathrm{d}t}{\mathrm{d}\tau }\right)^2+\left(1+\frac{M}{2r}\right)^4\left(\frac{\mathrm{d}r}{\mathrm{d}\tau }\right)^2.$$
(A1)
The quantity
$$\alpha g_{\mu \nu }\xi ^\mu u^\nu =\left(\frac{2rM}{2r+M}\right)^2\left(\frac{\mathrm{d}t}{\mathrm{d}\tau }\right)$$
(A2)
is a constant of the motion, where $`\xi ^\mu =(1,0,0,0)`$ is the static Killing field. The energy density observed by this observer is
$$T_{\mu \nu }u^\mu u^\nu =(u\varphi )^2\left(\frac{1}{2}\varphi \varphi +V[\varphi ]\right).$$
(A3)
Since we consider the static configuration and the spatial part of the metric is non-singular in the isotropic coordinates, the second term of Eq. (A3) is always finite. We have
$$u\varphi =\left(1+\frac{M}{2r}\right)^2\left[\alpha ^2\left(\frac{2r+M}{2rM}\right)^21\right]^{1/2}\frac{\varphi }{r}$$
(A4)
for the static axi-symmetric configuration $`\varphi (r,\vartheta )`$. Thus the requirement that Eq. (A3) is finite at the horizon is reduced to $`\varphi /r=0`$ at $`r=M/2`$, or equivalently Eq. (17).
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# hep-th/0005193 On Brane-World Cosmology
## I Introduction
The string-theory inspired idea of a four-dimensional universe as a brane world embedded into a multi-dimensional universe has now become very popular. In many papers, models of such kind are considered in the four-dimensional cosmological context; in particular, the theory of cosmological perturbations within the frames of such models is currently being developed (for a comprehensive list of references see, e.g., ). In several papers , an unusual law of cosmological expansion of a four-dimensional universe embedded into (or bounding ) a five-dimensional space has been reported. According to this law, the energy density of matter on the brane enters quadratically the right-hand side of the new equations for the brane world, in contrast with the standard cosmology, where it enters the similar equations linearly. Such a behaviour might strongly modify the standard cosmological model . There exist several approaches to solving this problem. Thus, for example, in it was shown that the standard cosmological evolution can be recovered in models with a large five-dimensional cosmological constant and brane tension and, in , it was also pointed out that, in the presence of a mechanism for stabilization of the radius of the extra dimension, the cosmologies are generically of the FriedmannโRobertsonโWalker ordinary four-dimensional form.
However, the basic equations of the theory were thus far obtained in the absence of curvature-dependent terms in the action for the four-dimensional brane world. The purpose of this article is to derive the field equations for a brane world embedded into (or bounding) a five-dimensional space in the case where such curvature-dependent terms are present. After that, we consider a particular example of a cosmological situation and discuss some of its solutions.
## II Field equations
Consider a theory with a four-dimensional hypersurface (brane) $`\mathrm{\Sigma }`$ which is the boundary of a five-dimensional manifold $``$.In fact, all the equations of this section are valid for a $`n`$-dimensional space $``$ bounded by a $`(n1)`$-dimensional hypersurface $`\mathrm{\Sigma }`$. We take the action of the theory to have the natural form
$$S=M_5^3\left[_{}\left({}_{}{}^{(5)}R2\mathrm{\Lambda }_5\right)+2_\mathrm{\Sigma }K\right]+_{}L_5(g_{ab},\mathrm{\Phi })+M_4^2_\mathrm{\Sigma }\left({}_{}{}^{(4)}R2\mathrm{\Lambda }_4\right)+_\mathrm{\Sigma }L_4(h_{ab},\varphi ).$$
(1)
Here, $`{}_{}{}^{(5)}R`$ is the scalar curvature of the Lorentzian five-dimensional metric $`g_{ab}`$ on $``$, and $`{}_{}{}^{(4)}R`$ is the scalar curvature of the induced metric $`h_{ab}=g_{ab}n_an_b`$ on $`\mathrm{\Sigma }`$, where $`n^a`$ is the vector field of the outer unit normal to $`\mathrm{\Sigma }`$. The boundary $`\mathrm{\Sigma }`$ is assumed to be timelike, so that the vector field $`n^a`$ is spacelike. The quantity $`K=K_{ab}h^{ab}`$ is the trace of the symmetric tensor of extrinsic curvature $`K_{ab}=h^c{}_{a}{}^{}_{c}^{}n_b`$ of $`\mathrm{\Sigma }`$ in $``$. The symbols $`L_5(g_{ab},\mathrm{\Phi })`$ and $`L_4(h_{ab},\varphi )`$ denote, respectively, the Lagrangian densities of the five-dimensional matter fields $`\mathrm{\Phi }`$ and of the four-dimensional matter fields $`\varphi `$ whose dynamics is restricted to the boundary $`\mathrm{\Sigma }`$ so that they interact only with the induced metric $`h_{ab}`$. Note that some of the fields $`\varphi `$ in principle may represent restrictions of some of the fields $`\mathrm{\Phi }`$ to the boundary $`\mathrm{\Sigma }`$. All integrations over $``$ and over $`\mathrm{\Sigma }`$ are taken, respectively, with the natural volume elements $`\sqrt{g}d^5x`$ and $`\sqrt{h}d^4x`$, where $`g`$ and $`h`$ are, respectively, the determinants of the matrices of components of the metric on $``$ and of the induced metric on $`\mathrm{\Sigma }`$ in a coordinate basis. The symbols $`M_n`$ and $`\mathrm{\Lambda }_n`$ denote, respectively, the $`n`$-dimensional Planck mass and cosmological constant.
In this paper, we freely use the notation and conventions of . In particular, following Sec. 10.2 of , we use the one-to-one correspondence between tensors in $`\mathrm{\Sigma }`$ and tensors in $``$ that are invariant under projection to the tangent space to $`\mathrm{\Sigma }`$, i.e., tensors $`T^{a_1\mathrm{}a_k}_{b_1\mathrm{}b_l}`$ such that
$$T^{a_1\mathrm{}a_k}{}_{b_1\mathrm{}b_l}{}^{}=h^{a_1}{}_{c_1}{}^{}\mathrm{}h^{a_k}{}_{c_k}{}^{}h_{b_1}^{}{}_{}{}^{d_1}\mathrm{}h_{b_l}{}_{}{}^{d_l}T_{}^{c_1\mathrm{}c_k}{}_{d_1\mathrm{}d_l}{}^{}.$$
(2)
The third term in (1) containing the four-dimensional curvature is often missing from the action, or its contribution is missing from the equations of motion. However, in general, this term seems to be essential since it is generated as a quantum correction to the matter action in (1). Note that this quantum correction typically involves an infinite number of terms of higher order in curvature (a similar situation in the context of the AdS/CFT correspondence is described in ). Thus, we assume that such terms are present and retain only the lowest-order ones in (1).
In this paper, we are interested only in the metric equation, so the Lagrangians for the matter fields $`\mathrm{\Phi }`$ and $`\varphi `$ will not be specified. The first variation of action (1) with respect to the metric $`g_{ab}`$ is equal toThose interested in the derivation of (3) may look into the appendix.
$`\delta S`$ $`=`$ $`M_5^3{\displaystyle _{}}\left({}_{}{}^{(5)}G_{ab}^{}+\mathrm{\Lambda }_5g_{ab}\right)\delta g^{ab}{\displaystyle _{}}T_{ab}\delta g^{ab}+M_5^3{\displaystyle _\mathrm{\Sigma }}S_{ab}\delta h^{ab}`$ (3)
$`+`$ $`M_4^2{\displaystyle _\mathrm{\Sigma }}\left({}_{}{}^{(4)}G_{ab}^{}+\mathrm{\Lambda }_4h_{ab}\right)\delta h^{ab}{\displaystyle _\mathrm{\Sigma }}\tau _{ab}\delta h^{ab},`$ (4)
where $`{}_{}{}^{(n)}G_{ab}^{}`$ denotes the $`n`$-dimensional Einsteinโs tensor, $`S_{ab}K_{ab}Kh_{ab}`$, and $`T_{ab}`$ and $`\tau _{ab}`$ define, respectively, the five-dimensional and four-dimensional stress-energy tensors of matter. Note that the variation $`\delta h^{ab}`$ on the brane is not an independent quantity, but is completely and uniquely determined by $`\delta g^{ab}`$. For simplicity, we also assume that the Lagrangian $`L_5(g_{ab},\mathrm{\Phi })`$ does not contain derivatives of the metric $`g_{ab}`$, the presence of which might contribute to the surface terms in (3).
On an extremal field configuration, variation (3) is equal to zero for arbitrary variations of the metric $`g_{ab}`$. Considering variations that leave the induced metric on $`\mathrm{\Sigma }`$ intact, i.e., for which $`h^a{}_{c}{}^{}h_{}^{b}{}_{d}{}^{}\delta h^{cd}0`$ and, hence, the surface integrals in (3) vanish, we obtain the equation of motion in the five-dimensional bulk:
$${}_{}{}^{(5)}G_{ab}^{}+\mathrm{\Lambda }_5g_{ab}=\frac{1}{M_5^3}T_{ab}.$$
(5)
It is important to stress that the GibbonsโHawking boundary term \[the second term in the square brackets in action (1)\] is required to obtain this equation in a consistent way . Now, considering arbitrary variations of the metric $`g_{ab}`$ and taking into account equation (5), we obtain the equation of motion on the boundary $`\mathrm{\Sigma }`$ in the form
$${}_{}{}^{(4)}G_{ab}^{}+\mathrm{\Lambda }_4h_{ab}+MS_{ab}=\frac{1}{M_4^2}\tau _{ab},$$
(6)
where $`M=M_5^3/M_4^2`$. It is the presence of the tensor $`S_{ab}`$ in the equation of motion (6) that makes the dynamics on the brane $`\mathrm{\Sigma }`$ unusual.
One of the GaussโCodacci relations, namely,
$$D_aS^a{}_{b}{}^{}={}_{}{}^{(5)}R_{cd}^{}n^dh^c{}_{b}{}^{},$$
(7)
where $`D_a`$ is the (unique) derivative on the brane $`\mathrm{\Sigma }`$ associated with the induced metric $`h_{ab}`$, together with equation (6) and the bulk equation (5) imply the relation
$$D_a\tau ^a{}_{b}{}^{}=T_{cd}n^dh^c{}_{b}{}^{}.$$
(8)
Thus, the four-dimensional stress-energy tensor is covariantly conserved if and only if the right-hand side of (8) is vanishing at the brane, in particular, if the stress-energy tensor $`T_{ab}`$ of the five-dimensional matter is a linear combination of $`g_{ab}`$ and $`h_{ab}`$ at the brane.
Thus far, we considered $`\mathrm{\Sigma }`$ to be a boundary of a five-dimensional manifold $``$. However, the theory can easily be extended to the case where $`\mathrm{\Sigma }`$ is embedded into $``$. In this case, it can be regarded as a common boundary of two pieces $`_1`$ and $`_2`$ of $``$, and we should simply add the actions of the form of the first term in (1) for these two pieces. In varying the resulting action, we must respect the condition that the metrics induced on the brane $`\mathrm{\Sigma }`$ by the metrics of these two pieces coincide; however, the extrinsic curvatures of $`\mathrm{\Sigma }`$ in $`_1`$ and in $`_2`$ are allowed to be different. Equation (5) remains valid in the bulk, and equation (6) will be modified to
$${}_{}{}^{(4)}G_{ab}^{}+\mathrm{\Lambda }_4h_{ab}+M\left(S_{ab}^{(1)}+S_{ab}^{(2)}\right)=\frac{1}{M_4^2}\tau _{ab},$$
(9)
where the tensors $`S_{ab}^{(1)}`$ and $`S_{ab}^{(2)}`$ are constructed with the use of the respective extrinsic curvatures. The analog of equation (8) also can easily be derived.
If several branes are embedded into the manifold $``$, equations of the form (9) are valid for each of these branes. If a brane is a common boundary of more than two bulk manifolds, the corresponding number of similar terms will be present inside the brackets of (9).
Neglecting the third term in action (1) amounts to taking the limit of $`M_40`$. In this case, equation (9) reduces to the Israelโs junction condition
$$M_5^3\left(S_{ab}^{(1)}+S_{ab}^{(2)}\right)=\tau _{ab}.$$
(10)
On the other hand, taking the limit of $`M_50`$ is equivalent to setting $`M=0`$ in equation (6) or (9). In this limiting case, the influence of the five-dimensional bulk vanishes, and we obtain the standard equation of general relativity.
## III Cosmological example
As an illustration, we consider a particular example of the cosmological situation described in . We take the five-dimensional metric in the static spherically-symmetric form
$$ds_5^2=f(r)dt^2+dr^2/f(r)+r^2d\mathrm{\Omega }_{(3)},$$
(11)
where $`d\mathrm{\Omega }_{(3)}`$ is the metric of the unit three-sphere. In the case of vanishing stress-energy tensor of the five-dimensional matter, as a solution of (5) we can take
$$f(r)=1\alpha r^2,$$
(12)
where $`\alpha =\mathrm{\Lambda }_5/6`$. The brane $`\mathrm{\Sigma }`$ is taken to be spherically-symmetrically embedded into this manifold according to the law $`r=a(t)`$. We then discard the exterior $`r>a(t)`$ and consider the resulting space with $`\mathrm{\Sigma }`$ as a boundary. The tensor $`S_{ab}`$ for this boundary can easily be calculated. Its nonzero components have the form (see also )
$$S_0^0=\frac{3\sqrt{f(a)+\dot{a}^2}}{a},S_j^i=\delta _j^i\frac{1}{a^2\dot{a}}\frac{d}{d\tau }\left(a^2\sqrt{f(a)+\dot{a}^2}\right),$$
(13)
where $`i,j=1,2,3`$ label the coordinates on the unit three-sphere and the overdot denotes the derivative with respect to the cosmological time $`\tau `$ on the brane, in terms of which the induced metric is given by the line element
$$ds_4^2=\left[f(a)\frac{1}{f(a)}\left(\frac{da}{dt}\right)^2\right]dt^2+a^2d\mathrm{\Omega }_{(3)}=d\tau ^2+a^2d\mathrm{\Omega }_{(3)}.$$
(14)
Equation (6) then yields the following two equations:
$$\frac{1+\dot{a}^2}{a^2}+\frac{M\sqrt{f(a)+\dot{a}^2}}{a}=\lambda +\kappa \rho ,$$
(15)
$$\frac{2\ddot{a}}{a}+\frac{1+\dot{a}^2}{a^2}+\frac{M}{a^2\dot{a}}\frac{d}{d\tau }\left(a^2\sqrt{f(a)+\dot{a}^2}\right)=3\lambda 3\kappa p,$$
(16)
where we made the notation $`\lambda =\mathrm{\Lambda }_4/3`$, $`\kappa =1/3M_4^2`$, and $`\rho `$ and $`p`$ denote the standard components of the four-dimensional stress-energy tensor (that may include its own cosmological-constant contribution). Introducing the Hubble parameter $`H\dot{a}/a`$, we have from (15)
$$H^2+\frac{f(a)}{a^2}=\left[\left(\frac{M^2}{4}+\frac{f(a)1}{a^2}+\lambda +\kappa \rho \right)^{1/2}\frac{M}{2}\right]^2.$$
(17)
Substituting this into equation (16), we obtain the conservation law
$$\dot{\rho }+3H(\rho +p)=0,$$
(18)
in accordance with the general equation (8). Together with the equation of state (e.g., in the form $`p=w\rho `$), equations (17), (18) constitute a closed system of cosmological equations.
If we neglect the third term in action (1) and thus take the limit of $`M_40`$, equation (17) will become
$$H^2+\frac{f(a)}{a^2}=\left(\frac{\rho }{3M_5^3}\right)^2,$$
(19)
with the quadratic dependence of the right-hand side on the energy density, noted previously . This result is approximately valid for finite values of $`M_4`$ provided $`\left|(f(a)1)/a^2+\lambda \right|\kappa \rho M^2`$.
In the limit of $`M_50`$, which corresponds to $`M0`$, the influence of the five-dimensional bulk vanishes, and we recover the equation of the standard cosmology
$$H^2+\frac{1}{a^2}=\lambda +\kappa \rho .$$
(20)
This equation is approximately valid for finite values of $`M_5`$ provided $`M^2(f(a)1)/a^2+\lambda +\kappa \rho `$.
We see that, in principle, the standard regime (20) might be realised at the early stages of the universe expansion, the standard theory of nucleosynthesis thus remaining intact, while, at later stages, the universe might evolve according to the nonstandard law (19).
The system of equations (15), (16) with the function $`f(r)`$ given by (12) admits a solution of the static empty (that is, $`\rho =p=0`$) universe with the scale factor
$$a_0=\left(\alpha +\frac{M^2}{4}\right)^{1/2},$$
(21)
provided the value of $`M_5`$ (hence, also of $`M`$) is negative and the values of parameters are tuned so that
$$\lambda =\alpha \frac{M^2}{4}.$$
(22)
Note that, in the case of negative $`\mathrm{\Lambda }_5`$, solution (11), (12) describes the anti-de Sitter five-dimensional space. In this case, the value of $`\alpha =\mathrm{\Lambda }_5/6`$ is negative, so that the size of the four-dimensional static universe given by (21) may be arbitrarily large, because the value of $`\alpha +M^2/4`$ may be arbitrarily close to zero. This is another instance of fine-tuning, similar to that discussed in , that countervails the effect of the four-dimensional cosmological constant which, according to (22), must be negative in the present case.
It should be noted, however, that, under relation (22), the system of equations (15), (16) is degenerate at the point
$$a=a_0,\dot{a}=0,$$
(23)
in the sense that the second-order derivative term in (16) vanishes at this point. This circumstance, in particular, has a consequence that the Cauchy problem (23) for system (15), (16) is ill-defined since, besides the static solution $`aa_0`$, it also has the solution
$$a(\tau )=a_0\mathrm{cosh}\left(\frac{\tau }{a_0}\right)$$
(24)
describing the de Sitter spacetime.
Gravitational perturbations of solutions to theory (1) and the resulting effective four-dimensional theory of gravity remain to be investigated. In this respect, some important results in the context of the AdS/CFT correspondence are discussed in .
## Note added
After this paper was already submitted to the LANL archive, the author became aware of the paper , in which all the essential results of the present work have been obtained.
## Acknowledgments
This work was supported in part by the Foundation of Fundamental Research of the Ministry of Science of Ukraine under grant No. 2.5.1/003.
## Variation of the action for gravity
Here, we derive the expression for the first variation of the action for gravity
$$S_g=_{}R+2_\mathrm{\Sigma }K.$$
(25)
Equations of this section will be valid for arbitrary dimension of spacetime. In contrast with the standard derivation, here we do not assume that the variation of $`g_{ab}`$ vanishes at the boundary $`\mathrm{\Sigma }`$, which is taken to be timelike.
We start from the standard expression (see, e.g., Appendix E of )
$$\delta \left(_{}R\right)=_{}G_{ab}\delta g^{ab}+_{}^av_a,$$
(26)
where
$$v_a=^b\left(\delta g_{ab}\right)g^{cd}_a\left(\delta g_{cd}\right).$$
(27)
The second integral in (26) can be transformed with the use of the Stokes theorem as
$$_{}^av_a=_\mathrm{\Sigma }v_an^a,$$
(28)
where
$$v_an^a=n^ag^{bc}\left[_c\left(\delta g_{ab}\right)_a\left(\delta g_{bc}\right)\right]=n^ah^{bc}\left[_c\left(\delta g_{ab}\right)_a\left(\delta g_{bc}\right)\right].$$
(29)
Then we have
$$\delta K=\delta \left(h^a{}_{b}{}^{}_{a}^{}n^b\right)=\delta h^a{}_{b}{}^{}_{a}^{}n^b+h^a{}_{b}{}^{}(\delta C)_{}^{b}{}_{ac}{}^{}n_{}^{c}+h^a{}_{b}{}^{}_{a}^{}\delta n^b,$$
(30)
where
$$(\delta C)^b{}_{ac}{}^{}=\frac{1}{2}g^{bd}[_a\left(\delta g_{cd}\right)+_c\left(\delta g_{ad}\right)_d\left(\delta g_{ac}\right)].$$
(31)
The first term in the right-hand side of (30) is identically zero. Indeed, we have $`\delta n_a=n_an_b\delta n^b`$, so that
$`\delta h^a{}_{b}{}^{}_{a}^{}n^b=(\delta n^an_b+n^a\delta n_b)_an^b=(\delta n^an^an_c\delta n^c)n_b_an^b`$ (32)
$`=\delta n^ch^a{}_{c}{}^{}n_{b}^{}_an^b=\delta n^cn_bK_c{}_{}{}^{b}=0.`$ (33)
Thus, variation of the second term of (25) is
$$\delta \left(2_\mathrm{\Sigma }K\right)=_\mathrm{\Sigma }[n^ch^{ab}_c\left(\delta g_{ab}\right)+2h^a{}_{b}{}^{}_{a}^{}\delta n^bKh_{ab}\delta h^{ab}],$$
(34)
where the last term in the square brackets stems from the variation of the volume element $`\sqrt{h}d^4x`$ in the integral over $`\mathrm{\Sigma }`$.
The total boundary term in the variation of action (25) is given by the sum of (28) and (34) with the result
$$(\text{Boundary term})=_\mathrm{\Sigma }[n^ah^{bc}_c\left(\delta g_{ab}\right)+2h^a{}_{b}{}^{}_{a}^{}\delta n^bKh_{ab}\delta h^{ab}].$$
(35)
We transform the first term in the integrand of the last expression:
$$n^ah^{bc}_c\left(\delta g_{ab}\right)=h^{bc}_c\left(n^a\delta g_{ab}\right)h^{bc}_cn^a\delta g_{ab}=h^{bc}_c\left(n^a\delta g_{ab}\right)+K_{ab}\delta h^{ab}.$$
(36)
Then
$$(\text{Boundary term})=_\mathrm{\Sigma }[h^{bc}_c\left(n^a\delta g_{ab}\right)+2h^a{}_{b}{}^{}_{a}^{}\delta n^b]+_\mathrm{\Sigma }(K_{ab}Kh_{ab})\delta h^{ab}.$$
(37)
Now we show that the integrand of the first integral in (37) is a divergence, so that this integral vanishes for variations of $`g_{ab}`$ with compact support in $`\mathrm{\Sigma }`$. Indeed,
$`h^{bc}_c\left(n^a\delta g_{ab}\right)+2h^a{}_{b}{}^{}_{a}^{}\delta n^b=h^{bc}_c(\delta n_bg_{ab}\delta n^a)+2h^a{}_{b}{}^{}_{a}^{}\delta n^b`$ (38)
$`=h^a{}_{b}{}^{}_{a}^{}(g^{bc}\delta n_c+\delta n^b)=h^a{}_{b}{}^{}_{a}^{}\left(h^b{}_{c}{}^{}\delta n^c\right)=D_b\left(h^b{}_{c}{}^{}\delta n^c\right),`$ (39)
where $`D_a`$ is the (unique) derivative on $`\mathrm{\Sigma }`$ associated with the induced metric $`h_{ab}`$, and the last equality in (38) is valid by virtue of Lemma 10.2.1 of .
As a final result, we obtain
$$\delta S_g=_{}G_{ab}\delta g^{ab}+_\mathrm{\Sigma }\left(K_{ab}Kh_{ab}\right)\delta h^{ab}.$$
(40)
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# CAN INSTANTONS SATURATE THE LARGE MASS OF THE ๐' MESON IN THE CHIRAL LIMIT?
## I Introduction
It is well known that one of the most important aspects of the famous $`U(1)`$ problem is the large mass of the $`\eta ^{}`$ meson. It does not vanish in the chiral limit, so the $`\eta ^{}`$ meson is not the Nambu-Goldstone (NG) boson. In Ref. (see also Ref. ) by using the large $`N_c`$ limit technique the expression for the mass of the $`\eta ^{}`$ meson was derived, namely
$$m_\eta ^{}^2=\frac{2N_f}{F_\pi ^2}\chi _t+\mathrm{\Delta },$$
(1)
where $`\mathrm{\Delta }=2m_K^2m_\eta ^2`$, $`N_f`$ is the number of light quarks and $`F_\pi `$ is the pion decay constant. However, the important quantity which enters this formula is the topological density operator (topological susceptibility), $`\chi _t`$ (for definition see section 3). In the chiral limit it is screened that is why it is defined for Yang-Mills (YM) fields, i.e., for pure gluodynamics ($`N_f=0`$). It is one of the main characteristics of the QCD nonperturbative vacuum where it measures the fluctuations of the topological charge.
The precise validity of the Witten-Veneziano (WV) formula (1.1) is, of course, not completely clear because of its origin. Nevertheless, let us regard it as exact for simplicity (in any case we have nothing better than Eq. (1.1). However, there are phenomenological reasons as well as some lattice indications to believe that QCD is close to $`SU(\mathrm{})`$). Using now experimental values of all physical quantities entering this formula, one obtains that the phenomenological (โexperimentalโ) value of the topological susceptibility is
$$\chi _t^{phen}=0.001058GeV^4=(180.36MeV)^4=0.1377GeV/fm^3.$$
(2)
In the chiral limit $`\mathrm{\Delta }=0`$ since $`K^\pm `$ and $`\eta `$ particles are NG excitations. It is worth noting further that neither the mass of the $`\eta ^{}`$ meson nor the pion decay constant in the chiral limit cannot exeed their experimental values. So the WV formula (1.1) provides an absolute lower bounds for the pion decay constant and the mass of the $`\eta ^{}`$ meson in the chiral limit, namely
$$854m_\eta ^{}^0957.77(MeV),$$
(3)
and
$$83.2F_\pi ^093.3(MeV),$$
(4)
respectively. They should be compared with their experimental values (upper bounds in the previous expressions). Let us note that the chiral perturbation theory value of the pion decay constant in the chiral limit, $`F_\pi ^0=(88.3\pm 1.1)MeV`$ , obviously satisfies these bounds, Eq. (1.4). Recent lattice result (see also brief review ) for the mass of the $`\eta ^{}`$ meson in the continuum chiral limit is $`m_\eta ^{}^0=863(86)MeV`$. It obviously satisfies our bounds, Eq. (1.3), as it should be.
One can conclude in that the mass of the $`\eta ^{}`$ meson remains large even in the chiral limit, which is real problem indeed. Thus the large mass of the $`\eta ^{}`$ meson in the chiral limit is due to the phenomenological value of the topological susceptibility. In other words, it is clear that through the topological susceptibility (i.e., via the WV formula (1.1)) the large mass of the $`\eta ^{}`$ meson even in the chiral limit is determined by the topological properties of the QCD ground state, its nonperturbative vacuum. It has a very rich dynamical and topological structure \[10-12\]. It is a very complicated medium and its dynamical and topological complexity means that its structure can be organized at various levels (quantum, classical). It can contain many different components and ingredients which may contribute to the truly nonperturbative vacuum energy density (VED). It is well known that VED in general is badly divergent , however the truly nonperturbative VED is finite, automatically negative and it has no imaginary part (stable vacuum). For gauge-invariant definition and concrete examples see recent papers . Precisely this quantity is one of the main characteristics of the QCD ground state and precisely it is related to the nonperturbative gluon condensate via the trace anomaly relation (see section 2) as well as to the above-mentioned topological susceptibility via the low energy โtheoremโ (relation) derived by Novikov, Schifman, Vanshtein and Zakharov (NSVZ) a long time ego and rederived quite recently by Halperin and Zhitnitsky (HZ) (see section 3). Let us remind that the truly nonperturbative VED is nothing else but the bag constant apart from the sign, by definition . It is much more general quantity than the string tension because it is relevant for light quarks as well.
Many models of the QCD vacuum involve some extra classical color field configurations (such as randomly oriented domains of constant color magnetic fields, background gauge fields, averaged over spin and color, stochastic colored background fields, etc.) and ingredients such as color-magnetic and Abelian-projected monopoles (see Refs. and references therein). The relevance of center vortices for QCD vacuum by both lattice and analytical methods was recently investigated as well.
However, the most elaborated classical models are the random and interacting instanton liquid models (RILM and IILM) of the QCD vacuum. They are based on the existence of the topologically nontrivial, instanton-type fluctuations of gluon fields there, which are nonperturbative, weak coupling limit solutions to the classical equations of motion in Euclidean space (and references therein). That is instantons may be $`qualitatively`$ responsible for the $`\eta ^{}`$ mass for the first time has been pointed out by โt Hooft . The WV formula (1.1) clearly shows that the topological susceptibility due to instantons should be nonzero. Here we would like to investigate this problem $`quantitatively`$ as well. In this paper we precisely address the problem whether instantons can $`quantitatively`$ saturate the large mass of the $`\eta ^{}`$ meson in the chiral limit or not. The reason is that the instanton-type fluctuations in the QCD vacuum are totally suppressed in the chiral limit and they are again restored due to dynamical breakdown of chiral symmetry (and references therein). Thus an important question immediately arises, namely is this restoration in terms of instanton number densities $`quantitatively`$ sufficient for the above-mentioned purpose or not.
In order to make the problem more transparent, let us elaborate on this point im more detail. There exists a โstandardโ, rather quilitative estimate of the topological susceptibility in terms of the instanton number density valid for dilute system of weakly interacting charges, namely $`\chi _t(N/V)`$ (see, for example a brief review in Ref. ). For its phenomenological value, $`n=(N/V)=1fm^4`$, one obtains $`\chi _t0.0015GeV^4`$, which rather roughtly reproduces the phenomenological value of the topological susceptibility shown in Eq. (1.2). However, in the $`chiral`$ limit it is impossible to use the phenomenological value of the instanton number density. As was mentioned above, the instanton-type fluctuations are strongly suppressed in the $`chiral`$ limit and it is obvious that the โrestoredโ value due to dynamical chiral symmetry breakdown (which is not due to instantons only) cannot be equal to its phenomenological value. Apparently it is substantially less. In any case one should look for a possible enhancement of the right hand side of the above-mentioned rough estimate in the chiral limit through the appropriate numerical coefficient in order to saturate the phenomenological value of the topological susceptibility in the left hand side of this estimate.<sup>*</sup><sup>*</sup>*It has been already well known to NSVZ that the right hand side of such type of relation should be in general somehow enhanced. However, their mechanism for a possible enhancement cannot be applied to the chiral limit case since it is completely opposite to the word with no light quarks (or, what is just the same, to the word were all quarks are very heavy). Moreover, in the same paper they argued that the gluon condensate in this limit by factor of two is less as compared to its phenomenological value (see discussion in sect. V below). The right hand side of this estimate in the chiral limit stands for $`n_0=(N/V)_0`$ (subscript โ0โ means the chiral limit) while the left hand side is the phenomenological topological susceptibility indeed (the quenched topological susceptibility in the chiral limit). Otherwise it is zero since topological charge is totally screened in the chiral limit. Thus the real problem arises how to find a correct $`quantitative`$ relation between the topological susceptibility and the instanton number density in the chiral limit. It becomes perfectly clear that more sophisticated methods are needed for this purpose. The above-discussed rough estimate even on account of a some well-justified value for $`(N/V)_0`$ (which is very doubtful) definately fails.
One of the main purposes in this paper is to develop a well-justified formalism in order to establish precisely the above-mentioned correct quantitative relation. For the readerโs convenience in sections 2, 3 and 4 we have modified an analytical formalism (along with lines of NSVZ paper ) which has been developed earlier . It is based on using the above-mentioned trace anomaly relation (section 2), low energy โtheoremโ (relation) (section 3) and WV formula (1.1) in the chiral limit (section 4). It allows one to directly calculate the topological susceptibility and the mass of the $`\eta ^{}`$ meson in the chiral limit as a functions of the truly nonperturbative YM VED due to instantons or equivalently as a functions of the instanton number density in the chiral limit. In section 5 we present its estimates in the chiral limit as it follows from phenomenology and lattice approach. The numerical results are shown in Tables I, II and III (section 6). Section 7 is devoted to discussion and conclusions.
## II The trace anomaly relation
The truly nonperturbative VED is important in its own right as one of the main characteristics of the QCD nonperturbative vacuum. Furthermore it assists in estimating such an important phenomenological parameter as the gluon condensate, introduced in the QCD sum rules approach to resonance physics . The famous trace anomaly relation in the general case (nonzero current quark masses $`m_f^0`$) is
$$\mathrm{\Theta }_{\mu \mu }=\frac{\beta (\alpha _s)}{4\alpha _s}G_{\mu \nu }^aG_{\mu \nu }^a+\underset{f}{}m_f^0\overline{q}_fq_f.$$
(5)
where $`\mathrm{\Theta }_{\mu \mu }`$ is the trace of the energy-momentum tensor and $`G_{\mu \nu }^a`$ being the gluon field strength tensor while $`\alpha _s=g^2/4\pi `$. Sandwiching Eq. (2.1) between vacuum states and on account of the obvious relation $`0|\mathrm{\Theta }_{\mu \mu }|0=4ฯต_t`$, one obtains
$$ฯต_t=\frac{1}{4}0|\frac{\beta (\alpha _s)}{4\alpha _s}G_{\mu \nu }^aG_{\mu \nu }^a|0+\frac{1}{4}\underset{f}{}m_f^00|\overline{q}_fq_f|0,$$
(6)
where $`ฯต_t`$ is the sum of all possible independent, truly nonperturbative contributions to VED (the total VED) and $`0|\overline{q}_fq_f|0`$ is the quark condensate. Since in what follows we want to saturate the total VED by instantons only, i.e., to put $`ฯต_t=ฯต_I+\mathrm{}`$, then it is legitimate to use weak coupling limit solution to the $`\beta `$-function
$$\beta (\alpha _s)=b\frac{\alpha _s^2}{2\pi }+O(\alpha _s^3),b=11\frac{2}{3}N_f,$$
(7)
where recalling $`N_f`$ is the number of light flavors. Then from Eq. (2.2) in the chiral limit ($`m_f^0=0`$), one obtains
$$ฯต_I=\frac{b}{4}\times \frac{1}{8}0|\frac{\alpha _s}{\pi }G_{\mu \nu }^aG_{\mu \nu }^a|0.$$
(8)
In general, of course, it is impossible to use the above-mentioned weak coupling limit solution to the $`\beta `$-function and one needs to introduce a new quantity, namely the gluon condensate in the strong coupling limit . Let us also emphasize that in the chiral limit the truly nonperturbative VED is nothing else but the gluon condensate apart from the overall numerical factor (see Eqs. (2.2) and (2.4)).
## III The topological susceptibility
One of the main characteristics of the QCD nonperturbative vacuum is the topological density operator (topological susceptibility) in gluodynamics ($`N_f=0`$)
$$\chi _t=\underset{q0}{lim}id^4xe^{iqx}\frac{1}{N_c^2}0|T\left\{q(x)q(0)\right\}|0,$$
(9)
where $`q(x)`$ is the topological charge density, defined as $`q(x)=(\alpha _s/4\pi )F(x)\stackrel{~}{F}(x)=(\alpha _s/4\pi )F_{\mu \nu }^a(x)\stackrel{~}{F}_{\mu \nu }^a(x)`$ and $`\stackrel{~}{F}_{\mu \nu }^a(x)=(1/2)ฯต^{\mu \nu \rho \sigma }F_{\rho \sigma }^a(x)`$ is the dual gluon field strength tensor while $`N_c`$ is the number of different colors. In the definition of the topological susceptibility (3.1) it is assumed that the corresponding regularization and subtraction of all types of the perturbative contributions have been already done in order Eq. (3.1) to stand for the renormalized, finite and the truly nonperturbative topological susceptibility (see Refs. ). The anomaly equation in the WV notations is
$$_\mu J_5^\mu =2N_f(2/N_c)(\alpha _s/4\pi )F\stackrel{~}{F}.$$
(10)
The topological susceptibility can be related to the nonperturbative gluon condensate via the low energy โtheoremโ in gluodynamics proposed by NSVZ (by using the dominance of self-dual fields hypothesis in the YM vacuum) as follows:
$$\underset{q0}{lim}id^4xe^{iqx}0|T\left\{\frac{\alpha _s}{8\pi }G\stackrel{~}{G}(x)\frac{\alpha _s}{8\pi }G\stackrel{~}{G}(0)\right\}|0=\frac{1}{2b}0|\frac{\alpha _s}{\pi }G_{\mu \nu }^aG_{\mu \nu }^a|0.$$
(11)
Quite recently it was discussed by HZ in Ref. (see also references therein) who noticed that it is not precisely a Ward identity, but rather is a relation between the corresponding correlation functions, indeed. That is why in what follows we call Eq. (3.3) as low energy relation or NSVZ-HZ relation. The anomaly equation in the NSVZ-HZ notations is
$$_\mu J_5^\mu =N_f(\alpha _s/4\pi )G\stackrel{~}{G},$$
(12)
with $`N_f=3`$. Thus to get the topological susceptibility in the WV form from this relation, it is necessary to make a replacement in its left hand side as follows: $`G\stackrel{~}{G}(2/N_c)F\stackrel{~}{F}`$. Then the WV topological susceptibility (3.1) finally becomes
$$\chi _t=\frac{1}{2b}0|\frac{\alpha _s}{\pi }G_{\mu \nu }^aG_{\mu \nu }^a|0=(\frac{4}{b})^2ฯต_{YM},$$
(13)
where $`b=11`$. The second equality comes from Eq. (2.4) by denoting the truly nonperturbative VED due to instantons at $`N_f=0`$ as $`ฯต_{YM}`$. The significance of this formula is that it gives the topological susceptibility as a function of the truly nonperturbative VED for pure gluodynamics, $`ฯต_{YM}`$.
## IV The $`U(1)`$ problem
The topological susceptibility (3.1) assists in the resolution of the above-mentioned $`U(1)`$ problem via the WV formula for the mass of the $`\eta ^{}`$ meson (1.1). Within our notations it is expressed as follows $`f_\eta ^{}^2m_\eta ^{}^2=4N_f\chi _t`$, where $`f_\eta ^{}`$ is the $`\eta ^{}`$ residue defined as $`0|_{q=u,d,s}\overline{q}\gamma _\mu \gamma _5q|\eta ^{}=i\sqrt{N_f}f_\eta ^{}p_\mu `$ and $`0|N_f\frac{\alpha _s}{4\pi }G\stackrel{~}{G}|\eta ^{}=(N_c\sqrt{N_f}/2)f_\eta ^{}m_\eta ^{}^2`$. Using the normalization relation $`f_\eta ^{}=\sqrt{2}F_\pi ^0`$, one finally obtains
$$F_\pi ^2m_\eta ^{}^2=2N_f\chi _t.$$
(14)
Eq. (3.5) then leads to
$$m_\eta ^{}^2=2N_f\left(\frac{4}{bF_\pi }\right)^2ฯต_{YM},$$
(15)
which expresses the mass of the $`\eta ^{}`$ meson as a function of the truly nonperturbative YM VED, in particularly due to instantons in this case. In previous expressions we omit for simplicity the superscript โ0โ in the pion decay constant as well as in $`m_\eta ^{}^2`$.
In order to directly apply this formalism to the instanton liquid model we need the realistic estimate of the corresponding truly nonperturbative chiral VED in this model.
## V The truly nonperturbative VED due to instantons
The contribution to the truly nonperturbative VED at the classical level from the instanton-type nonperturbative fluctuations of gluon fields can be estimated as follows. Let us consider the trace anomaly relation (2.4) in the chiral limit again. The phenomenological analysis of QCD sum rules for the gluon condensate implies
$$G_2G^2\frac{\alpha _s}{\pi }G^20|\frac{\alpha _s}{\pi }G_{\mu \nu }^aG_{\mu \nu }^a|00.012GeV^4,$$
(16)
which can be changed within a factor of two . From the phenomenological estimate (5.1), one easily can calculate
$$\frac{1}{8}0|(\alpha _s/\pi )G_{\mu \nu }^aG_{\mu \nu }^a|00.0015GeV^41.0fm^4.$$
(17)
Having in mind this and assuming that the gluon condensate in the weak coupling limit is determined by the instanton-type fluctuations only, Shuryak (see also references therein) has concluded in that the โaverage separationโ between instantons was $`1.0fm`$, so the corresponding density of the instanton-type fluctuations should be $`1.0fm^4`$. Let us note that the second parameter of the instanton liquid model of the QCD vacuum, the instanton size $`\rho _01/3`$, was chosen to reproduce standard (also (as gluon condensate) phenomenologically estimated from QCD sum rules ) value of the quark condensate. This contribution to VED via the trace anomaly relation (2.1-2.2) vanishes in the chiral limit. However, due to all reasonable estimates of light quark masses, numerically its contribution is at $`20\%`$ and thus comparable with the systematic error in the determination of the gluon condensate itself .
Using the above-mentioned estimate (5.2), from Eq. (2.4) for dilute ensemble, one finally obtains
$$ฯต_I=\frac{b}{4}\times n_0,$$
(18)
where, let us remind, we denote the instanton number density in the chiral limit as $`n_0=(N/V)_0`$ while in the general (nonchiral case) it is $`n=(N/V)`$. It is well known that density of instanton-type fluctuations is strongly suppressed in the chiral limit and is again restored bacause of dynamical breakdown of chiral symmetry (see Ref. and references therein). In any case it can not be large in the chiral limit, so the functional dependence of VED on the instanton number density in the chiral limit, established in Eq. (5.3) due to dilute gas approximation, seems to be justified in this case. The only problem is its numerical value which, in general of course, can not be equal to its phenomenological value, $`1.0fm^4`$. Let us emphasize that the instanton contribution to the truly nonperturbative VED was $`notcalculated`$ independently but was $`postulated`$ via the trace anomaly relation (2.2) using the phenomenological value of the gluon condensate (5.1) as well as weak coupling limit solution to the $`\beta `$-function (2.3). The significance of this expression is that it allows one to estimate instantonโs contribution to the truly nonperturbative VED as a function of the instanton number density in the chiral limit, which can be estimated either from phenomenology or taken from lattice simulations (see below).
In Ref. it has been argued that the gluon condensate in the chiral limit is approximately two times less than the above-mentioned its phenomenological (empirical) value (5.1), i. e. $`G^2_{ch}0.5G^2_{phen}`$. This means that in this case the instanton number density in the chiral limit is $`n_00.5fm^4`$. However, it has been already pointed out that QCD sum rules substantially underestimate the value of the gluon condensate. The most recent phenomenological calculation of the gluon condensate is given by Narison in Ref. , where a brief review of many previous calculations is also presented. His analysis leads to the update average value as
$$0|\frac{\alpha _s}{\pi }G_{\mu \nu }^aG_{\mu \nu }^a|0=(0.0226\pm 0.0029)GeV^4.$$
(19)
This means that instanton density is approximately two time bigger than it was estimated by Shuryak for instanton liquid model , but in the chiral limit we are again left with $`n_01.0fm^4`$. In Ref. the dimensionless instanton density has been evaluated through the QCD asymptotic scale parameter, $`\mathrm{\Lambda }_{QCD}`$. At a value $`\mathrm{\Lambda }_{QCD}=280MeV`$ estimated in the $`\overline{MS}`$ scheme for $`N_f=3`$ (compatible with DIS and other data), one gets the instanton density about $`1.0fm^4`$. Thus in the chiral limit we are left with about half of this number again. Unfortunately, the conversion to physical units is rather ambiguous in the pure gauge theory. It strongly depends on the renormalization scheme chosen for calculation.
In lattice QCD situation with instanton density and their sizes is also ambigious. In quenched ($`N_f=0`$) lattice QCD by using the so-called โcoolingโ method the role of the instanton-type fluctuations in the QCD vacuum was investigated . In particular, it was found that the instanton density should be $`n=(1+\delta )fm^4`$, where $`\delta 0.30.6`$ depending on cooling steps. Moreover, by studying the topological content of the vacuum of $`SU(2)`$ pure gauge theory using a method of RG mapping , it is concluded that the average radius of an instanton is about $`0.2fm`$, at a density of about $`2fm^4`$. However, in Ref. the topological content of the $`SU(3)`$ vacuum was studied using the same method as for $`SU(2)`$ gauge theory earlier and was obtained a fair agreement with Shuryakโs phenomenologically estimated numbers for the instanton liquid model. At the same time, in Refs. \[34-36\] considerably larger values were reported and advocated. Thus at this stage it is rather difficult to choose some well-justified numerical value for the instanton number density in the chiral limit. If $`n2fm^4`$, then we are left with half of this value in the chiral limit, i.e., $`n_01fm^4`$, but if $`n1fm^4`$, we will be left with half of this, i.e., $`n_00.5fm^4`$. That is why we will perform all our calculations on account of the two different values for the instanton number densities in the chiral limit, $`n_0=0.5fm^4,1.0fm^4`$.
In conclusion, let us note that for densities $`n>2fm^4`$ (which means $`n_0>1fm^4`$) the applicability of the dilute gas approximation becomes doubtful.
## VI Numerical results
Lattice approach shows that in instanton calculus much more convenient to express all results as a functions of the instanton number density, $`n`$ . Thus let us also express all our results as a functions of the instanton number density in the chiral limit, $`n_0`$. In this case, the truly nonperturbative VED is given in Eq. (5.3) and its numerical results are shown in Table I. It explicitly depends on $`N_f`$ and instanton number density $`n_0`$. The quark part is one order of magnitude less than gluon part and is of opposite sign. Combining Eqs. (2.4) and (5.3), one obtains
$$G_2\frac{\alpha _s}{\pi }G^2=\frac{32}{b}ฯต_t\frac{32}{b}ฯต_I=8\times n_0,$$
(20)
i.e., the gluon condensate in the weak coupling limit does not explicitly depend on $`N_f`$ (for numerical results see Table I as well). As it was mentioned above, precisely this gluon condensate was introduced a long time ago . This unphysical situation takes place because in instanton calculus there is no other way to estimate the truly nonperturbative VED than the trace anomaly relation (2.4) which becomes finally Eq. (6.1) as it was described above. In this case it is preferable to have the $`N_f`$ dependent VED than the gluon condensate since the former mainly characterisizes the detail structure of the nonperturbative vacuum while the latter is one of its average (global) characteristics, indeed.
The numerical values of the topological susceptibility due to instanton number density in the chiral limit are shown in Table II. They are determinded by the following relation
$$\chi _t=\frac{4}{b}\times n_0,b=11,$$
(21)
which, obviously follows from Eqs. (3.5) and (5.3).
The numerical values of the $`\eta ^{}`$ meson in the chiral limit due to instanton number density in the chiral limit are shown in Table III. They are determinded by the following relation
$$m_\eta ^{}^2=\frac{8N_f}{b}\frac{1}{F_\pi ^2}\times n_0,b=11,N_f=3,$$
(22)
which follows from Eqs. (4.2) and again (5.3). As in previous expression (4.2) we omit for simplicity the superscript โ0โ in the pion decay constant as well as in $`m_\eta ^{}^2`$.
## VII Discussion and conclusions
In summary, using the trace anomaly relation (2.4), The NSVZ-HZ low energy relation (3.3) and WV formula for the mass of the $`\eta ^{}`$ meson (4.1), the chiral topology of the QCD nonperturbative instanton vacuum has been numerically evaluated. The NSVZ-HZ low energy relation (3.3) (which is valid in the $`chiral`$ limit as well (with the corresponding value of the gluon condensate in this limit)) in combination with the trace anomaly relation in the chiral limit (2.4) provides precisely the well-justified formalism to find the above-mentioned correct $`quantitative`$ relation between the topological susceptibility and the instanton nunber density in the chiral limit, Eq. (6.2). It turned out that the correct coefficient in this relation is even substantially less than one ($`4/b`$ with $`b=11`$), so there is no way for instantons to saturate the phenomenological value of the topological susceptibility in the chiral limit.
There are, of course, very strong theoretical reasons to think that the low energy relation (3.3) is very important. Let us remind the redear that behind this low energy theorem-relation is a beatiful physical idea , namely that the self-dual gluon fields may be dominant in the YM vacuum. This idea is not only beatiful but it is very powerful as well. As it is shown here (see also Ref. ), it allows to relate many important quantities such as the truly nonperturbative VED, gluon condensate, topological susceptibility, density numbers of different types of excitations and fluctuations (not only those of instantons) of the gluon field configurations, etc. to each other in a well-justified way, i.e., by providing well-justified coefficients between the above-mentioned quantities of the same dimension. At the same time, the rough estimate discussed in the Introduction is rather artificial. In fact, its right hand side has been adjusted by Shuryak by using the pnenomenological value of the gluon condensate which precisely via the relation (3.3) is related to the topological susceptibility. Thus precisely it justifies (and not vice versa) the existence of the above-mentioned estimate. At the same time, it is, of course, much more general since it can be applied to the chiral limit while the estimate certainly fails in this limit as it has been explained in the Introduction.
The topological susceptibility due to instanton number densities in the chiral limit (Table II) is substantially less than its phenomenological value shown in Eq. (1.2). Therefore they (instantons) cannot account for the large mass of the $`\eta ^{}`$ meson in the chiral limit alone (compare bounds (1.3) with numerical results presented in Table III). This means that the truly nonperturbative VED (the bag constant, apart from the sign, by definition) due to instantons (see Table I) is an order of magnitude less than it is required to saturate the large mass of the $`\eta ^{}`$ meson in the chiral limit as well as the phenomenological value of the topological susceptibility. In turn, this means that larger (than it can be reasonably estimated from phenomenology or lattice simulations) value of the instanton number density in the chiral limit is required for this purpose. This is impossible since in the chiral limit it cannot be large as it was underlined a few times in the main body of the text.
It is well known from recent and old lattice calculations that instantons by themselves do not confine quarks contributing only no more than ten percent into the phenomenological value of the string tension. It has been already known for a long time that instantons give rise to the constant (not linear) potential at large distances in continuum theory . Obviously, if instantons are not able to saturate the large mass of the $`\eta ^{}`$ meson in the chiral limit (see Table III and Eq. (1.2)), this means that they have nothing to do with its experimental value. Thus one can conclude in that instantons provide resolution of the large mass of the $`\eta ^{}`$ meson, indeed, but $`onlyquilitatively`$. Their $`quantitative`$ contribution is not clearly sufficient.
If instantons by themselves cannot saturate the phenomenological value of the topological susceptibility (and consequently the large mass of the $`\eta ^{}`$ meson) in continuum theory, then a question immediately arises how they can saturate these quantities in lattice calculations. There are no doubts left that all lattice data which nicely saturate the phenomenological value of the topological susceptibility (see for example, a brief recent review in Ref. ) clearly indicate the existence and importance of such types of the nonperturbative excitations of gluon field configurations in the QCD vacuum which $`cannot`$ be treated as instantons. The presense of other (than instantons) nonperturbative fluctuations and excitations in the QCD true ground state has been pointed out and discussed in Ref. as well.
Concluding, let us emphasize once more that instantons by themselves do not confine quarks and cannot saturate the phenomenological value of the topological susceptibility either. In order to explain confinement and saturate the large mass of the $`\eta ^{}`$ meson in the chiral limit, we need completely $`differentfrominstantons`$ types of the nonperturbative excitations of gluon field configurations in the QCD true vacuum. These types of vacuum excitations should have at least an order of magnitude larger amplitudes than those of instantons can provide at all and they should be closely related to the excitations responsible for quark confinement .
Obvious simplicity of analytical calculations of the above-discussed quantities in comparison with obvious complexity of their calculations by lattice method (see, for example, Refs. ) should be also mentioned.
The author is grateful to M. Faber and A. Ivanov for many interesting remarks and useful discussion in the Institute of Nuclear Physics of the Thechnical University in Wien. It is also a pleasure to thank Gy. Kluge for many useful remarks and help.
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# Boundary conditions for the spinor field in Rindler spacetime and the quantum field theoretical basis of the Unruh effect
## 1 Introduction
The โUnruh effectโ could be expressed by the following statements: 1) the Minkowski vacuum state, from the point of view of an accelerated observer, is a particle state described by a density matrix at the temperature
$$T=\frac{a}{2\pi k_B}$$
(1)
called the Unruh-Davies temperature,where $`a`$ is the (constant) acceleration of the observer and $`k_B`$ is the Boltzmann constant; 2)an accelerated observer in the empty Minkowski space will detect a thermal bath of particles at the temperature 1.
As was stressed in (in particular is recommended for an especially clear account) the crucial point here is not that an accelerator detector (observer) would in someway react to the vacuum state of the field in Minkowski space, rather that this response is universal, i.e. independent from the structure of the detector itself, from the quantum field considered, and from the details of the interaction between them. This situation could remind the response of a probe massive body to the gravitational field, which is indeed universal. This universality is in this case determined by the pure geometrical nature of the gravitational field, so that the interaction between a body and the gravity is determined uniquely by the geometry of the spacetime in which the body moves, and not by its inner structure. If a similar universality would appear in the different context of accelerated observer in Minkowski space, this would mean that here only the quantum properties of Minkowski vacuum state matter.
Consequently, two problems are involved here, which in principle are different, but are claimed to be equivalent in the literature: the physical properties of a quantum field when restricted to a submanifold (the Rindler Space, (RS)) of Minkowski space (MS); the behaviour of a constantly accelerated particle detector in empty flat space.
The first one deals only with basic principles of quantum theory and appears to be more fundamental, whether the second one should in general involve also a description of structure and characteristics of the detector, and details of interaction with the quantum field. Again, it is the analysis of the first problem that made possible to claim for universality of the Unruh-like detector response.
We will treat here only the first problem above, which is a particular example of the analysis of the behaviour of a field in a submanifold of a maximally analytically extended manifold.
The procedure used by Unruh is based on a quantization scheme for a free field in MS, alternative but claimed to be equivalent to the standard one, which uses as the Hilbert space of solutions of the wave equation
$$_U=_R_L$$
(2)
where $`_R`$ consists of solutions which are non-zero everywhere but in the L sector, which have positive frequency with respect to the โRindler timeโ $`\eta `$, and which reduce themselves to the well known Fulling modes , and $`_L`$ is given by the solutions which are non-zero everywhere but in R and with negative frequency with respect to $`\eta `$. Then it is obtained a representation of the Minkowski vacuum state as a state in $`_s(_U)`$, i.e. the Fock space constructed on $`_U`$. Finally it is derived the particle content of Minkowski vacuum and the expression for the density matrix associated to it, when expressed as a mixed state in the Fock space $`_s(_R)`$ having traced out the degrees of freedom related to the L sector, which is unaccessible to the Rindler observer.
The usual explanation of the Unruh effect is based exactly on the presence, for a Rindler observer, which is confined inside the Rindler wedge, of an event horizon which prevent him from having part of the informations about the quantum field, so that he sees Minkowski vacuum state as a mixed state. But this explanation is (as was indicated in and also in our opinion) not entirely satisfying, for several reasons: on the one hand, the existence of horizons is due to overidealization of the problem, since for physical accelerations (which last finite amount of time) no horizon should be present, so that the response of an accelerated detector, if of the Unruh-Davies type, cannot be caused by it; on the other hand, in the purely quantum theoretical treatment of the problem, we would expect that the presence of event horizons affects deeply the fields, from the point of view of Rindler observer, in the form of some kind of boundary condition, which instead are totally absent in the Unruh scheme, and in the usual quantization in RS.
What we are going to do in this paper is to extend to the spinor field the results obtained in for the scalar field and to show the conditions to (and only to) which the quantum theory of the spinor field in the Rindler spacetime is well posed, to study the relationship between this construction and the usual one in MS, in order to find which role is played by these conditions in the derivation and interpretation of the Unruh effect in the spinor case, to analyse the Unruh quantization scheme and to understand finally what is its physical significance.
What we will find is that a correct quantization procedure for the spinor field in Rinlder space requires the boundary condition
$$\underset{\rho \mathrm{\hspace{0.17em}0}}{lim}\rho ^{\frac{1}{2}}\mathrm{\Psi }(0,\rho )=\mathrm{\hspace{0.17em}0}$$
(3)
i.e. the field should not grow up too rapidly at the origin of Rindler (and Minkowski) space. From a more general point of view, this means that the quantization on a background manifold which is not maximally extended requires a boundary condition which is absent in the quantization procedure over the extended manifold; in the literature this is not recognized clearly enough, and the usual procedure is to restrict the fields just considering a smaller domain of definition.
Moreover, we show that also for the spinor field the Unruh quantization scheme is not valid in the whole Minkowski space, but only in a sector of it, namely the double wedge $`RL`$. Consequently, the Unruh quantization implies the same kind of boundary condition, and cannot be used as a proof of the Unruh effect.
These results, as we said, represent a generalization to the spinor field of similar ones obtained for the scalar field (and also recently extended also to the electromagnetic field ); consequently, they appear to be consequences of very general properties of Quantum Field Theory, and seem to be firmly established. It is worth to note here that the Unruh effect is often identified (mainly by mathematical physicists) with the so-called Bisognano-Wichmann Theorem, in the context of the algebraic approach to quantum field theory. We wonโt deal in this paper with the algebraic approach, but the interested reader could found the analog of our result for the scalar case extended to the algebraic framework in . There it is shown that the physical interpretation of the Bisognano-Wichmann Theorem, in terms of accelerated observers in MS occurs in the same kind of problems encountered in the conventional approach, that we will now discuss.
In addition to the main results cited, we obtained some minor but original results, necessary to achieve the first, namely: the explicit expression of the Lorentz boost generator (or Lorentz Momentum) for the spinor field, its eigenfunctions and their analytical representation holding in the whole MS (except for the origin), which is in turn a generalization of the Gerlachโs Minkowski Bessel Modes .
## 2 Rindler Spacetime
Letโs consider a particle moving in MS with constant acceleration $`a`$ along the x-axis; it will follow the trajectory given by (parameter $`\tau `$):
$$t=a^1\mathrm{sinh}(a\tau )x=a^1\mathrm{cosh}(a\tau )y=y(0)z=z(0)$$
(4)
This is an hyperbola in the (t,x) plane. The lines $`t=\pm x`$ represent asymptotes for it and event horizons for the moving particle. Varying $`a`$ we obtain different hyperbolas with the same characteristics. Letโs now perform, starting from the Minkowski metric, the change in coordinates given by:
$`t=\rho \mathrm{sinh}\eta x=\rho \mathrm{cosh}\eta `$ (5)
$`\rho =\sqrt{x^2t^2}\eta =arctgh\left({\displaystyle \frac{t}{x}}\right)`$ (6)
with the other coordinates left unchanged. The metric assumes the form:
$$ds^2=\rho ^2d\eta ^2d\rho ^2dy^2dz^2$$
(7)
Note that it describes a stationary spacetime. The worldlines $`\rho =const,y=const,z=const`$ correspond to uniformly accelerated observers, with $`a=\rho ^1`$ and proper time $`\tau =\rho \eta `$, as it can be seen comparing 4 and 5. We can think at the Rindler Space as the collection of these worldlines, and this is the reason why RS is generally regarded as the โnaturalโmanifold in which to describe accelerate motion. The hypersurfaces $`\eta =const`$ describe events which are simultaneous from the point of view of a โRindler (uniformly accelerated) observerโ. The Rindler manifold cannot be extended to negative values of $`\rho `$ trough $`\rho =0`$, for RS is no more rigid beyond the limit $`\rho 0`$. This hypersurface, in fact, represents, as we said, an event horizon for Rindler observers, which cannot see any event located beyond it. Nevertheless the horizon is a regular surface, of course, and the metric singularity is due only to the choice of the coordinates.
The ( 5) cover only a sector of the whole MS and the others are covered by the following charts:
L :
$$t=\rho \mathrm{sinh}\eta x=\rho \mathrm{cosh}\eta $$
(8)
$$\rho =\sqrt{x^2t^2}\eta =arctgh(\frac{t}{x})$$
(9)
F :
$$t=\rho \mathrm{cosh}\eta x=\rho \mathrm{sinh}\eta $$
(10)
$$\rho =\sqrt{t^2x^2}\eta =arctgh(\frac{x}{t})$$
(11)
P :
$$t=\rho \mathrm{cosh}\eta x=\rho \mathrm{sinh}\eta $$
(12)
$$\rho =\sqrt{t^2x^2}\eta =arctgh(\frac{x}{t})$$
(13)
## 3 Quantization in Rindler space
We now turn to the problem of quantization of the spinor field in the Rindler wedge. The procedure will be the standard one, so we will first solve the Dirac equation looking for hamiltonian eigenfunctions.
### 3.1 The Dirac equation and its solutions in RS
Using the tetrad formalism, the Dirac equation in a generic curved spacetime is given by:
$$(i\gamma ^\mu _\mu m)\mathrm{\Psi }=\mathrm{\hspace{0.17em}0}$$
(14)
where: $`\gamma ^\mu =\theta _{\overline{\mu }}^\mu \gamma ^{\overline{\mu }}`$ are the analogous in curved spaces of the usual Dirac gamma matrices $`\gamma ^{\overline{\mu }}`$, and they satisfy:
$$\gamma ^\mu \gamma ^\nu +\gamma ^\nu \gamma ^\mu =g^{\mu \nu }$$
(15)
la $`\theta _{\overline{\mu }}^\mu `$ is the inverse of the tetrad vector $`\theta _\mu ^{\overline{\mu }}`$ with vectorial index $`\mu `$ and tetradic index $`\overline{\mu }`$; $`_\mu =_\mu \mathrm{\Gamma }_\mu `$ is the spinorial covariant derivative, defined in such a way that $`_\mu \mathrm{\Psi }`$ is a covariant vector;
$$\mathrm{\Gamma }_\mu =\frac{1}{8}\theta _{\overline{b}}^b\theta _{\overline{c}b;\mu }[\gamma ^{\overline{c}},\gamma ^{\overline{b}}]$$
(16)
is the spinorial connection.
Letโs consider the particular case of the Rindler metric:
$$ds^2=\rho ^2d\eta ^2d\rho ^2dy^2dz^2$$
(17)
The Dirac equation assumes the form:
$`\left(i_\eta +i\rho \gamma ^0\gamma ^1_\rho +i\rho \gamma ^0\gamma ^2_y+i\rho \gamma ^0\gamma ^3_z+{\displaystyle \frac{i}{2}}\gamma ^0\gamma ^1m\rho \gamma ^0\right)\mathrm{\Psi }=0`$ (18)
$`i_\eta \mathrm{\Psi }=\left(i\rho \alpha _i_i{\displaystyle \frac{1}{2}}i\alpha _1+m\rho \beta \right)\mathrm{\Psi }`$ (19)
that is a Shroedinger-like form with an hamiltonian given by:
$`H_R=i\rho \alpha _i_i{\displaystyle \frac{1}{2}}i\alpha _1+m\rho \beta =`$ (20)
$`=i\rho \gamma ^0\gamma ^1_\rho i\rho \gamma ^0\gamma ^2_yi\rho \gamma ^0\gamma ^3_z{\displaystyle \frac{1}{2}}i\gamma ^0\gamma ^1+m\rho \gamma ^0`$ (21)
We now look for solutions of the Dirac equation which are simultaneously eigenfunctions of the Rindler hamiltonian ( 21) and of the operators $`P_y`$ and $`P_z`$ (of course they should also have the correct behaviour at infinity). So we expect to find a degeneracy of these solutions, because we know that three operators are not sufficient to completely characterize the states of the spinor field. This degeneracy is however of no relevance for our pourposes, so we will not deal with it.
The solutions are:
$$\mathrm{\Psi }_1^R=N_{}\left(X_1^RK_{i\frac{1}{2}}(\kappa \rho )+Y_1^RK_{i+\frac{1}{2}}(\kappa \rho )\right)e^{i\eta }e^{ik_2y+ik_3z}$$
(22)
with:
$$X_1^R=\left(\begin{array}{c}k_3\\ i\left(k_2+im\right)\\ i\left(k_2+im\right)\\ k_3\end{array}\right)Y_1^R=\left(\begin{array}{c}0\\ i\kappa \\ i\kappa \\ 0\end{array}\right)$$
(23)
and
$$\mathrm{\Psi }_2^R=N_{}^{}\left(X_2^RK_{i\frac{1}{2}}(\kappa \rho )+Y_2^RK_{i+\frac{1}{2}}(\kappa \rho )\right)e^{i\eta }e^{ik_2y+ik_3z}$$
(24)
with:
$$X_2^R=\left(\begin{array}{c}0\\ i\kappa \\ i\kappa \\ 0\end{array}\right)Y_2^R=\left(\begin{array}{c}k_3\\ i\left(k_2im\right)\\ i\left(k_2im\right)\\ k_3\end{array}\right)$$
(25)
where $`k_2`$,$`k_3`$ are eigenvalues of $`P_2`$, $`P_3`$, $`\kappa =\sqrt{k_2+k_3+m^2}`$, $``$ is the eigenvalue of the hamiltonian, and $`N_{}`$, $`N_{}^{}`$ are normalization factors which are found to be:
$$N_{}=N_{}^{}=\frac{1}{4\pi ^2\sqrt{\kappa }}\sqrt{\mathrm{cosh}\pi }$$
(26)
### 3.2 Physical characterization of the solutions
We now want to understand better the physical nature of the solutions ( 22)( 24), and this means to characterize them in a clearer way than just saying they are eigenfunctions of the Rindler hamiltonian.
For this pourpose we need to find the expressions for the solutions in minkowski coordinates.
In order to do it, it is necessary to consider carefully the way in which spinors transform under coordinate transformations.
It is well known that also in curved spacetimes spinors are characterized by their transformation properties under the action of the Lorentz group, but restricting the attention to the local minkowskian neighbour of the point in which the spinors have to be calculated, i.e. considering local Lorentz transformations on the tangent space of each point. Consequently, under a general coordinate transformation, the spinor will undergo a Lorentz transformation, but with a โvelocity parameterโwhich will be a function of the coordinates; this local Lorentz transformation has to be determined linearizing the coordinate transformation in which we are interested. In our case (transformation between Minkowskian and Rindler coordinate systems 5) the result is that the spinor transformation is given by:
$$\mathrm{\Psi }(t,x)=S(\eta )\mathrm{\Psi }(\eta ,\rho )=\mathrm{exp}\left(\frac{1}{2}\gamma ^0\gamma ^1\eta \right)\mathrm{\Psi }(\eta ,\rho )$$
(27)
Note that the operator we found, $`S=\mathrm{exp}\left(\frac{1}{2}\gamma ^0\gamma ^1\eta \right)`$, has the form of an operator resulting from a Lorentz coordinate transformation, with โvelocityโ parameter $`\eta `$.
The solutions ( 22)( 24) in minkowskian coordinates take the form:
$`\mathrm{\Psi }_{i,}^R(t,x,y,z)=N_{}(X_i^RK_{i\frac{1}{2}}(\kappa \rho )e^{\left(i\frac{1}{2}\right)\eta }+`$
$`+Y_i^RK_{i+\frac{1}{2}}(\kappa \rho )e^{\left(i+\frac{1}{2}\right)\eta })e^{ik_2y+ik_3z}`$ (28)
It must be clear that these are again defined only in the RS, but expressed in minkowskian coordinates, and this is why the normalization factor is still $`N`$.
Now we can turn to the anticipated physical characterization of the solutions: they are found to be eigenfunctions of the Boost Generator Operator, or Lorentz Momentum. This could be aspected, because 1)$`\mathrm{\Psi }_i`$ were eigenfunctions of the Rindler Hamiltonian $`H_R`$ and 2)$`H_R`$ is precisely the Lorentz Boost Generator written in Rindler coordinates, since (3)the time evolution of a Rindler (uniformly accelerated) observer is properly a infinite succession of infinitesimal boost transformations.
It is easy to verify this statement, once known the Lorentz Momentum operator, and this in turn can be obtained from the classical theory of fields.
In fact, given the conserved quantity:
$$M^{0i}=d^3x\left[\left(x^0T^{i0}x^iT^{00}\right)+S^{0i0}\right]$$
(29)
corresponding to the invariance of the Lagrangian under boost transformations along the i-axis, (which are isometries of the Rindler spacetime), the explicit calculation of $`T^{\mu \nu }`$ and of $`S^{\mu \nu }`$ gives:
$$M^{01}=d^3x\mathrm{\Psi }^{}\left[it^i+x^i\left(i\gamma ^0\gamma ^i_im\gamma ^0\right)+\frac{i}{2}\gamma ^0\gamma ^i\right]\mathrm{\Psi }$$
(30)
Interpreting it as a mean value of a quantum operator, we obtain for the Lorentz Momentum the expression (for contravariant and covariant components):
$$M^{0i}=it_ix_i\left(i\gamma ^0\gamma ^j_j+m\gamma ^0\right)+\frac{i}{2}\gamma ^0\gamma ^i$$
(31)
$$M_{0i}=+it_i+x_i\left(i\gamma ^0\gamma ^j_j+m\gamma ^0\right)\frac{i}{2}\gamma ^0\gamma ^i$$
(32)
Given this expression it can be verified that our solutions are eigenfunctions of the operator $`M_{01}`$ with eigenvalue $``$, and this represent their physical characterization.
### 3.3 The second quantization in RS
We now possess all the necessary elements to perform the second quantization of the spinor field in RS, i.e. a set of normalized functions which are solutions of the equation of motion. We then expand the field in terms of them:
$`\mathrm{\Psi }^R(\eta ,\rho ,y,z)={\displaystyle \underset{i=1,2}{}}{\displaystyle _{\mathrm{}}^+\mathrm{}}๐{\displaystyle _{\mathrm{}}^+\mathrm{}}๐k_2{\displaystyle _{\mathrm{}}^+\mathrm{}}๐k_3a_i(k_2,k_3)\mathrm{\Psi }_{i,,k_2,k_3}(t,x,y,z)=`$ (33)
$`=`$ $`{\displaystyle \underset{i=1,2}{}}{\displaystyle _0^+\mathrm{}}d{\displaystyle _{\mathrm{}}^+\mathrm{}}dk_2{\displaystyle _{\mathrm{}}^+\mathrm{}}dk_3\times `$
$`\times `$ $`\left(a_{,i}(k_2,k_3)\mathrm{\Psi }_{i,,k_2,k_3}^R(\eta ,\rho ,y,z)+b_{,i}^{}(k_2,k_3)\mathrm{\Psi }_{i,,k_2,k_3}^R(\eta ,\rho ,y,z)\right)`$
The second quantization is now performed considering the coefficients $`a_{i,}`$ and $`b_{i,}`$ (and their hermitian coniugated) as operators, and requiring for their anticommutators:
$`\{a_{,i}(k_2,k_3),a_{^{},j}^{}(k_2^{},k_3^{})\}=\delta _{ij}\delta (^{})\delta (k_2k_2^{})\delta (k_3k_3^{})`$ (34)
$`\{b_{,i}(k_2,k_3),b_{^{},j}^{}(k_2^{},k_3^{})\}=\delta _{ij}\delta (^{})\delta (k_2k_2^{})\delta (k_3k_3^{})`$ (35)
$`\{a_{,i}(k_2,k_3),b_{^{},j}^{}(k_2^{},k_3^{})\}=\mathrm{\hspace{0.17em}0}i,j,,^{},k_2,k_2^{},k_3,k_3^{}`$ (36)
$`\{a_{,i}(k_2,k_3),b_{^{},j}(k_2^{},k_3^{})\}=\mathrm{\hspace{0.17em}0}i,j,,^{},k_2,k_2^{},k_3,k_3^{}`$ (37)
and the quantum states of the field are constructed from the Rindler vacuum state $`0_R`$, defined by:
$$a_{,i}(k_2,k_3)0_R=\mathrm{\hspace{0.17em}0}i,,k_2,k_3$$
(38)
Now we will study if this quantum construction is well posed and at which conditions.
### 3.4 Conditions for the quantization in RS
Letโs consider again the Rindler hamiltonian:
$`H_R=i\rho \alpha _i_i{\displaystyle \frac{1}{2}}i\alpha _1+m\rho \beta =`$ (39)
$`=`$ $`i\rho \gamma ^0\gamma ^1_\rho i\rho \gamma ^0\gamma ^2_yi\rho \gamma ^0\gamma ^3_z{\displaystyle \frac{1}{2}}i\gamma ^0\gamma ^1+m\rho \gamma ^0`$ (40)
and letโs check whether it represents an hermitian operator, that is a necessary and sufficient condition for the completeness and orthonormality of the modes used. We should verify the condition
$$(H_R\mathrm{\Phi },\mathrm{\Psi })=(\mathrm{\Phi },H_R\mathrm{\Psi })$$
(41)
with scalar product given by:
$$(\mathrm{\Phi },\mathrm{\Psi })=๐\mathrm{\Sigma }_\mu \overline{\varphi }\gamma ^\mu \psi =_{\mathrm{}}^+\mathrm{}๐y_{\mathrm{}}^+\mathrm{}๐z_0^+\mathrm{}๐\rho \mathrm{\Phi }^{}\mathrm{\Psi }$$
(42)
The explicit calculation shows that:
$$(H_R\mathrm{\Phi },\mathrm{\Psi })=(\mathrm{\Phi },H_R\mathrm{\Psi })+_{\mathrm{}}^+\mathrm{}๐y_{\mathrm{}}^+\mathrm{}๐z\left[i\mathrm{\Phi }^{}\alpha _1\rho \mathrm{\Psi }\right]_{\rho =\mathrm{\hspace{0.17em}0}}^{\rho =+\mathrm{}}$$
(43)
Then it is evident that the hermiticity of the hamiltonian is assured if and only if
$$\underset{\rho \mathrm{\hspace{0.17em}0}}{lim}\rho ^{\frac{1}{2}}\mathrm{\Psi }(\eta ,\rho ,y,z)=\mathrm{\hspace{0.17em}0}\eta $$
(44)
and of course with analogous condition at $`\rho +\mathrm{}`$, i.e. the usual requirement of vanishing of the fields at spatial infinity.
We emphasize that, since the field $`\mathrm{\Psi }`$ is to be considered as an operator-valued distribution, this condition should be interpreted in the weak sense, this meaning that every matrix element of the quantity $`\rho ^{\frac{1}{2}}\mathrm{\Psi }`$ calculated with respect to any pair of physical states has to go to zero as $`\rho 0`$.
We stress again that a similar condition was already found for the scalar field in RS , and for the vector field .
If the condition ( 44) is necessary for the hermiticity of the hamiltonian, it is expected to appear also in the analysis of the coefficients of the expansion ( 33). So we write the explicit expression of the coefficients $`a_{i,}(k_2,k_3)`$:
$`a_{,i}(k_2,k_3)=(\mathrm{\Psi }_{i,,k_2,k_3},\mathrm{\Psi })_R=`$
$`=`$ $`+{\displaystyle _{\mathrm{}}^+\mathrm{}}dy{\displaystyle _{\mathrm{}}^+\mathrm{}}dz{\displaystyle _0^+\mathrm{}}d\rho N_{}^{}\left[(X_i^{}K_{i+\frac{1}{2}}+Y_i^{}K_{i\frac{1}{2}})\mathrm{\Psi }\right]\times `$
$`\times `$ $`e^{(i\eta +ik_2y+ik_3z)}`$ (46)
and consider the behaviour of the $`K_{i\pm \frac{1}{2}}(\kappa \rho )`$ for $`\rho 0`$.
We have:
$`K_{i\pm \frac{1}{2}}(\kappa \rho ){\displaystyle \frac{\pi }{2}}{\displaystyle \frac{1}{\mathrm{sin}\left(\pi \left(i\pm \frac{1}{2}\right)\right)}}\times `$ (47)
$`\times `$ $`\left[{\displaystyle \frac{\kappa ^{\left(i\pm \frac{1}{2}\right)}}{\mathrm{\Gamma }\left(\left(i\pm \frac{1}{2}\right)+\mathrm{\hspace{0.17em}1}\right)}}\left({\displaystyle \frac{\rho }{2}}\right)^{\left(i\pm \frac{1}{2}\right)}{\displaystyle \frac{\kappa ^{\left(i\pm \frac{1}{2}\right)}}{\mathrm{\Gamma }\left(\left(i\pm \frac{1}{2}\right)+\mathrm{\hspace{0.17em}1}\right)}}\left({\displaystyle \frac{\rho }{2}}\right)^{\left(i\pm \frac{1}{2}\right)}\right]`$
We see that a divergence is present for $`\rho 0`$ (we note also that, for $`=0`$, we have the exact expression:
$$K_{\frac{1}{2}}(\kappa \rho )=\sqrt{\frac{\pi }{2\kappa \rho }}e^{\kappa \rho }$$
(48)
again with the same type of divergence).
It is so evident that, in order the integral ( 46) to converge, is necessary to require the boundary condition ( 44) on the field. Otherwise, the coefficients $`a_{i,}(k_2,k_3)`$ are not defined and consequently the fundamental operators of quantum field theory like energy or particle number, which are built in terms of the annihilation and creation operators, are similarly not defined.
We want now to prove the necessity of the condition ( 44) in an even more apparent way, i.e. showing that the requirement of finiteness of the mean value of the energy in a generic state implies indeed ( 44). We work for simplicity in the plane $`(\eta ,\rho )`$ Consider the state of the field $`g`$ given by
$$g=c^{}(g)0_R=\underset{i}{}_0^{\mathrm{}}\frac{d}{^{\frac{1}{2}}}g()c_{i,}^{}0_R$$
(49)
i.e. a generic sovrapposition of eigenstates of the hamiltonian with weight function $`g()`$, and consider also the related one-particle amplitude for the field in RS given by
$$\mathrm{\Psi }_g^R=_R0\mathrm{\Psi }_Rg=e^{iH_R\eta }\varphi _g$$
(50)
where
$$\varphi _g=\underset{i}{}_0^{\mathrm{}}\frac{d}{^{\frac{1}{2}}}\frac{g()}{4\pi ^2\sqrt{\kappa }}\sqrt{\mathrm{cosh}\pi }\left(X_i^RK_{i\frac{1}{2}}+Y_i^RK_{i+\frac{1}{2}}\right)$$
(51)
Consequently we have the following translation of physical requirements (normalization of the states and finiteness of the mean value of the energy) into mathematical requirements on the weight function $`g()`$:
$$gg=_0^{\mathrm{}}\frac{d}{}g()^2=\mathrm{\hspace{0.17em}1}$$
(52)
$`gHg=`$ (53)
$`=`$ $`{\displaystyle _0^{\mathrm{}}}๐\rho \varphi _g^{}H\varphi _g={\displaystyle _0^{\mathrm{}}}๐\rho \varphi _g^{}\left\{i\rho \alpha _1_\rho {\displaystyle \frac{i}{2}}\alpha _1+m\rho \beta \right\}\varphi _g=`$ (54)
$`=`$ $`\mathrm{\hspace{0.17em}2}{\displaystyle _0^{\mathrm{}}}๐g()^2<\mathrm{}`$ (55)
From the equation ( 54) it is not immediately manifest what should be the behaviour of the funtions $`\varphi _g`$ with respect to $`\rho `$, and in order to understand it we have to analyse directly the expression ( 51), but taking into proper account the physical requirements ( 52) and ( 55).
Letโs consider the quantity $`\rho ^{\frac{1}{2}}\varphi _g`$ for $`\rho 0`$, using the expression 51 and the formula 47 for the modified Bessel function for $`\rho 0`$. Simple manipulations lead to
$`\rho ^{\frac{1}{2}}\varphi _g_{\rho 0}{\displaystyle \underset{i}{}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d}{^{\frac{1}{2}}}}N_{}g(){\displaystyle \frac{\pi }{\sqrt{2\kappa }}}\times `$
$`\times `$ $`\{X_i^R{\displaystyle \frac{1}{\mathrm{sin}(\pi (i\frac{1}{2}))\mathrm{\Gamma }\left(\frac{1}{2}+i\right)}}\left({\displaystyle \frac{\kappa \rho }{2}}\right)^i+`$
$`+`$ $`Y_i^R{\displaystyle \frac{1}{\mathrm{sin}\left(\pi \left(i+\frac{1}{2}\right)\right)\mathrm{\Gamma }\left(\frac{1}{2}i\right)}}\left({\displaystyle \frac{\kappa \rho }{2}}\right)^i\}=`$
$`=`$ $`{\displaystyle \underset{i}{}}\left\{X_i^RG_A(\rho ,\kappa )+Y_i^RG_B(\rho ,\kappa )\right\}`$ (57)
Now we want to study in details these quantities $`G_A`$ and $`G_B`$; we will do an explicit calculation only for the first one, since the argument for the second one is analogous.
First letโs write $`G_A`$ as:
$`G_A={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d}{^{\frac{1}{2}}}}N_{}g(){\displaystyle \frac{\pi }{\sqrt{2\kappa }}}{\displaystyle \frac{1}{\mathrm{sin}(\pi (i\frac{1}{2}))\mathrm{\Gamma }\left(\frac{1}{2}+i\right)}}\left({\displaystyle \frac{\kappa \rho }{2}}\right)^i=`$
$`=G_{A1}+G_{A2}+G_{A3}`$ (58)
where we have just splitted the integration domain into three parts, i.e. $`(0,\mathrm{})=(0,_1)(_1,_2)(_2,\mathrm{})`$, with $`_1<<1`$ and $`_2>>1`$.
Consider the term $`G_{A3}`$. Using the explicit form of the normalization factor $`N_{}`$, the asymptotic expression for the Gamma function (for $`\mathrm{}`$), and basic formulas for hyperbolic functions, we obtain:
$$G_{A3}^2__2\mathrm{}\frac{1}{32\pi ^3\kappa ^2_2^3}__2^{\mathrm{}}๐g()^2$$
(59)
now, given the finiteness condition for the mean value of the energy ( 55), we have that $`G_{A3}`$ could be made as small as we want with sensible choices of $`_2\mathrm{}`$.
Consider the term $`G_{A2}`$. We can easily obtain the inequality:
$$G_{A2}__1^_2๐C()g()$$
(60)
where $`C()`$ is a non singular function in the interval of integration. Taking into account the inequality $`g()\frac{1}{2}(1+g()^2)`$, and the nomalization condition ( 52), it is easy to see that the integral above should converge. Now applying the Riemann-Lebesgue Lemma, we conclude that $`G_{A2}`$ vanishes for $`\rho 0`$.
Coming to the term $`G_{A1}`$, we have:
$`G_{A1}={\displaystyle _0^_1}d{\displaystyle \frac{1}{^{\frac{1}{2}}}}{\displaystyle \frac{1}{2^{\frac{5}{2}}\pi \kappa }}\sqrt{\mathrm{cosh}\pi }g(){\displaystyle \frac{1}{\mathrm{sin}(\pi (i\frac{1}{2}))\mathrm{\Gamma }(i+\frac{1}{2})}}\times `$
$`\times \left({\displaystyle \frac{\kappa \rho }{2}}\right)^i{\displaystyle \frac{1}{2^{\frac{3}{2}}\pi ^{\frac{3}{2}}\kappa }}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d}{^{\frac{1}{2}}}}g()\left({\displaystyle \frac{\kappa \rho }{2}}\right)^i`$ (61)
being $`_1<<1`$. Letโs restrict the calculation to the case in which $`g()`$ vanish with $`0`$ as a suitable high power of $``$, i.e.
$$g()a^\alpha \alpha \frac{1}{2}0$$
(62)
(the results could be generalised to other cases). Note that the vanishing of the weight function is required also by the condition ( 52).
We obtain the inequality:
$$G_{A1}\frac{a}{2^{\frac{3}{2}}\pi ^{\frac{3}{2}}\kappa }\frac{1}{\mathrm{ln}\frac{\kappa \rho }{2}}$$
(63)
so that $`G_{A1}0`$ as $`\rho 0`$.
Consequently we have proved that for the generic physical state ( 49) to be normalised and to have finite energy, the boundary condition ( 44) is necessary. In other words, since ( 52) and ( 55) imply $`\rho ^{\frac{1}{2}}\varphi _g0`$ for $`\rho 0`$, we know that, if this condition is not satisfied, then ( 52) or ( 55) doesnโt hold, and consequently the state $`g`$ is not a physical state.
### 3.5 Discussion
We have found that the quantum theory of the spinor field in Rindler space is well defined if and only if we have the condition ( 44). This condition means that the field has to be quantized in a different way in MS and RS, because the horizons are not only of a causal but also of a physical significance for it. It also means that the usual procedure of quantize the field in RS, which just restrict its domain of definition, is not correct, because this kind of restriction is not enough to have a well posed theory. If one would like to study the spinor field in RS and work with physical states and modes, the two possible ways are to construct suitable wave packets made as combinations of the modes $`\mathrm{\Psi }_i`$, or to consider from the beginning a constrained hamiltonian, different from $`H_R`$ which automatically assures that the condition ( 44) is satisfied. The crucial point is however that this condition prevents any relationship between the quantization in RS and that in MS, and it means that RS should be treated as an manifold on its own, and not as a submanifold of MS (the Rindler wedge), i.e. the two quantizations define two different physical systems. This is because in the latter case we would not be free to choose any boundary condition for the field at the origin, because the state of the field would be determined from the beginning as a state in $`_M`$ (the Minkowski vacuum state $`0_M`$). Consequently, it is not possible to describe any state of the spinor field quantized in Minkowski space, and then restricted to the Rindler wedge, as a state in $`_R`$, and we could exspect that the necessity of boundary conditions on the field would manifest itself also in the analysis of the Unruh effect, which concerns exactly this relationship between RS and MS quantum constructions. To show this will be our next problem.
## 4 Analysis of the Unruh effect
In order to analyse this relationship, it is convenient to perform a quantization in Minkowski spacetime which is different from the standard one (that in plane waves) and which could be more easily compared with the Rindler one; namely, a quantization in terms of the Lorentz Momentum eigenfunctions defined in the whole Minkowski space, that represents analytical continuation of the $`\mathrm{\Psi }_{i,}`$ we found before, spinorial anologous of the Fulling modes for the scalar field, and that reduce themselves to these when restricted to the Rindler wedge., A physical motivation for this choice could be seen in the fact that the trajectories of a Rindler observer are the orbits of the Lorentz group, and that the R sector of MS is left invariant by the action of this group. The determination of these functions will require some preliminary steps.
### 4.1 Solutions in the other sectors and their relationship
First of all we will study the relations between the solutions of Dirac equation in the different sectors F, L, P (see 2). Of course, we omit the passages and just write down the form of Dirac equation in the sector and its possible solutions . That is the following.
F sector:
$$\left(i\gamma ^0_\rho +\frac{i}{\rho }\gamma ^1_\eta +\frac{i}{2}\frac{1}{\rho }\gamma ^0+i\gamma ^2_y+i\gamma ^3_zm\right)\mathrm{\Psi }=\mathrm{\hspace{0.17em}0}$$
(64)
$$\mathrm{\Psi }_i^{F,\pm }=M_i^\pm \left(X_i^FK_{i\frac{1}{2}}(\pm i\kappa \rho )+Y_i^FK_{i+\frac{1}{2}}(\pm i\kappa \rho )\right)e^{i\eta }e^{ik_2y+ik_3z}$$
(65)
with $`i=1,2`$ and
$$X_1^F=\left(\begin{array}{c}k_3\\ i\left(k_2+im\right)\\ i\left(k_2+im\right)\\ k_3\end{array}\right)Y_1^F=\left(\begin{array}{c}0\\ \kappa \\ \pm \kappa \\ 0\end{array}\right)$$
(66)
$$X_2^F=\left(\begin{array}{c}0\\ \pm \kappa \\ \kappa \\ 0\end{array}\right)Y_2^F=\left(\begin{array}{c}k_3\\ i\left(k_2im\right)\\ i\left(k_2im\right)\\ k_3\end{array}\right)$$
(67)
L sector:
$$\left(i\frac{1}{\rho }\gamma ^0_\eta +i\gamma ^1_\rho i\gamma ^2_yi\gamma ^3_z+\frac{1}{2}i\frac{1}{\rho }\gamma ^1+m\right)\mathrm{\Psi }=\mathrm{\hspace{0.17em}0}$$
(68)
$$\mathrm{\Psi }_i^{L,\pm }=O_i^\pm \left(X_i^LK_{i\frac{1}{2}}(\kappa \rho )+Y_i^LK_{i+\frac{1}{2}}(\kappa \rho )\right)e^{i\eta }e^{ik_2y+ik_3z}$$
(69)
with
$$X_1^L=\left(\begin{array}{c}k_3\\ i\left(k_2+im\right)\\ i\left(k_2+im\right)\\ k_3\end{array}\right)Y_1^L=\left(\begin{array}{c}0\\ i\kappa \\ i\kappa \\ 0\end{array}\right)$$
(70)
$$X_2^L=\left(\begin{array}{c}0\\ i\kappa \\ i\kappa \\ 0\end{array}\right)Y_2^L=\left(\begin{array}{c}k_3\\ i\left(k_2im\right)\\ i\left(k_2im\right)\\ k_3\end{array}\right)$$
(71)
P sector:
$$\left(i\gamma ^0_\rho +\frac{i}{\rho }\gamma ^1_\eta +\frac{i}{2}\frac{1}{\rho }\gamma ^0i\gamma ^2_yi\gamma ^3_z+m\right)\mathrm{\Psi }=\mathrm{\hspace{0.17em}0}$$
(72)
$$\mathrm{\Psi }_i^{P,\pm }=P_i^\pm \left(X_i^PK_{i\frac{1}{2}}(\pm i\kappa \rho )+Y_i^PK_{i+\frac{1}{2}}(\pm i\kappa \rho )\right)e^{i\eta }e^{ik_2y+ik_3z}$$
(73)
with
$$X_1^P=\left(\begin{array}{c}k_3\\ i\left(k_2+im\right)\\ i\left(k_2+im\right)\\ k_3\end{array}\right)Y_1^P=\left(\begin{array}{c}0\\ \pm \kappa \\ \kappa \\ 0\end{array}\right)$$
(74)
$$X_2^P=\left(\begin{array}{c}0\\ \kappa \\ \kappa \\ 0\end{array}\right)Y_2^P=\left(\begin{array}{c}k_3\\ i\left(k_2im\right)\\ i\left(k_2im\right)\\ k_3\end{array}\right)$$
(75)
As regards to the relationship between these solutions, this could be found by analytically continuing the functions across the event horizons, which represent branch points for these functions. Using the variables
$$x_+=x+tx_{}=tx$$
(76)
the passages trough the horizons is given by the substitutions: $`x_{}x_{}e^{\pm i\pi }`$ ($`RF`$) $`x_+x_+e^{\pm i\pi }`$ ($`FL`$), $`x_{}x_{}e^{\pm i\pi }`$ ($`LP`$), $`x_+x_+e^{\pm i\pi }`$ ($`PR`$). The result is that the solutions of the Dirac equation are linked by two possible paths of analytical continuation, namely that one corresponding (we name it A) to the transformation $`x_\pm x_\pm e^{+i\pi }`$ and linking in succession the functions
$$\mathrm{\Psi }_{i,}^R\mathrm{\Psi }_{i,}^{F,+}\mathrm{\Psi }_{i,}^{L,+}\mathrm{\Psi }_{i,}^{P,}\mathrm{\Psi }_{i,}^R$$
(77)
and that one corresponding to the transformation $`x_\pm x_\pm e^{i\pi }`$ (we name it B) and linking in succession the functions
$$\mathrm{\Psi }_{i,}^R\mathrm{\Psi }_{i,}^{F,}\mathrm{\Psi }_{i,}^{L,}\mathrm{\Psi }_{i,}^{P,+}\mathrm{\Psi }_{i,}^R$$
(78)
Moreover, it is possible to demonstrate that the normalization factors of the different solutions are related each one other in such a way that, once determined $`N`$, all the others are determined consequently.
### 4.2 Lorentz Momentum eigenfunctions in MS
We now turn to the problem of finding a unified representation for the eigenfunctions of the Boost Generator. Letโs consider the integral representations of the Bessel functions $`K_\nu (\rho )`$ given by:
$`K_\nu (\rho )={\displaystyle \frac{1}{2}}e^{\frac{i\pi \nu }{2}}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{i\rho \mathrm{sinh}\vartheta }e^{\nu \vartheta }๐\vartheta `$ (79)
$`K_\nu (\rho )={\displaystyle \frac{1}{2}}e^{\frac{i\pi \nu }{2}}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{i\rho \mathrm{sinh}\vartheta }e^{\nu \vartheta }๐\vartheta `$ (80)
Using the second one to express in an integral form and in minkowski coordinates the functions $`K_\nu (\rho (t,x))e^{\nu \eta (t,x)}`$ (using the coordinate transformation 5), we obtain:
$`K_\nu (\kappa \rho )e^{\nu \eta }=`$ (81)
$`=`$ $`{\displaystyle \frac{1}{2}}e^{\frac{i\pi \nu }{2}}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{i\kappa \rho \mathrm{sinh}\vartheta }e^{\nu \left(\vartheta +\eta \right)}๐\vartheta =`$ (82)
$`=`$ $`{\displaystyle \frac{1}{2}}e^{\frac{i\pi \nu }{2}}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{i\kappa \rho \mathrm{sinh}\left(\vartheta \eta \right)}e^{\nu \vartheta }๐\vartheta =`$ (83)
$`=`$ $`{\displaystyle \frac{1}{2}}e^{\frac{i\pi \nu }{2}}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{\left[i\kappa \rho \mathrm{sinh}\vartheta \mathrm{cosh}\eta i\kappa \rho \mathrm{cosh}\vartheta \mathrm{sinh}\eta \right]}e^{\nu \vartheta }๐\vartheta =`$ (84)
$`=`$ $`{\displaystyle \frac{1}{2}}e^{\frac{i\pi \nu }{2}}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{i\kappa \left[x\mathrm{sinh}\vartheta t\mathrm{cosh}\vartheta \right]}e^{\nu \vartheta }๐\vartheta =`$ (85)
$`=`$ $`{\displaystyle \frac{1}{2}}e^{\frac{i\pi \nu }{2}}{\displaystyle _{\mathrm{}}^+\mathrm{}}P_\vartheta ^{}(t,x)e^{\nu \vartheta }๐\vartheta `$ (86)
where $`P_\vartheta ^{}(t,x)`$ represent (2-dim) positive frequency plane waves with $`\omega =\kappa \mathrm{cosh}\vartheta `$ and $`k_x=\kappa \mathrm{sinh}\vartheta `$. Using the first one, and with the same procedure, we have:
$`K_\nu (\kappa \rho )e^{\nu \eta }={\displaystyle \frac{1}{2}}e^{\frac{i\pi \nu }{2}}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{i\kappa \left[x\mathrm{sinh}\vartheta +t\mathrm{cosh}\vartheta \right]}e^{\nu \vartheta }๐\vartheta =`$ (87)
$`={\displaystyle \frac{1}{2}}e^{\frac{i\pi \nu }{2}}{\displaystyle _{\mathrm{}}^+\mathrm{}}P_\vartheta ^+(t,x)e^{\nu \vartheta }๐\vartheta `$ (88)
where now $`P_\vartheta ^+(t,x)`$ are (2-dim) plane waves with negative frequency $`\omega =\kappa \mathrm{cosh}\vartheta `$. What we found means that the functions $`K_\nu (\rho (t,x))e^{\nu \eta (t,x)}`$ can be expressed equivalently as linear (integral) combination of positive or negative two dimensional plane waves.
Inserting these formulas into the expression for $`\mathrm{\Psi }_{i,}^R(t,x,y,z)`$, not taking care of the normalization factor, gives
$`\mathrm{\Psi }_i^{}(t,x,y,z)={\displaystyle \frac{1}{2}}N^{}[X_i^Re^{\pm \frac{i\pi }{2}\left(i\frac{1}{2}\right)}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{i\kappa \left[x\mathrm{sinh}\vartheta t\mathrm{cosh}\vartheta \right]}e^{\left(i\frac{1}{2}\right)\vartheta }d\vartheta +`$
$`+Y_i^Re^{\pm \frac{i\pi }{2}\left(i+\frac{1}{2}\right)}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{i\kappa \left[x\mathrm{sinh}\vartheta t\mathrm{cosh}\vartheta \right]}e^{\left(i+\frac{1}{2}\right)\vartheta }d\vartheta ]\times `$
$`\times e^{ik_2y+ik_3z}`$ (89)
These are the global functions we were looking for. For them, the following properties hold true:
* they are well behaved (analytical) on the entire Minkowski manifold, except for the origin;
* they are solutions of the Dirac equation;
* they are eigenfunctions of the boost generator operator $`M_{01}`$, with eigenvalue $``$;
* they reduce themselves to the correct solutions of the Dirac equations in the different sectors;
* they correspond each to one of the two possible paths of analytical continuation we mentioned before, namely $`\mathrm{\Psi }_i^{}`$ corresponds to the path A, and $`\mathrm{\Psi }_i^+`$ corresponds to the path B, so we could say that the reason for the existence of two different global representation is the existence of two different paths of analytical continuation across the horizons;
* they are orthonormalized with respect to the ordinary scalar product in MS, with normalization factors given by:
$$N^{}=\frac{e^{\pm \frac{1}{2}\pi }}{2\pi \sqrt{\kappa }}$$
(90)
We note also that these functions represent the analogous in the spinor case of the Gerlachโs Minkowski Bessel modes for the scalar field . Before considering the quantization of the field in terms of these modes, it is worth to notice that we could expect to have additional difficulties in using them instead of the standard plane wave basis. The reason for that is the divergence of these modes in the origin of Minkowski space, which doesโt affect, of course, the quantization procedure in any deep way, as will be proved in the following, but requires additional care in the calculations.
### 4.3 Alternative quantization in MS
Having obtained these global functions, we can perform the quantization of the spinor field in terms of them. We remember that the usual plane wave expansion is given by:
$$\mathrm{\Psi }=\underset{r=\mathrm{\hspace{0.17em}1},2}{}d^3k\left[a_r(k)\mathrm{\Psi }_r^+(k)+b_r^{}(k)\mathrm{\Psi }_r^{}(k)\right]$$
(91)
where the $`\mathrm{\Psi }_r^+(k)`$ are positive frequency plane waves and $`\mathrm{\Psi }_r^{}(k)`$ are negative frequency ones, and that the quantum vacuum state is defined by the relation:
$$a_r(k)0_M=b_r(k)0_M=\mathrm{\hspace{0.17em}0}r,\stackrel{}{k}$$
(92)
But we know that our $`\mathrm{\Psi }_i^{}`$ are linear combinations of positive frequency plane waves and $`\mathrm{\Psi }_i^+`$ of negative frequency ones, so we can have this kind of expansion:
$$\mathrm{\Psi }(t,x,y,z)=\underset{i=1,2}{}_{\mathrm{}}^+\mathrm{}๐๐k_2๐k_3\left[c_i(k_2,k_3)\mathrm{\Psi }_i^{}+d_i^{}(k_2,k_3)\mathrm{\Psi }_i^+\right]$$
(93)
then imposing the usual anticommutations rules
$$\{c_i(k_2,k_3),c_j^{}^{}(k_2^{},k_3^{})\}=\delta _{ij}\delta \left(^{}\right)\delta \left(k_2k_2^{}\right)\delta \left(k_3k_3^{}\right)\mathrm{}\mathrm{}$$
(94)
so defining a vacuum state $`0`$ by means of:
$$c_i(k_2,k_3)0=d_i(k_2,k_3)0=\mathrm{\hspace{0.17em}0}i,,k_2,k_3$$
(95)
It is easy now to show that this quantization in equivalent to the usual one and so that the state $`0`$ is the usual Minkowski vacuum state $`0_M`$. In fact we have the following relations:
$`c_i(k_2,k_3)=(\mathrm{\Psi }_i^{},\mathrm{\Psi })_M={\displaystyle _M}d^3x\mathrm{\Psi }_i^{}\mathrm{\Psi }=`$ (96)
$`=`$ $`{\displaystyle \underset{r}{}}{\displaystyle _0^+\mathrm{}}dk_1^{}{\displaystyle \frac{2\pi ^2}{\omega ^{}}}N^{}N_k^{}[e^{i\frac{\pi }{2}\left(i+\frac{1}{2}\right)}e^{\left(i+\frac{1}{2}\right)\vartheta ^{}}X_i^{}u_r(k^{})+`$
$`+`$ $`e^{i\frac{\pi }{2}\left(i\frac{1}{2}\right)}e^{\left(i\frac{1}{2}\right)\vartheta ^{}}Y_i^{}u_r(k^{})]\times a_r(k_1^{},k_2,k_3)=`$
$`=`$ $`{\displaystyle \underset{r}{}}{\displaystyle _0^+\mathrm{}}๐k_1^{}F_{ir}(k_1^{},)a_r(k_1^{},k_2,k_3)`$ (97)
with $`\vartheta ^{}=\frac{1}{2}\mathrm{ln}\left(\frac{\omega ^{}+k_1^{}}{\omega ^{}k_1^{}}\right)`$
and, for $`d_i^{}(k_2,k_3)`$:
$`d_i^{}(k_2,k_3)=(\mathrm{\Psi }_i^+,\mathrm{\Psi })={\displaystyle _M}d^3x\mathrm{\Psi }_i^+\mathrm{\Psi }=`$ (98)
$`=`$ $`{\displaystyle \underset{r}{}}{\displaystyle _0^+\mathrm{}}dk_1^{}{\displaystyle \frac{2\pi ^2}{\omega ^{}}}N^+N_k^{}[e^{i\frac{\pi }{2}\left(i+\frac{1}{2}\right)}e^{\left(i+\frac{1}{2}\right)\vartheta ^{}}X_i^{}v_r(k^{})+`$
$`+`$ $`e^{i\frac{\pi }{2}\left(i\frac{1}{2}\right)}e^{\left(i\frac{1}{2}\right)\vartheta ^{}}Y_i^{}v_r(k^{})]\times b^{}_r(k_1^{},k_2,k_3)=`$
$`=`$ $`{\displaystyle \underset{r}{}}{\displaystyle _0^+\mathrm{}}๐k_1^{}G_{ir}(k_1^{},)b_r^{}(k_1^{},k_2,k_3)`$ (99)
So it is demonstrated that the vacuum states defined by the two quantization procedures are the same. Moreover, by explicit calculation it is possible to show that
$`\{a_s(k_1^{\prime \prime },k_2,k_3),a_r^{}(k_1^{},k_2^{},k_3^{})\}=\delta _{rs}\delta \left(k_1^{}k_1^{\prime \prime }\right)\delta \left(k_2k_2^{}\right)\delta \left(k_3k_3^{}\right)`$
$``$ $`\{c_i(k_2,k_3),c_j^{}^{}(k_2^{},k_3^{})\}=Cost\times \delta _{ij}\delta \left(^{}\right)\delta \left(k_2k_2^{}\right)\delta \left(k_3k_3^{}\right)`$
so the two quantum constructions are totally equivalent.
### 4.4 The Unruh construction and the Unruh effect
We will now derive the Unruh effect for the spinor field following the standard procedure first used by Unruh himself.
Letโs first recall that an attempt in finding the relationship between the quantum construction in RS and that in MS was made by Fulling (for the scalar field), who simply identified the RS with the R-sector of the MS, consequently considered the Rindler vacuum state as a state in the Minkowski Hilbert space, and tried to express the Rindler annihilation (and creation) operators in terms of the usual plane waves ones. He then argued that Minkowski vacuum state could be considered as a particle state with respect to the Rindler vacuum state. Of course, because of the boundary condition we found necessary for the quantization in RS, this procedure is meaningless, since, as we already stressed, RS cannot be identified with R-sector of MS, but should be considered as a manifold on its own. Apart from this, however, there is another reason why the Fulling scheme is not valid. This scheme, in fact, implies to consider field modes in MS which corresponds to the Fulling ones in R and are zero everywhere else. This is equivalent to use a representation of the boost modes 89 given by
$`\mathrm{\Psi }_{}=\theta (x_+)\theta (x_{})\mathrm{\Psi }_{}^R+\theta (x_+)\theta (x_{})\mathrm{\Psi }_{}^F+\theta (x_+)\theta (x_{})\mathrm{\Psi }_{}^L+`$
$`+\theta (x_+)\theta (x_{})\mathrm{\Psi }_{}^P`$ (100)
with $`x_\pm =t\pm x`$, but then using for the quantization only the first term of this expression, so to be restricted in the R-wedge. This procedure cannot be valid, since physically these modes correspond to solutions of the field equation when infinite sources of energy are placed at the horizon, because of the presence of the theta function.
The procedure used by Unruh to compare the Minkowski quantization to the Rindler one is to construct a new quantization scheme, which should be valid in the whole MS, but should also reproduce the Fulling quantization in RS when restricted to to the R-sector. The idea is to build the Hilbert space of Minkowski states $`_M`$ out of the Hilbert space of solutions of the wave equation of the form $`_R_L`$, i.e. sum of solution of Minkowski equation of motion which are the same as the Fulling modes in the R (L) sector, but vanish identically in the L (R) sector.
Letโs perform the Unruh construction for quantization in MS, using our globally defined functions $`\mathrm{\Psi }_i^{}`$. Consider the functions:
$`R_i={\displaystyle \frac{1}{\sqrt{2\mathrm{cosh}\pi }}}\left(e^{\frac{\pi }{2}}\mathrm{\Psi }_i^{}+e^{\frac{\pi }{2}}\mathrm{\Psi }_i^+\right)`$ (101)
$`L_i`$ $`={\displaystyle \frac{1}{\sqrt{2\mathrm{cosh}\pi }}}\left(e^{\frac{\pi }{2}}\mathrm{\Psi }_i^{}e^{\frac{\pi }{2}}\mathrm{\Psi }_i^+\right)`$ (102)
which are solutions of Dirac equations, are eigenfunctions of $`M_{01}`$, are analytical in the whole Minkowski space, except fro the origin, and orthonormalzed in MS. Moreover, it happens that the $`R_i`$ are defined everywhere but in L sector and reduce themselves to the $`\mathrm{\Psi }_i^R`$ in the R one, while the $`L_i`$ manifest the inverse behaviour. Inverting these relations, and inserting into the expansion 93, we have:
$`\mathrm{\Psi }(t,x,y,z)={\displaystyle \underset{i=\mathrm{\hspace{0.17em}1},2}{}}{\displaystyle _{\mathrm{}}^+\mathrm{}}๐{\displaystyle ๐k_2๐k_3\left[c_{i,}(k_2,k_3)\mathrm{\Psi }_{i,}^{}+d_{i,}^{}(k_2,k_3)\mathrm{\Psi }_{i,}^+\right]}=`$ (103)
$`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle _{\mathrm{}}^+\mathrm{}}๐\left[r_{i,}R_{i,}+l_{i,}^{}L_{i,}\right]=`$
$`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle _0^+\mathrm{}}๐\left[r_{i,}R_{i,}+l_{i,}L_{i,}+r_{i,}^{}R_{i,}+l_{i,}^{}L_{i,}\right]`$ (104)
having introduced the coefficients:
$`r_{i,}={\displaystyle \frac{c_{i,}e^{\frac{\pi }{2}}+d_{i,}^{}e^{\frac{\pi }{2}}}{\sqrt{2\mathrm{cosh}\pi }}}`$ (105)
$`l_{i,}^{}={\displaystyle \frac{c_{i,}e^{\frac{\pi }{2}}d_{i,}^{}e^{\frac{\pi }{2}}}{\sqrt{2\mathrm{cosh}\pi }}}`$ (106)
These relations represent the Bogolubov transformation between the quantum constructions ( 93) and ( 104).
Now, consider carefully the nature of the operators $`r_{i,}`$ and $`l_{i,}`$ (and $`r_{i,}^{}`$ and $`l_{i,}^{}`$).
First of all suppose that the ( 104) represents an expansion of the spinor field, that leads to a correct quantization of it in MS, providing we impose the conditions
$$\{r_{i,},r_{j,^{}}^{}\}=\delta _{ij}\delta (^{})\{r_{i,},r_{j,}\}=\mathrm{\hspace{0.17em}0}\{r_{i,},l_{j,}\}=\mathrm{\hspace{0.17em}0}$$
(107)
and the analogous for $`l_{i,}`$ and $`l_{i,}^{}`$.
Suppose also that the operators above could be considered as annihilation (and creation) operators for R-particles and L-particles.
We stress that these hypothesis are crucial for the following derivation being physically meaningful.
In fact, if these both hold, the operators $`r_{i,}`$, defined as scalar products in MS, could be also expressed as integrals over the surface ($`t=0`$, $`x>0`$) which is a Cauchy (hyper)surface for the R wedge.
Moreover, this, together with the particular functional behaviour of the R-function (that reduce to the $`\mathrm{\Psi }_i^R`$ in the R sector), would mean that we can identify the operators $`r_{i,}`$ ($`r_{i,}^{}`$) with the Rindler annihilation (creation) operators $`a_{i,}`$ ($`a_{i,}^{}`$). Consequently, the R-particles would be identified with the Rindler particles, constructed in terms of the $`a_{i,}^{}`$, that an accelerated observer detects.
Letโs now consider the operator $`N_{i,j,,^{}}=r_i^{}r_j^{}`$ (we can of course not consider the quantum numbers $`k_2`$ and $`k_3`$, because they donโt play any significant role here). We remind that $``$ is the eigenvalue of $`M_{01}`$ but also the energy of a Rindler observer. It is easy to calculate the mean value of this operator in the Minkowski vacuum state, having:
$`{}_{M}{}^{}0N_{i,j,,^{}}0_{M}^{}=_M0r_j^{}^{}r_i0_M=`$ (108)
$`=`$ $`{\displaystyle \frac{1}{e^{2\pi }+\mathrm{\hspace{0.17em}1}}}\delta _{ij}\delta \left(^{}\right)\delta \left(k_2k_2^{}\right)\delta \left(k_3k_3^{}\right)`$ (109)
If we now identify the operators $`r_{i,}`$ with the operators $`a_{i,}`$, then we can intepret $`N_{i,j,,^{}}`$ as a Rindler particle number operator. This is exactly what is usually done in literature. The reasons for this identification were explained before and appear to be quite convincing. Nevertheless this passage is not trivial at all, as we will show.
Anyway, once identified the operators $`r_{i,}`$ with the operators $`a_{i,}`$, the result ( 109) can be interpreted as meaning that an inertial observer and a Rindler observer donโt share the same vacuum state, and that Minkowski vacuum state is a particle state for a Rindler observer.
We can also calculate the total number of particles that a Rindler observer will perceive if the field is in Minkowski vacuum state, for any quantum number and for unity of proper time. The result is:
$`dN={\displaystyle \underset{i,j}{}}{\displaystyle _0^+\mathrm{}}๐{\displaystyle _0^+\mathrm{}}๐^{}{\displaystyle ๐k_2๐k_2^{}๐k_3๐k_3^{}N_{i,j,,^{}}๐\eta }=`$
$`=`$ $`{\displaystyle \underset{i,j}{}}{\displaystyle _0^+\mathrm{}}d{\displaystyle _0^+\mathrm{}}d^{}{\displaystyle }dk_2{\displaystyle }dk_2^{}{\displaystyle }dk_3{\displaystyle }dk_3^{}{}_{M}{}^{}0r_j^{}^{}r_i0_{M}^{}d\eta =`$
$`=`$ $`{\displaystyle \underset{i,j}{}}{\displaystyle _0^+\mathrm{}}d{\displaystyle _0^+\mathrm{}}d^{}{\displaystyle }dk_2{\displaystyle }dk_2^{}{\displaystyle }dk_3{\displaystyle }dk_3^{}{\displaystyle \frac{1}{e^{2\pi }+\mathrm{\hspace{0.17em}1}}}\times `$ (111)
$`\times `$ $`\delta _{ij}\delta \left(^{}\right)\delta \left(k_2k_2^{}\right)\delta \left(k_3k_3^{}\right)d\eta =`$ (112)
$`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle _0^+\mathrm{}}๐{\displaystyle ๐k_2๐k_3\frac{1}{e^{2\pi }+\mathrm{\hspace{0.17em}1}}๐\eta }=`$ (113)
$`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle _0^+\mathrm{}}๐h_R{\displaystyle ๐k_2๐k_3\frac{1}{e^{\frac{2\pi h_R}{a}}+\mathrm{\hspace{0.17em}1}}๐\tau }=`$ (114)
$`=`$ $`2{\displaystyle _0^+\mathrm{}}๐h_R{\displaystyle ๐k_2๐k_3\frac{1}{e^{\frac{2\pi h_R}{a}}+\mathrm{\hspace{0.17em}1}}๐\tau }`$ (115)
where we have used the quantities: $`h_R=a`$, Rindler energy, and $`\tau =\frac{\eta }{a}`$, proper time, with $`a`$ acceleration of the Rindler observer. This result could be stated saying that the particle distribution of the Minkowski vacuum state, with respect to the quantization performed using the $`R_i`$ modes, is given by a thermal spectrum, according to Fermi-Dirac statistics with temperature:
$$T=\frac{1}{\beta k_B}=\frac{a}{2\pi k_B}$$
(116)
We saw that the identification between the operators $`r_{i,}`$ and $`a_{i,}`$ is crucial in its interpretation. We also pointed out that this identification is based on two strong hypotheses: that the Unruh construction is valid in the whole MS, and that the operators $`r_{i,}`$ ($`r_{i,}^{}`$), and $`l_{i,}`$ ($`l_{i,}^{}`$) can be considered as annhilation (creation) operators.
Consequently, we will now analyse in detail these crucial points.
### 4.5 Analysis of the Unruh construction
We are going now to show that the same arguments against the validity of the Unruh construction in the whole Minkowski Space, that were indicated in for the case of the scalar field, are preserved also for the spinor field and support the conclusion that the Unruh scheme is valid instead in the disjont union of the R and L wedge. This conclusion will be also proved by an explicit calculation.
First of all we argue that the Unruh quantization scheme is not suitable for the quantization of the spinor field in MS. This can be easily seen by looking at the initial expansion of the field ( 104). This is based on a separation of the integration interval into two parts corresponding to positive and negative values of the eigenvalue $``$, and this separation is necessary in order to have an expansion in terms of R and L modes. But we should recall the divergence of the functions $`\mathrm{\Psi }_{i,}`$ (hence of the Unruh R and L modes) at the origin of MS, as we already noticed. This, we said, doesnโt affect the quantization in MS in terms of the boost modes, but nevertheless implies that we cannot perform the separation of the integration interval as in the Unruh procedure. In fact, discarding the eigenfunction correspondent to the eigenvalue $`=0`$ (as any other eigenfunction, because the divergence at the origin is a common property of all of them) it would mean to discard an infinite number of degrees of freedom of the field, if the origin is in the domain of definition of our field, so commuting from our initial system to a phisically different one. Consequently we could say that the Unruh construction cannot be valid in the whole MS, which of course include the origin.
Moreover, the Unruh operators $`r`$ and $`l`$ (and their conjugates) cannot be considered as annihilation (and creation) operators for the field in MS, even if they satisfy the anticommutation relations ( 107). For this being possible, it is necessary the existence of a stationary ground state for the field in MS, which is defined with respect to a global timelike variable in the whole space, and which is annihilated by the $`r`$ and $`l`$ operators. But such a ground state is definitely missing in MS. In fact, there is no timelike variable with respect to which the Unruh modes (or their adjoints) are positive frequency solutions of the Dirac equation in MS (remember that the Killing vector $`\frac{}{\eta }`$ is spacelike in the $`F`$ and $`P`$ sectors) and consequently the Unruh operators will always be linear combinations of operators which create or annihilate particles with opposite frequency signs, i.e. there is not a stationary vacuum state in MS with respect to r-particles or l-particles.
These problems disappear if we consider the field only defined in the double Rindler wedge, i.e. in the disjoint union of the R and L sectors. Here the origin of MS is not considered, and so the divergence of the $`R`$ and $`L`$ modes has no physical consequencies for the separation of the integration domain in the expansion ( 104). In this manifold there exists a timelike killing vector with respect to which the $`R`$ and $`L`$ modes are positive frequency solutions of the Dirac equation, consequently there exists a stationary ground state for the system defined by the Unruh expansion, and the $`r`$ and $`l`$ operators could be interpreted as annihilation operators for the field in this double wedge. The result is that the Unruh construction is well defined in this case but, we emphasize, it refers to a field restricted to the union of the R and L sectors of MS and so physically different from the field defined in MS. In addition, since the R and L sectors are totally independent of each other, because they are separated by a spacelike interval, the quantization in these two wedges should be carried on separately. In other words we have the expansion:
$`\mathrm{\Psi }_{RL}(x)={\displaystyle \underset{i}{}}{\displaystyle _0^{\mathrm{}}}๐\left\{r_{i,}R_{i,}(x)+r_{i,}^{}R_{i,}(x)\right\}+`$
$`+{\displaystyle \underset{i}{}}{\displaystyle _0^{\mathrm{}}}๐\left\{l_{i,}L_{i,}(x)+l_{i,}^{}L_{i,}(x)\right\}xRL`$ (117)
We already showed that the quantization of the spinor field in the R sector requires the boundary condition ( 44), so we expect it to be manifest also for the field $`\mathrm{\Psi }_{RL}`$. In particular, if the Unruh espansion is valid (only) for the field in the double Rindler wedge, the Unruh operators $`r`$ should coincide with the operators $`a_{i,}`$ in terms of which the quantization of the field in RS is performed, provided that the field satisfy the boundary condition ( 44). Letโs now prove with an explicit calculation this statement.
First of all, we are going to find the explicit expression of the coefficients $`r_i`$ as functions of the values of the field $`\mathrm{\Psi }_{RL}`$. We recall that these are defined as:
$$r_{i,}=\frac{c_{i,}e^{\frac{\pi }{2}}+d_{i,}^{}e^{\frac{\pi }{2}}}{\sqrt{2\mathrm{cosh}\pi }}$$
(118)
so we first need to find the expression for the coefficients $`c_{i,}`$ and $`d_{i,}^{}`$, in terms of the field and of its spatial derivative. This task implies the proper treatment of integral whose hypersurface of integration is the $`t=0`$ hypersurface, which pass across the origin of Minkowski space, that is the intersection of the branch points where the functions we used are not well defined; this requires great attention.
By performing this calculation it is also possible to see that, as it should be, the coefficients $`c_{i,}`$ and $`d_{i,}^{}`$ are well defined without the need of any additional boundary condition other than the vanishing of the field at spatial infinity, in contrast to the โRindler coefficientsโ in ( 33).
Inserting the expressions so found for $`c_{i,}`$ and $`d_{i,}^{}`$ into the ( 118), we obtain:
$`r_{i,}=`$ (119)
$`=`$ $`{\displaystyle \frac{1}{2\pi \sqrt{\kappa }}}\sqrt{2\mathrm{cosh}\pi }{\displaystyle _{\mathrm{}}^+\mathrm{}}dye^{ik_2y}{\displaystyle _{\mathrm{}}^+\mathrm{}}dze^{ik_3z}\times `$
$`\times `$ $`\{X_i^{}{\displaystyle _0^+\mathrm{}}dx[K_{i+\frac{1}{2}}(\kappa x)\mathrm{\Psi }_{RL}(0,x)`$
$``$ $`{\displaystyle \frac{1}{2}}\mathrm{\Gamma }\left({\displaystyle \frac{1}{2}}+i\right)\left({\displaystyle \frac{\kappa x}{2}}\right)^{i\frac{1}{2}}\mathrm{\Psi }_{RL}(0,x)`$
$``$ $`{\displaystyle \frac{1}{\kappa }}{\displaystyle \frac{1}{\frac{1}{2}i}}\mathrm{\Gamma }({\displaystyle \frac{1}{2}}+i)\left({\displaystyle \frac{\kappa x}{2}}\right)^{i+\frac{1}{2}}{\displaystyle \frac{d}{dx}}\mathrm{\Psi }(0,x)]+`$
$`+`$ $`Y_i^{}{\displaystyle _0^+\mathrm{}}dx[K_{i\frac{1}{2}}(\kappa x)\mathrm{\Psi }_{RL}(0,x)`$
$``$ $`{\displaystyle \frac{1}{2}}\mathrm{\Gamma }\left({\displaystyle \frac{1}{2}}i\right)\left({\displaystyle \frac{\kappa x}{2}}\right)^{+i\frac{1}{2}}\mathrm{\Psi }_{RL}(0,x)`$
$``$ $`{\displaystyle \frac{1}{\kappa }}{\displaystyle \frac{1}{\frac{1}{2}+i}}\mathrm{\Gamma }({\displaystyle \frac{1}{2}}i)\left({\displaystyle \frac{\kappa x}{2}}\right)^{+i+\frac{1}{2}}{\displaystyle \frac{d}{dx}}\mathrm{\Psi }_{RL}(0,x)]\}`$
Changing coordinates to the Rindler ones (R sector), we have:
$`r_{i,}=`$ (120)
$`=`$ $`\sqrt{2}a_{i,}+{\displaystyle \frac{1}{2\pi \sqrt{\kappa }}}\sqrt{2\mathrm{cosh}\pi }{\displaystyle _{\mathrm{}}^+\mathrm{}}dye^{ik_2y}{\displaystyle _{\mathrm{}}^+\mathrm{}}dze^{ik_3z}\times `$
$`\times `$ $`\{\underset{\rho 0}{lim}{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{\frac{1}{2}i}}\mathrm{\Gamma }({\displaystyle \frac{1}{2}}+i)\left({\displaystyle \frac{\kappa \rho }{2}}\right)^{i+\frac{1}{2}}X_i^{}\mathrm{\Psi }^R(0,\rho )+`$
$`+`$ $`\underset{\rho 0}{lim}{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{\frac{1}{2}+i}}\mathrm{\Gamma }({\displaystyle \frac{1}{2}}i)\left({\displaystyle \frac{\kappa \rho }{2}}\right)^{+i+\frac{1}{2}}Y_i^{}\mathrm{\Psi }^R(0,\rho )\}`$
It is so clear that the coefficients $`r_{i,}`$ and $`a_{i,}`$ cannot be identified unless
$$\underset{\rho \mathrm{\hspace{0.17em}0}}{lim}\rho ^{\frac{1}{2}}\mathrm{\Psi }^R(0,\rho )=\mathrm{\hspace{0.17em}0}$$
(121)
which is just the boundary condition we found in sec. 3.
Letโs now discuss this result. We saw that it is not possible to identify the operators $`r_{i,}`$ and $`a_{i,}`$, unless ( 121); but there are no physical reasons to impose this boundary condition on the spinor field in MS; we have seen indeed that we obtain a meaningful quantum theory in terms of the Unruh modes only if we perform the โUnruh constructionโ just in the R and L sectors of MS, which are completely disjoint, from the causal and physical point of view; moreover, for an observer living in the Rindler wedge it is not possible, because of condition ( 121), to perform measurement in the whole MS, in order to put the field in the Minkowski vacuum state; in other words, the Unruh construction outlined in 4.4 does not represent a valid quantization scheme for the whole MS; consequently the operator $`r^{}r`$ cannot be interpreted as a particle number operator in MS and therefore, the relation ( 115) cannot be interpreted in any sense as a proof of the โUnruh effectโ.
## 5 Conclusions
Here we come to our conclusions. We can say that our analysis of the spinor field confirms the conclusion made in that the basic principles of quantum field theory imply that the Unruh procedure does not represent a correct derivation of the Unruh effect.
We saw that the reason for this conclusion is the existence of boundary conditions for the quantization of the spinor field in Rindler spacetime, preventing any relationship between this quantum construction and that one in Minkowski spacetime.The role played by this boundary conditions was analytically showed in our analysis of the Unruh procedure. We already noted that the existence of such boundary conditions could be expected since a Rinlder observer is confined inside the Rindler wedge by event horizons, so he would see these like spatial infinity of an inertial observer. For the same reason a Rindler observer has no relationship with MS and he cannot in any way prepare the quantum field in the Minkowski vacuum state.
We have also shown that the Unruh quantization scheme is valid only in the double Rindler wedge $`RL`$, and so it is also in the spinor case, because of the boundary condition ( 44); consequently the Unruh construction cannot be used to analyse the relationship between Rindler and Minkowski quantization schemes.
It seems to us that there are enough reasons to assert that this is a general feature of the analysed problem and that it holds true for any quantum field, and we are supported in this conclusion by the previously mentioned and similar results obtained for the scalar field and for the electromagnetic field.
On the other hand, the appearence of the fermion factor in the distribution ( 115) is entirely due to the particular form of the Bogolubov transformation ( 105), and there is no real need to interpretate it as an proof of a thermal nature of the spectrum; a similar situation is encounterd in other physical problems where the concept of temperature doesnโt arise at all, as it is explained in details in .
Of course, there are many different approaches to the Unruh problem and many derivations of the effect have been proposed. It is then worth to study these carefully in order to see if the difficulties found here for the Unruh procedure survive also in these cases, or they can be considered as correct derivations of the Unruh effect and of its consequences. Particularly interesting are the results in and , which deserve further study and attention.
As regards to the other aspect of the Unruh problem, that we mentioned in the beginning, namely the behaviour of an accelerated detector in Minkowski space, we could only say that it remains an open question. This should be clear also just considering that the Unruh effect is generally explained using the key role of the event horizons, but there are no horizon at all for a non-ideal accelerating detector, whose acceleration lasts a finite amount of time.
Our results show, however, that there are no reasons to expect that it will be of the type predicted by Unruh, at least as far as only the conventional derivation of it is considered, and that it should not be expected to be universal and independent from the nature and characteristics of the detector itself. It is worth to note that, as a partial confirmation of what we are saying, it was showed in great details in that elementary particles accelerated by a constant electric field donโt follow, in general, the Unruh behaviour, i.e. the thermal response with temperature (1).
We are sure, of course, that an accelerated detector behaves, in general, in a different way from an inertial one, and we admit that in some cases it could follow the Unruh behaviour, but we think that there are no quantum theoretical reasons to believe that this is the universal one.
## Acknowledgements
We would like to thank here V. A. Belinski, for enlightening discussions and advices, and for having read the manuscript of this paper, giving useful comments. A special thank goes also to B. Narozhny, who noticed an error occurred in a previous version of this work, now corrected, as well as U. Gerlach for useful comments.
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# The Theory of Fractal Time: Field Equations (the Theory of Almost Inertial Systems and Modified Lorentz Transformations)
## I Introduction
In the theory of multifractal time - for the case when the fractional dimensions of time $`d_t`$ almost coincide with integer value (equal unit) ($`d_t=1+\beta _iL_i(๐ซ(t),t)=1+\epsilon `$, $`|\epsilon |<<1`$) was shown (see ) that for modifying Lorentz transformations in the theory of fractal time it is necessary to change factor $`\beta =\sqrt{1v^2/c^2}`$ by factor $`\beta ^{}=\sqrt[4]{\beta ^4+4a_0^2}`$. Then relativistic energy has form
$$E=mc^2$$
(1)
$`m=\beta ^1m_0,\beta ^2=1{\displaystyle \frac{v^2}{c^2}},vc`$ (2)
$`m=\beta ^1m_0\sqrt{1+\beta ^2+\beta ^2},\beta ^2={\displaystyle \frac{v^2}{c^2}}1,vc`$ (3)
$$p=\beta ^1m_0v$$
(4)
$$\beta ^{}=\sqrt[4]{\beta ^4+4a_0^2(t)}$$
(5)
$$a_0=\underset{i}{}\beta _iF_{0,i}\frac{v}{c}ct,\underset{i}{}F_{0,i}=\underset{i}{}\frac{}{๐ซ}L_i$$
(6)
On the base of fractal theory in this paper we presented the new equations. These equations are analogies of relativistic equations (scalar,vector, tensor, spinor) and valid in the domain of arbitrary velocities including speed of light if external physical fields donโt equal zero (but so small that fractal dimensions of time is near unit). These equations in the case of absence of physical fields coincide with the usual relativistic equations, but give the new fields with imaginary energies (including the case of rest energies). If the external fields are not equal zero, the received equations give also a new spin characteristics.
## II Equations for Fields on the Base of Modified Relations for Energy
Let us write formulas (1)-(2) in the form (using the hypothesis about approximate conservation of the energy-momentum vector in the space with fractal time )
$$E^2=\frac{E_0^2}{\beta ^2}=\frac{๐ฉ_0^2c^2}{\beta ^2}+E_0^2$$
(7)
We introduce now the new designation for relativistic energy and momentum similar to those that have been used in SR ($`\beta ^2=E^2E_0^2,E=m_0c^2\beta ^2)`$ (see (2)). Then equation without roots has form
$$\frac{(E^2๐ฉ^2c^2)^2}{(1+4a_0^2E^4E_0^4)}=E_0^4$$
(8)
$$E=m_0\beta ^1,p=m_0\text{v}\beta ^1$$
(9)
The equation (8) is the base equation for describing the energy in the space with multifractal time. For $`E`$ we receive in the absence of fields and the momentum equal zero the four solutions:$`E=\pm E_0`$, $`E=\pm iE_0`$. Thus, there are particles and anti-particles with real ($`E=\pm E_0`$) and imaginary ($`E=\pm iE_0`$) masses. The last are new sort of particles with imaginary masses. These particles are not taxions because they exist in the domain of velocities $`v0`$ (including velocity $`v=c`$). For receiving equations for fields and particles we use the ordinary method replacing the energy and the momentum by derivatives. This method consists in changing the energy $`E=E_0/(1v^2/c^2)^{1/2}`$ and the momentum $`P=m_0v/(1v^2/c^2)^{1/2}`$ by derivatives with respect to time and space coordinates that used in quantum mechanics: $`\mathrm{}=c=1`$, $`\frac{E}{\beta }`$ $`i\frac{}{t}=\widehat{E}`$, $`\frac{P_j}{\beta }i_j=\widehat{๐ฉ}`$, (j = 1,2,3). So we obtain ( in case $`vc`$) for function $`\mathrm{\Phi }(๐ซ,t)`$ the integral -differential equation
$$\frac{(\widehat{E}^2\widehat{๐ฉ}^2)^2}{(1+4a_0^2\widehat{E}^4E_0^4)}\mathrm{\Phi }(๐ซ,t)=E_0^4\mathrm{\Phi }(๐ซ,t)$$
(10)
where $`\widehat{E}`$ and $`\widehat{p}`$ are differential operators and was determined early. For simplifying this equation we multiply it by operator $`(1+4a_0^2)\frac{\widehat{E}^4}{E_0^4})`$ (thou it may introduce non-physical solutions but it is makes the equation (10)only differential). So we obtain
$`(\mathrm{}^24a_0^2{\displaystyle \frac{^4}{t^4}})\mathrm{\Phi }(๐ซ,t)=E_0^4\mathrm{\Phi }(๐ซ,t)`$ (11)
where $`\mathrm{}`$ is DโAlamber operator ($`\mathrm{}=\mathrm{\Delta }\frac{^2}{t^2})`$, $`\mathrm{\Delta }`$ is Laplasian), $`\mathrm{\Phi }`$ are functions describing particles or fields. For scalar $`\mathrm{\Phi }`$ equation (10) describes the scalar field in the space with fractal dimensions that originated by the presence of the external physical fields ($`a_00`$ . The corrections in (11) to the usual DโAlamber equation are the result of modifying the Lorentz transformation. The last is consequences of fractal nature of time. For receiving the equations in which taken into account the influence of multifractal structure of time on using the derivatives and receive more correct equation it is necessary to use the generalized Riemann-Liouville fractional derivatives (GFD),. In that case equation (11)take the form
$`(D_{,t}^{d_t}D_{+,t}^{d_t}`$ $``$ $`\mathrm{\Delta })^2\mathrm{\Phi }(๐ซ,t)=[E_0^4+`$ (12)
$`+`$ $`4a_0^2(D_{,t}^{d_t}D_{+,t}^{d_t})^2]\mathrm{\Phi }(๐ซ,t)`$ (13)
where
$$D_{+,t}^df(t)=\left(\frac{d}{dt}\right)^n_a^t\frac{f(t^{})dt^{}}{\mathrm{\Gamma }(nd(t^{}))(tt^{})^{d(t^{})n+1}}$$
(14)
$$D_{,t}^df(t)=\left(\frac{d}{dt}\right)^n_t^b\frac{(1)^nf(t^{})dt^{}}{\mathrm{\Gamma }(nd(t^{}))(t^{}t)^{d(t^{})n+1}}$$
(15)
where $`\mathrm{\Gamma }(x)`$ is Eulerโs gamma function, and $`a`$ and $`b`$ are some constants from $`[0,\mathrm{})`$. In these definitions, as usually, $`n=\{d\}+1`$ , where $`\{d\}`$ is the integer part of $`d`$ if $`d0`$ (i.e. $`n1d<n`$) and $`n=0`$ for $`d<0`$. In this paper we donโt consider equations with fractal derivatives and restrict consideration only by calculation of corrections from alterations of Lorentz transformation
## III equations of four and second order in derivatives
It is useful rewrite (11) in the form
$$\mathrm{}^2\mathrm{\Phi }(๐ซ,t)=(E_0^4+4a_0^2\frac{^4}{t^4})\mathrm{\Phi }(๐ซ,t)$$
(16)
Now introduce the four component unit matrix $`I`$ and the Dirac type matrices $`\alpha _i`$: ($`\alpha _i^2=1`$, $`\alpha _i\alpha _j+\alpha _j\alpha _i=0`$, $`ij`$; $`i,j=1,2,3,4`$ ). Than after usual splitting equation procedure to the equations (16) we receive
$$[\mathrm{}I+2a_0\frac{^2}{t^2}\alpha _2]\mathrm{\Phi }=\alpha _1E_0^2\mathrm{\Phi }$$
(17)
where $`\mathrm{\Phi }`$ is a four element bispinor column
$$\mathrm{\Phi }=\left(\begin{array}{c}\mathrm{\Phi }_1\\ \mathrm{\Phi }_2\\ \mathrm{\Phi }_3\\ \mathrm{\Phi }_4\end{array}\right)$$
(18)
So, we have four equations for $`\mathrm{\Phi }_1,\mathrm{}\mathrm{\Phi }_4`$ (18).
$$\mathrm{}\mathrm{\Phi }_12a_0\frac{^2}{t^2}\mathrm{\Phi }_4=E_0^2\mathrm{\Phi }_1$$
(19)
$$\mathrm{}\mathrm{\Phi }_22a_0\frac{^2}{t^2}\mathrm{\Phi }_3=E_0^2\mathrm{\Phi }_2$$
(20)
$$\mathrm{}\mathrm{\Phi }_32a_0\frac{^2}{t^2}\mathrm{\Phi }_2=E_0^2\mathrm{\Phi }_3$$
(21)
$$\mathrm{}\mathrm{\Phi }_42a_0\frac{^2}{t^2}\mathrm{\Phi }_1=E_0^2\mathrm{\Phi }_4$$
(22)
In the equations (19)-(22) the first two equations describe the particles or fields with real energies, the last two equations describe the new particles or fields with imaginary energies.The energies of particles with real energies depends on behavior of fields with imaginary energies and vice versa. The new spin characteristics consequences of these equations.
## IV equations of first order in derivatives
For receiving of first order equations in derivatives it is necessary to introduce the Dirac matrices $`\gamma _i`$ $`(i=0,1,2,3)`$ and Dirac type matrices $`\sigma _j`$ ($`j=0,1,2,3`$). These matrices may be used for splitting of the left-hand side and the right-hand side of each of the equations (19)-(22). We write the function $`a_0`$ in the form
$`a_0=a_g+a_e+a_n`$ $`=`$ $`\beta _gL_g(๐ซ,t)+\beta _eL_e(๐ซ,t)+`$ (23)
$`+\beta _nL_n(๐ซ,t)`$ (24)
where $`L_g`$, $`L_e`$, $`L_n`$ are Lagrangians density of energies for gravitational, electro-weak and strong fields. Let a module of complex function $`\mathrm{\Phi }_i`$ is the probability to find the particle in a moment $`t`$ in a point $`๐ซ`$. In that case the square root $`\sqrt{\mathrm{\Phi }_i}=\psi _i(๐ซ,t)`$ gives the function with characteristics: $`\psi ^{}\psi `$ has sense of probability. Now if change in equations (19)-(22) the differential operators by $`E`$ and $`P`$ , translate in the right-hand side of equations all members with $`a_0`$ and than extract a square root from left-hand side and right-hand side of these equations we obtain (after splitting we change $`E`$ and $`p`$ by ordinary differential operators and use the designations $`\psi _j(i)`$ where indexes $`i`$ corresponds to the indexes $`i`$ at $`\mathrm{\Phi }_i`$ and indexes $`j`$ corresponds to each splitting component of $`\mathrm{\Phi }_i`$
$`i\gamma _i_i\sqrt{\mathrm{\Phi }_1}`$ $`=`$ $`E_0\sigma _1\sqrt{\mathrm{\Phi }_1}+\sqrt{2a_g}\sigma _2i_4\sqrt{\mathrm{\Phi }_4}`$ (25)
$`+`$ $`\sqrt{2a_e}\sigma _3i_4\sqrt{\mathrm{\Phi }_4}+\sqrt{2a_n}\sigma _4i_4\sqrt{\mathrm{\Phi }_4}`$ (26)
$`i\gamma _i_i\sqrt{\mathrm{\Phi }_2}`$ $`=`$ $`E_0\sigma _1\sqrt{\mathrm{\Phi }_2}+\sqrt{2a_g}\sigma _2i_4\sqrt{\mathrm{\Phi }_3}`$ (27)
$`+`$ $`\sqrt{2a_e}\sigma _3i_4\sqrt{\mathrm{\Phi }_3}+\sqrt{2a_n}\sigma _4i_4\sqrt{\mathrm{\Phi }_3}`$ (28)
$`i\gamma _i_i\sqrt{\mathrm{\Phi }_3}`$ $`=`$ $`E_0\sigma _1\sqrt{\mathrm{\Phi }_3}+\sqrt{2a_g}\sigma _2i_4\sqrt{\mathrm{\Phi }_2}`$ (29)
$`+`$ $`\sqrt{2a_e}\sigma _3i_4\sqrt{\mathrm{\Phi }_2}+\sqrt{2a_n}\sigma _4i_j\sqrt{\mathrm{\Phi }_2}`$ (30)
$`i\gamma _i_i\sqrt{\mathrm{\Phi }_4}`$ $`=`$ $`E_0\sigma _1\sqrt{\mathrm{\Phi }_4}+\sqrt{2a_g}\sigma _2i_4\sqrt{\mathrm{\Phi }_1}`$ (31)
$`+`$ $`\sqrt{2a_e}\sigma _3i_4\sqrt{\mathrm{\Phi }_1}+\sqrt{2a_n}\sigma _4i_j\sqrt{\mathrm{\Phi }_1}`$ (32)
$$\sqrt{\mathrm{\Phi }_j}=\left(\begin{array}{c}\psi _1(j)\\ \psi _2(j)\\ \psi _3(j)\\ \psi _4(j)\end{array}\right),(j=1,2,3,4)$$
(33)
The equations (25)- (31) are generalized Dirac equations based on taking into account only the modified by fractal nature of time Lorentz transformations ( and possibility of moving with speed of light as the consequences of it). The eight equations (25) and (27) describe the particles (or fields) with spin $`\frac{\mathrm{}}{2}`$ ($`\mathrm{\Phi }_i`$ are bispinors), real energy and new characteristics โquasi-spinโ originated by influence of the fields with imaginary energy. The last described by the equations (29) and (31). Thus there are two sorts of particles described by equations (25)-(31): with real energies and with imaginary energies. Each sort of particles have own anti-particles and โquasi-spinโ. In these equations taken into account the influences of all known fields (gravitational, electro-weak, strong).
## V Electromagnetic fields with a rest mass (equations Proca)
Let us consider the case when each of functions $`\mathrm{\Phi }_i`$ is four- vector of electro-magnetic field ($`\mathrm{\Phi }_i(0),\mathrm{\Phi }_i(1),\mathrm{\Phi }_i(2),\mathrm{\Phi }_i(3)`$) with a rest energy $`E_0`$. The equations (16)take form
$$\mathrm{}\mathrm{\Phi }_\mu (i)=[\alpha _1E_0^2+2a_0\frac{^2}{t^2}\alpha _2]\mathrm{\Phi }_\mu (i),(\mu =0,1,2,3)$$
(34)
So there are four sorts of electro-magnetic Proca fields: fields of Proca photons and Proca anti-photons with real and imaginary energies and different โquasi-spinsโ . For $`a_0=0)`$ all the equations coincide with Proca equations.
## VI Maxwell equations
If the Maxwell equations of electro-magnetic field may be considered as the equations of the Proca in the limits of a rest mass equal zero, in that case the new sorts of electro-magnetic fields (anti-photons fields and fields with imaginary energies) as consequences of the Proca fields appear. Let us write equations for 4-vector electro-magnetic potentials $`A_\mu (i)`$ ( which are consequences of equations (11) and (19)-(22)) for electro-magnetic fields where role $`E_0^2`$ plays 4-vector of electric current $`j_\mu `$, $`(\mu =0,1,2,3)`$
$$\mathrm{}A_\mu (i)=[\alpha _1j_\mu (i)+2a_0\frac{^2}{t^2}\alpha _2]A_\mu (i),\mu =0,1,2,3$$
(35)
In equation (35) $`A_\mu (i)`$ is a 4-column with respect to $`i`$. The equations (35) are generalized Maxwell equations for photon and anti-photon fields, with usual and โquasiโ spins and coincide for $`a_0=0`$ with usual Maxwell equations. The matrices $`\alpha _1`$ and $`\alpha _2`$ may be chosen for example as
$$\alpha _1=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right),\alpha _2=\left(\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& 1& 0\\ 0& 1& 0& 0\\ 1& 0& 0& 0\end{array}\right)$$
(36)
The last eight equations (35) describe anti-photon fields and its energy (if electrical charges for anti-photon fields are real) may be imaginary. In these fields the role of minus-sign charges will play plus-sign charges and vice versa. There are difference between speed of photons light and anti-photons light. The speeds of photons light and anti-photons light are equal only if $`a_0=0`$.
## VII Conclusion
There are new consequences of equations based on the modified Lorentz transformations in the multifractal time theory and now we stress it: $`a`$) the existence of two imaginary solutions for the rest energy ($`iE_0`$ and $`iE_0`$) and thus the existence of new class of particles with imaginary mass. Because its velocities are arbitrary ( $`0v\mathrm{}`$ they are not taxions ); $`b`$) the appearance of new โquasi-spinโ characteristics. The cause of it lays in the more high order ( four order) differential equations then usual equations and in additional splitting of square roots and taking into account the existence of external fields ($`a_0`$). The physical sense and physical nature of new spin characteristics are not clear. The nature of the additional โspinsโ originated by additional decompositions of square roots for equations of four order and physical sense of it needs in special investigation; ; $`c`$) all equations coincide with known physical equation if fractal dimension of time is integer ( so in that sense the theory is not contradicts known physical theories). In the fractal theory of time the last case corresponds to vanishing of physical fields (the last originates the fractional dimensions of time); $`d`$) the theory use the improved classical relations for relativistic energy which take into attention the fractional dimensions of time and allow motions with arbitrary velocities (including velocities equal the speed of light); e) it may be shown that the equations for case $`vc`$ some differs at equations used in the paper but main results the theory based on them coincide qualitative with the results used in this paper(if use the equations of this paper for case $`vc`$ ).
The presented in this paper the theory based on a relative motions in almost inertial systems which in turn based on the multifractal time theory and gives the new describing for characteristics of moving bodies (energy, momentum, mass and so on). The main results of this theory used in this paper are: $`a`$) the possibility of moving with arbitrary velocities without appearance of infinitum energy and imaginary mass; $`b`$) existence of maximum energy if $`v=c`$; $`c`$) possibility of experimental verification the main results of the theory.
The theory - describes the Universe as an open systems (the theory of open systems see in ). This theory coincides with SR after transition to inertial systems (if neglect by the fractional dimensions of time) or almost coincides (the differences are non-essential) for velocities $`v<c`$. The movement of bodies with velocities that exceed the speed of light is accompanied by a series of physical effectโs which can be found by experiments (these effects was considered in the separate papers (, , ) in more details). It is shown in these papers the necessity to receive the particles with energies $`E_0`$ $`10^3`$ $`1/\sqrt{t}`$ and $`\frac{}{x}a_e<\frac{}{x}a_g`$ for verification of the theory.If accelerate the particles by electric fields then $`E_0[2\stackrel{~}{E}(Mc_2)^1t]^{\frac{1}{2}}<E_0(\sqrt{2a_gt})^1)`$. In this formula $`\stackrel{~}{E}`$ is the electric field strength, $`M`$ is the mass of electric charges originated the $`\stackrel{~}{E}`$ . It is useful to pay attention to the problem of receiving the particles with such energies of the physicians of known accelerate centers. If the possibilities for organizing such works will be found, the results are useful for development of our views on the nature of an energy, the time and the space.
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# Absence of Static, Spherically Symmetric Black Hole Solutions for Einstein-Dirac-Yang/Mills Equations with Complete Fermion Shells
## 1 Introduction
Recently the Einstein-Dirac-Yang/Mills (EDYM) equations were studied for a static, spherically symmetric system of a Dirac particle interacting with both a gravitational field and an $`SU(2)`$ Yang-Mills field . In these papers, the Dirac particle had no angular momentum, and we could make a consistent ansatz for the Dirac wave function involving two real spinor functions. In the present paper, we allow the Dirac particles to have non-zero angular momentum $`j`$, $`j=1,2,\mathrm{}`$. Similar to , we can build up a spherically symmetric system out of $`(2j+1)`$ such Dirac particles. In this case however, a reduction to real 2-spinors is no longer possible, but we can obtain a consistent ansatz involving four real spinor functions.
We show that the only black hole solutions of our 4-spinor EDYM equations are those for which the spinors vanish identically outside the black hole; thus these EDYM equations admit only the Bartnik-McKinnon (BM) black hole solutions of the $`SU(2)`$ Einstein-Yang/Mills equations . This result extends our work in to the case with angular momentum; it again means physically that the Dirac particles must either enter the black hole or escape to infinity. This generalization comes as a surprise because if one thinks of the classical limit, then classical point particles with angular momentum can โrotate aroundโ the black hole on a stable orbit. Our result thus shows that the non-existence of black hole solutions is actually a quantum mechanical effect. A simple way of understanding the difference between classical and quantum mechanical particles is that for classical particles, the centrifugal barrier prevents the particles from falling into the black hole, whereas quantum mechanical particles can tunnel through this barrier. In our system, tunnelling alone does not explain the non-existence of black-hole solutions, because the Dirac particles are coupled to the classical fields; that is, they can influence the potential barrier. Our results are established by analyzing in detail the interaction between the matter fields and the gravitational field.
In Section 2, we derive the static, spherically symmetric $`SU(2)`$ EDYM equations with non-zero angular momentum of the Dirac particles. By assuming the BM ansatz for the YM potential (the vanishing of the electric component), the resulting system consists of 4 first-order equations for the spinors, two first-order Einstein equations, and a second-order equation for the YM potential. This EDYM system is much more complicated than the system considered in , and in order to make possible a rigorous mathematical analysis of the equations, we often assume (as in ) a power ansatz for the metric functions and the YM potential. Our analysis combines both geometrical and analytic techniques.
## 2 Derivation of the EDYM Equations
We begin with the separation of variables for the Dirac equation in a static, spherically symmetric EYM background. As in , we choose the line element and the YM potential $`๐`$ in the form
$`ds^2`$ $`=`$ $`{\displaystyle \frac{1}{T(r)^2}}dt^2{\displaystyle \frac{1}{A(r)}}dr^2r^2d\vartheta ^2r^2\mathrm{sin}^2\vartheta d\phi ^2`$ (1)
$`๐`$ $`=`$ $`w(r)\tau ^1d\vartheta +(\mathrm{cos}\vartheta \tau ^3+w(r)\mathrm{sin}\vartheta \tau ^2)d\phi `$ (2)
with two metric functions $`A`$, $`T`$, and the YM potential $`w`$. The Dirac operator was computed in \[1, Section 2\] to be
$`G`$ $`=`$ $`iT\gamma ^t_t+\gamma ^r\left(i\sqrt{A}_r+{\displaystyle \frac{i}{r}}(\sqrt{A}1){\displaystyle \frac{i}{2}}\sqrt{A}{\displaystyle \frac{T^{}}{T}}\right)`$ (3)
$`+i\gamma ^\vartheta _\vartheta +i\gamma ^\phi _\phi +{\displaystyle \frac{2i}{r}}(w1)(\stackrel{}{\gamma }\stackrel{}{\tau }\gamma ^r\tau ^r)\tau ^r.`$
This Dirac operator acts on 8-component wave functions, which as in we denote by $`\mathrm{\Psi }=(\mathrm{\Psi }^{\alpha ua})_{\alpha ,u,a=1,2}`$, where $`\alpha `$ are the two spin orientations, $`u`$ corresponds to the upper and lower components of the Dirac spinor (usually called the โlargeโ and โsmallโ components, respectively), and $`a`$ is the YM index. The Dirac equation is
$$(Gm)\mathrm{\Psi }=\mathrm{\hspace{0.33em}0},$$
(4)
where $`m`$ is the rest mass of the Dirac particle, which we assume to be positive ($`m>0`$).
As explained in , the Dirac operator (3) commutes with the โtotal angular momentum operatorsโ
$$\stackrel{}{J}=\stackrel{}{L}+\stackrel{}{S}+\stackrel{}{\tau },$$
(5)
where $`\stackrel{}{L}`$ is angular momentum, $`\stackrel{}{S}`$ the spin operator, and $`\stackrel{}{\tau }`$ the standard basis of $`\text{su}(2)_{\text{YM}}`$. Thus the Dirac operator is invariant on the eigenspaces of total angular momentum, and we can separate out the angular dependence by restricting the Dirac operator to suitable eigenspaces of the operators $`\stackrel{}{J}`$. Since (5) can be regarded as the addition of angular momentum and two spins $`\frac{1}{2}`$, the eigenvalues of $`\stackrel{}{J}`$ are integers. In , the Dirac equation was considered on the kernel of the operator $`J^2`$; this leads to the two-component Dirac equation \[1, (2.23),(2.24)\]. Here we want to study the effect of angular momentum and shall thus concentrate on the eigenspaces of $`J^2`$ with eigenvalues $`j(j+1)`$, $`j=1,2,\mathrm{}`$. Since the eigenvalues of $`J_z`$ merely describe the orientation of the wave function in space, it is furthermore sufficient to restrict attention to the eigenspace of $`J_z`$ corresponding to the highest possible eigenvalue. Thus we shall consider the Dirac equation on the wave functions $`\mathrm{\Psi }`$ with
$$J^2\mathrm{\Psi }=j(j+1)\mathrm{\Psi }\text{and}J_z\mathrm{\Psi }=j\mathrm{\Psi }(j=1,2,\mathrm{}).$$
(6)
Since (6) involves only angular operators, it is convenient to analyze these equations on spinors $`\mathrm{\Phi }^{\alpha a}(\vartheta ,\phi )`$ on $`S^2`$ (where $`\alpha `$ and $`a`$ are again the spin and YM indices, respectively). Let us first determine the dimension of the space spanned by the vectors satisfying (6). Using the well-known decomposition of two spins $`\frac{1}{2}`$ into a singlet and a triplet, we choose a spinor basis $`\mathrm{\Phi }_{st}`$ with $`s=0,1`$ and $`sts`$ satisfying
$$(\stackrel{}{S}+\stackrel{}{\tau })^2\mathrm{\Phi }_{st}=s(s+1)\mathrm{\Phi }_{st},(S_z+\tau _z)\mathrm{\Phi }_{st}=t\mathrm{\Phi }_{st}.$$
The spherical harmonics $`(Y_{lk})_{l0,lkl}`$, on the other hand, are a basis of $`L^2(S^2)`$. Using the rules for the addition of angular momentum, the wave functions satisfying (6) must be linear combinations of the following vectors,
$`Y_{jj}\mathrm{\Phi }_{\mathrm{0\hspace{0.25em}0}}`$ (7)
$`Y_{j1j1}\mathrm{\Phi }_{\mathrm{1\hspace{0.25em}1}}`$ (8)
$`Y_{jj1}\mathrm{\Phi }_{\mathrm{1\hspace{0.25em}1}},Y_{jj}\mathrm{\Phi }_{10}`$ (9)
$`Y_{j+1j1}\mathrm{\Phi }_{11},Y_{j+1j}\mathrm{\Phi }_{10},Y_{j+1j+1}\mathrm{\Phi }_{11}.`$ (10)
These vectors all satisfy the second equation in (6), but they are not necessarily eigenfunctions of $`J^2`$. We now use the fact that a vector $`\mathrm{\Psi }0`$ satisfying the equation $`J_z\mathrm{\Psi }=j\mathrm{\Psi }`$ is an eigenstate of $`J^2`$ with eigenvalue $`j(j+1)`$ if and only if it is in the kernel of the operator $`J^+=J_x+iJ_y`$. Thus the dimension of the eigenspace (6) coincides with the dimension of the kernel of $`J^+`$, restricted to the space spanned by the vectors (7)-(10). A simple calculation shows that this dimension is four (for example, we have $`J^+(Y_{jj1}\mathrm{\Phi }_{\mathrm{1\hspace{0.25em}1}})=Y_{jj}\mathrm{\Phi }_{\mathrm{1\hspace{0.25em}1}}=J^+(Y_{jj}\mathrm{\Phi }_{\mathrm{1\hspace{0.25em}0}})`$, and thus $`J^+`$ applied to the vectors (9) has a one-dimensional kernel).
We next construct a convenient basis for the angular functions satisfying (6). We denote the vector (7) by $`\mathrm{\Phi }_0`$. It is uniquely characterized by the conditions
$`L^2\mathrm{\Phi }_0`$ $`=`$ $`j(j+1)\mathrm{\Phi }_0,L_z\mathrm{\Phi }_0=j\mathrm{\Phi }_0`$
$`(\stackrel{}{S}+\stackrel{}{\tau })\mathrm{\Phi }_0`$ $`=`$ $`0,\mathrm{\Phi }_0_{S^2}=\mathrm{\hspace{0.33em}1}.`$
We form the remaining three basis vectors by multiplying $`\mathrm{\Phi }_0`$ with spherically symmetric combinations of the spin and angular momentum operators, namely
$`\mathrm{\Phi }_1`$ $`=`$ $`2S^r\mathrm{\Phi }_0=2\tau ^r\mathrm{\Phi }_0`$
$`\mathrm{\Phi }_2`$ $`=`$ $`{\displaystyle \frac{2}{c}}(\stackrel{}{S}\stackrel{}{L})\mathrm{\Phi }_0={\displaystyle \frac{2}{c}}(\stackrel{}{\tau }\stackrel{}{L})\mathrm{\Phi }_0`$
$`\mathrm{\Phi }_3`$ $`=`$ $`{\displaystyle \frac{4}{c}}S^r(\stackrel{}{S}\stackrel{}{L})\mathrm{\Phi }_0={\displaystyle \frac{4}{c}}\tau ^r(\stackrel{}{\tau }\stackrel{}{L})\mathrm{\Phi }_0`$
where
$$c=\sqrt{j(j+1)}\mathrm{\hspace{0.33em}0}.$$
Since the operators $`S^r`$, $`\tau ^r`$, and $`(\stackrel{}{S}\stackrel{}{L})`$ commute with $`\stackrel{}{J}`$, it is clear that the vectors $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_3`$ satisfy (6). Furthermore, using the standard commutation relations between the operators $`\stackrel{}{L}`$, $`\stackrel{}{x}`$, and $`\stackrel{}{S}`$ , we obtain the relations
$`(\stackrel{}{S}\stackrel{}{L})^2`$ $`=`$ $`{\displaystyle \frac{1}{4}}L^2+{\displaystyle \frac{i}{2}}ฯต_{jkl}S_lL_jL_k={\displaystyle \frac{1}{4}}L^2{\displaystyle \frac{1}{4}}ฯต_{jkl}S_lฯต_{jkm}L_m`$
$`=`$ $`{\displaystyle \frac{1}{4}}L^2{\displaystyle \frac{1}{2}}\stackrel{}{S}\stackrel{}{L}`$
$`(\stackrel{}{S}\stackrel{}{\tau })^2`$ $`=`$ $`{\displaystyle \frac{1}{4}}\tau ^2{\displaystyle \frac{1}{2}}\stackrel{}{S}\stackrel{}{\tau }={\displaystyle \frac{3}{16}}{\displaystyle \frac{1}{2}}\stackrel{}{S}\stackrel{}{\tau }`$
$`2S^r\mathrm{\Phi }_0`$ $`=`$ $`2\tau ^r\mathrm{\Phi }_0=\mathrm{\Phi }_1`$
$`2S^r\mathrm{\Phi }_1`$ $`=`$ $`2\tau ^r\mathrm{\Phi }_1=\mathrm{\Phi }_0`$
$`2S^r\mathrm{\Phi }_2`$ $`=`$ $`2\tau ^r\mathrm{\Phi }_2=\mathrm{\Phi }_3`$
$`2S^r\mathrm{\Phi }_3`$ $`=`$ $`2\tau ^r\mathrm{\Phi }_3=\mathrm{\Phi }_2`$
$`(\stackrel{}{S}\stackrel{}{\tau })\mathrm{\Phi }_0`$ $`=`$ $`S^2\mathrm{\Phi }_0={\displaystyle \frac{3}{4}}\mathrm{\Phi }_0`$
$`(\stackrel{}{S}\stackrel{}{\tau })\mathrm{\Phi }_1`$ $`=`$ $`2S_k\tau _kS^r\mathrm{\Phi }_0=2S_k(\stackrel{}{x}\stackrel{}{S})S_k\mathrm{\Phi }_0`$
$`=`$ $`2S^2S^r\mathrm{\Phi }_0S_kx_k\mathrm{\Phi }_0={\displaystyle \frac{1}{2}}S^r\mathrm{\Phi }_0={\displaystyle \frac{1}{4}}\mathrm{\Phi }_1`$
$`(\stackrel{}{S}\stackrel{}{\tau })\mathrm{\Phi }_2`$ $`=`$ $`{\displaystyle \frac{2}{c}}(\stackrel{}{S}\stackrel{}{\tau })(\stackrel{}{S}\stackrel{}{L})\mathrm{\Phi }_0={\displaystyle \frac{1}{c}}(\stackrel{}{L}\stackrel{}{\tau })\mathrm{\Phi }_0{\displaystyle \frac{2}{c}}(\stackrel{}{S}\stackrel{}{L})(\stackrel{}{S}\stackrel{}{\tau })\mathrm{\Phi }_0`$
$`=`$ $`{\displaystyle \frac{1}{c}}(\stackrel{}{S}\stackrel{}{L})\mathrm{\Phi }_0+{\displaystyle \frac{3}{2c}}(\stackrel{}{S}\stackrel{}{L})\mathrm{\Phi }_0={\displaystyle \frac{1}{4}}\mathrm{\Phi }_2`$
$`(\stackrel{}{S}\stackrel{}{\tau })\mathrm{\Phi }_3`$ $`=`$ $`{\displaystyle \frac{4}{c}}(\stackrel{}{S}\stackrel{}{\tau })S^r(\stackrel{}{S}\stackrel{}{\tau })\mathrm{\Phi }_0={\displaystyle \frac{2}{c}}\tau ^r(\stackrel{}{S}\stackrel{}{L})\mathrm{\Phi }_0{\displaystyle \frac{4}{c}}S^r(\stackrel{}{S}\stackrel{}{\tau })^2\mathrm{\Phi }_0`$
$`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Phi }_3{\displaystyle \frac{1}{2}}S^r\mathrm{\Phi }_2={\displaystyle \frac{1}{4}}\mathrm{\Phi }_3`$
$`(\stackrel{}{S}\stackrel{}{L})\mathrm{\Phi }_0`$ $`=`$ $`{\displaystyle \frac{c}{2}}\mathrm{\Phi }_2`$
$`(\stackrel{}{S}\stackrel{}{L})\mathrm{\Phi }_1`$ $`=`$ $`2(\stackrel{}{S}\stackrel{}{L})S^r\mathrm{\Phi }_0`$
$`=`$ $`2S_j[L_j,S^r]\mathrm{\Phi }_0+\mathrm{\hspace{0.25em}2}\{S_j,S^r\}L_j\mathrm{\Phi }_0\mathrm{\hspace{0.25em}2}S^r(\stackrel{}{S}\stackrel{}{L})\mathrm{\Phi }_0`$
$`=`$ $`2iS_jฯต_{jkl}x_kS_l\mathrm{\Phi }_0{\displaystyle \frac{c}{2}}\mathrm{\Phi }_3`$
$`=`$ $`ฯต_{jkl}x_jฯต_{klm}S_m\mathrm{\Phi }_0{\displaystyle \frac{c}{2}}\mathrm{\Phi }_3=\mathrm{\Phi }_1{\displaystyle \frac{c}{2}}\mathrm{\Phi }_3`$
$`(\stackrel{}{S}\stackrel{}{L})\mathrm{\Phi }_2`$ $`=`$ $`{\displaystyle \frac{2}{c}}(\stackrel{}{S}\stackrel{}{L})^2\mathrm{\Phi }_0={\displaystyle \frac{2}{c}}\left({\displaystyle \frac{1}{4}}L^2{\displaystyle \frac{1}{2}}\stackrel{}{S}\stackrel{}{L}\right)\mathrm{\Phi }_0={\displaystyle \frac{c}{2}}\mathrm{\Phi }_0{\displaystyle \frac{1}{2}}\mathrm{\Phi }_2`$
$`(\stackrel{}{S}\stackrel{}{L})\mathrm{\Phi }_3`$ $`=`$ $`{\displaystyle \frac{4}{c}}(\stackrel{}{S}\stackrel{}{L})S^r(\stackrel{}{S}\stackrel{}{L})\mathrm{\Phi }_0={\displaystyle \frac{2}{c}}S^r(\stackrel{}{S}\stackrel{}{L})\mathrm{\Phi }_0cS^r\mathrm{\Phi }_0`$
$`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Phi }_3{\displaystyle \frac{c}{2}}\mathrm{\Phi }_1`$
and thus
$`2(\stackrel{}{S}\stackrel{}{\tau }S^r\tau ^r)\tau ^r\mathrm{\Phi }_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Phi }_1`$
$`2(\stackrel{}{S}\stackrel{}{\tau }S^r\tau ^r)\tau ^r\mathrm{\Phi }_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Phi }_0`$
$`2(\stackrel{}{S}\stackrel{}{\tau }S^r\tau ^r)\tau ^r\mathrm{\Phi }_2`$ $`=`$ $`0`$
$`2(\stackrel{}{S}\stackrel{}{\tau }S^r\tau ^r)\tau ^r\mathrm{\Phi }_3`$ $`=`$ $`0.`$
Using these relations, it is easy to verify that the vectors $`\mathrm{\Phi }_0,\mathrm{},\mathrm{\Phi }_3`$ are orthonormal on $`L^2(S^2)`$. We take for the wave function $`\mathrm{\Psi }`$ the ansatz
$`\mathrm{\Psi }^{\alpha ua}(t,r,\vartheta ,\phi )`$ (11)
$`=`$ $`e^{i\omega t}{\displaystyle \frac{\sqrt{T(r)}}{r}}(\alpha (r)\mathrm{\Phi }_0^{\alpha a}(\vartheta ,\phi )\delta _{u,1}+\gamma (r)\mathrm{\Phi }_2^{\alpha a}(\vartheta ,\phi )\delta _{u,1}`$
$`+i\beta (r)\mathrm{\Phi }_1^{\alpha a}(\vartheta ,\phi )\delta _{u,2}+i\delta (r)\mathrm{\Phi }_3^{\alpha a}(\vartheta ,\phi )\delta _{u,2})`$
with real functions $`\alpha `$, $`\beta `$, $`\gamma `$, and $`\delta `$, where $`\omega >0`$ is the energy of the Dirac particle. This is the simplest ansatz for which the Dirac equation (4) reduces to a consistent set of ODEs. Namely, we obtain the following system of ODEs for the four-component wave function $`\mathrm{\Phi }:=(\alpha ,\beta ,\gamma ,\delta )`$,
$$\sqrt{A}\mathrm{\Phi }^{}=\left(\begin{array}{cccc}\frac{w}{r}& \omega Tm& \frac{c}{r}& 0\\ \omega Tm& \frac{w}{r}& 0& \frac{c}{r}\\ \frac{c}{r}& 0& 0& \omega Tm\\ 0& \frac{c}{r}& \omega Tm& 0\end{array}\right)\mathrm{\Phi }.$$
(12)
Similar to , we consider the system of $`(2j+1)`$ Dirac wave functions obtained from (11) by applying the ladder operators $`J_\pm `$. Substituting the Einstein and YM equations and using the ansatz (1) and (2), we get the following system of ODEs,
$`rA^{}`$ $`=`$ $`1A{\displaystyle \frac{1}{e^2}}{\displaystyle \frac{(1w^2)^2}{r^2}}`$ (13)
$`\mathrm{\hspace{0.25em}2}\omega T^2(\alpha ^2+\beta ^2+\gamma ^2+\delta ^2){\displaystyle \frac{2Aw_{}^{}{}_{}{}^{2}}{e^2}}`$
$`{\displaystyle \frac{2rA}{T}}T^{}`$ $`=`$ $`A1+{\displaystyle \frac{1}{e^2}}{\displaystyle \frac{(1w^2)^2}{r^2}}`$ (14)
$`\mathrm{\hspace{0.25em}2}\omega T^2(\alpha ^2+\beta ^2+\gamma ^2+\delta ^2){\displaystyle \frac{2Aw_{}^{}{}_{}{}^{2}}{e^2}}`$
$`+T\left[2m(\alpha ^2\beta ^2+\gamma ^2\delta ^2)+{\displaystyle \frac{4c}{r}}(\alpha \delta +\beta \gamma )+{\displaystyle \frac{4w}{r}}\alpha \beta \right]`$
$`r^2Aw^{\prime \prime }`$ $`=`$ $`w(1w^2)+e^2rT\alpha \beta {\displaystyle \frac{1}{2}}r^2A^{}w^{}+{\displaystyle \frac{r^2AT^{}w^{}}{T}}.`$ (15)
Here (13) and (14) are the Einstein equations, and (15) is the YM equation. Notice that the YM equation does not depend on $`\gamma `$ and $`\delta `$; moreover the lower two rows in the Dirac equation (12) are independent of $`w`$. This means that the Dirac particles couple to the YM field only via the spinor functions $`\alpha `$ and $`\beta `$. Indeed, a main difficulty here as compared to the two-spinor problem will be to control the behavior of $`\gamma `$ and $`\delta `$.
For later use, we also give the equations for the following composite functions,
$`r^2(Aw^{})^{}`$ $`=`$ $`w(1w^2)+e^2rT\alpha \beta +{\displaystyle \frac{1}{2}}r^2{\displaystyle \frac{(AT^2)^{}}{T^2}}w^{}`$ (16)
$`r(AT^2)^{}`$ $`=`$ $`\mathrm{\hspace{0.25em}4}\omega T^4(\alpha ^2+\beta ^2+\gamma ^2+\delta ^2){\displaystyle \frac{4AT^2w_{}^{}{}_{}{}^{2}}{e^2}}`$ (17)
$`+T^3[2m(\alpha ^2\beta ^2+\gamma ^2\delta ^2)+{\displaystyle \frac{4c}{r}}(\alpha \delta +\beta \gamma )+{\displaystyle \frac{4w}{r}}\alpha \beta ].`$
Also, it is quite remarkable and will be useful later that for $`\omega =0`$, the squared Dirac equation splits into separate equations for $`(\alpha ,\gamma )`$ and $`(\beta ,\delta )`$; namely from (12),
$`\sqrt{A}_r(\sqrt{A}_r\mathrm{\Phi })=\left(m^2+{\displaystyle \frac{c^2}{r^2}}\right)\mathrm{\Phi }`$ (26)
$`+\left[\sqrt{A}\left(\begin{array}{cccc}\left(\frac{w}{r}\right)^{}& 0& \frac{c}{r^2}& 0\\ 0& \left(\frac{w}{r}\right)^{}& 0& \frac{c}{r^2}\\ \frac{c}{r^2}& 0& 0& 0\\ 0& \frac{c}{r^2}& 0& 0\end{array}\right)+\left(\begin{array}{cccc}\frac{w^2}{r^2}& 0& \frac{wc}{r^2}& 0\\ 0& \frac{w^2}{r^2}& 0& \frac{wc}{r^2}\\ \frac{wc}{r^2}& 0& 0& 0\\ 0& \frac{wc}{r^2}& 0& 0\end{array}\right)\right]\mathrm{\Phi }.`$
## 3 Non-Existence Results
As in , we consider the situation where $`r=\rho >0`$ is the event horizon of a black hole, i.e. $`A(\rho )=0`$, and $`A(\rho )>0`$ if $`r>\rho `$. We again make (cf. ) suitable assumptions on the regularity of the event horizon:
The volume element $`\sqrt{|detg_{ij}}|=|\mathrm{sin}\vartheta |r^2A^1T^2`$ is smooth and non-zero on the horizon; i.e.
$$T^2A^1,T^2AC^1([\rho ,\mathrm{})).$$
The strength of the Yang-Mills field $`F_{ij}`$ is given by
$$\text{Tr}(F_{ij}F^{ij})=\frac{2Aw_{}^{}{}_{}{}^{2}}{r^2}+\frac{(1w^2)^2}{r^4}.$$
We assume that it is bounded near the horizon; i.e.
$`w`$ and $`Aw_{}^{}{}_{}{}^{2}`$ are bounded for $`\rho <r<\rho +\epsilon `$. (27)
Furthermore, the spinors should be normalizable outside and away from the event horizon, i.e.
$$_{\rho +\epsilon }^{\mathrm{}}|\mathrm{\Phi }|^2\frac{T}{\sqrt{A}}<\mathrm{}\text{for every }\epsilon >0.$$
(28)
Finally, we assume that the metric functions and the YM potential satisfy a power ansatz near the event horizon. More precisely, setting
$$ur\rho ,$$
we assume the ansatz
$`A(r)`$ $`=`$ $`A_0u^s+o(u^s)`$ (29)
$`ww_0`$ $`=`$ $`w_1u^\kappa +o(u^\kappa )`$ (30)
with real coefficients $`A_00`$ and $`w_1`$, powers $`s,\kappa >0`$ and $`w_0=lim_{r\rho }w(r)`$. Here and in what follows,
$$f(u)=o(u^\nu )\text{means that }\delta >0\text{ with }\underset{r\rho }{lim\; sup}|u^{\nu \delta }f(u)|<\mathrm{}.$$
Also, we shall always assume that the derivatives of a function in $`o(u^\nu )`$ have the natural decay properties; more precisely,
$$f(u)=o(u^\nu )\text{implies that}f^{(n)}(u)=o(u^{\nu n}).$$
According to $`๐_๐`$, (29) yields that $`T`$ also satisfies a power law, more precisely
$$T(r)u^{\frac{s}{2}}+o(u^{\frac{s}{2}}).$$
(31)
Our main result is the following.
###### Theorem 3.1.
Under the above assumptions, the only black hole solutions of the EDYM equations (12)โ(15) are either the Bartnik-McKinnon black hole solutions of the EYM equations, or
$$s=\frac{4}{3}\text{and}\kappa =\frac{1}{3}.$$
(32)
In (32), the so-called exceptional case, the spinors behave near the horizon like
$$(\alpha \beta )(r)u^{\frac{1}{3}}+o(u^{\frac{1}{3}}),\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}<(\gamma \delta )(\rho )<\mathrm{}.$$
(33)
Our method for the proof of this theorem is to assume a black hole solution with $`\mathrm{\Phi }0`$, and to show that this implies (32) and (33). The proof, which is split up into several parts, is given in Sections 47.
In Section 8, we will analyze the exceptional case. It is shown numerically that the ansatz (32),(33) does not yield global solutions of the EDYM equations. From this we conclude that for all black hole solutions of our EDYM system, the Dirac spinors must vanish identically outside of the event horizon.
## 4 Proof that $`\omega =0`$
Let us assume that there is a solution of the EDYM equations where the spinors are not identically zero, $`\mathrm{\Phi }0`$. In this section we will show that then $`\omega `$ must be zero. First we shall prove that the norm of the spinors $`|\mathrm{\Phi }|`$ is bounded from above and below near the event horizon. We distinguish between the two cases where $`A^{\frac{1}{2}}`$ is or is not integrable near the event horizon.
###### Lemma 4.1.
If $`A^{\frac{1}{2}}`$ is integrable near the event horizon $`r=\rho `$, then there are positive constant $`k`$ and $`\epsilon `$ such that
$$\frac{1}{k}|\mathrm{\Phi }(r)|^2k,\text{if }\rho <r<\rho +\epsilon \text{.}$$
(34)
Proof. Writing (12) as $`\sqrt{A}\mathrm{\Phi }^{}=M\mathrm{\Phi }`$, we have
$`{\displaystyle \frac{1}{2}}\sqrt{A}{\displaystyle \frac{d}{dr}}|\mathrm{\Phi }|^2={\displaystyle \frac{1}{2}}\mathrm{\Phi }^t(M+M^{})\mathrm{\Phi }`$ (35)
$`=`$ $`{\displaystyle \frac{w}{r}}(\alpha ^2\beta ^2)\mathrm{\hspace{0.25em}2}m(\alpha \beta +\gamma \delta )+{\displaystyle \frac{2c}{r}}(\alpha \gamma \beta \delta )`$
$``$ $`{\displaystyle \frac{w}{r}}(\alpha ^2\beta ^2)+m(\alpha ^2+\beta ^2+\gamma ^2+\delta ^2)+{\displaystyle \frac{2c}{r}}(\alpha \gamma \beta \delta )`$
$``$ $`c_1|\mathrm{\Phi }|^2.`$
Here the constant $`c_1`$ is independent of $`r(\rho ,\rho +1]`$, since $`w`$ is bounded near the horizon according to assumption $`๐_{\mathrm{๐๐}}`$. Since we are assuming that $`\mathrm{\Phi }0`$ in $`r>\rho `$, the uniqueness theorem for solutions of ODEs yields that $`|\mathrm{\Phi }|^2>0`$ on $`(\rho ,\rho +1]`$. Then dividing (35) by $`\frac{1}{2}\sqrt{A}|\mathrm{\Phi }|^2`$ and integrating from $`r_1`$ to $`r_2`$, $`\rho <r_1<r_2`$, we get
$$\left|\mathrm{log}|\mathrm{\Phi }(r_2)|^2\mathrm{log}|\mathrm{\Phi }(r_1)|^2\right|\mathrm{\hspace{0.33em}2}c_1_{r_1}^{r_2}A^{\frac{1}{2}}(r)๐r.$$
Taking the limit $`r_1\rho `$ in this last inequality gives the desired result.
###### Lemma 4.2.
If $`A^{\frac{1}{2}}`$ is not integrable near the event horizon $`r=\rho `$ and $`\omega 0`$, then there are positive constants $`k`$ and $`\epsilon `$ such that
$$\frac{1}{k}|\mathrm{\Phi }(r)|^2k\text{if }\rho <r<\rho +\epsilon \text{.}$$
(36)
Proof. Define the matrix $`J`$ by
$$J=\left(\begin{array}{cccc}1\frac{m}{\omega T}& \frac{w}{r\omega T}& 0& \frac{c}{r\omega T}\\ \frac{w}{r\omega T}& 1+\frac{m}{\omega T}& \frac{c}{r\omega T}& 0\\ 0& \frac{c}{r\omega T}& 1\frac{m}{\omega T}& 0\\ \frac{c}{r\omega T}& 0& 0& 1+\frac{m}{\omega T}\end{array}\right)$$
and notice that, since $`T(r)\mathrm{}`$ as $`r\rho `$, $`J`$ is close to the identity matrix for $`r`$ near $`\rho `$. If we let
$$F(r)=<\text{ }\mathrm{\Phi }(r),J(r)\mathrm{\Phi }(r)>,$$
then a straightforward calculation yields that
$$F^{}=<\text{ }\mathrm{\Phi }(r),J^{}(r)\mathrm{\Phi }(r)>.$$
In a manner similar to that in , we can prove that $`|J^{}|`$ is integrable near $`r=\rho `$, and as in , it follows that (36) holds.
###### Lemma 4.3.
If $`\mathrm{\Phi }0`$ for $`r>\rho `$, then $`\omega =0`$.
Proof. Assume that $`\omega 0`$. We write the $`(AT^2)^{}`$ equation (17) as
$`r(AT^2)^{}=4\omega T^4|\mathrm{\Phi }|^2{\displaystyle \frac{4Aw_{}^{}{}_{}{}^{2}}{e^2}}T^2`$ (37)
$`+\left[2m(\alpha ^2\beta ^2+\gamma ^2\delta ^2)+{\displaystyle \frac{4w}{r}}\alpha \beta +{\displaystyle \frac{4c}{r}}(\alpha \delta +\beta \gamma )\right]T^3.`$
According to hypothesis $`๐_{\mathrm{๐๐}}`$, the left side of this equation is bounded near the event horizon. The Lemmas 4.1 and 4.2 together with $`๐_{\mathrm{๐๐}}`$ imply that the coefficients of $`T^4`$, $`T^3`$, and $`T^2`$ in this equation are all all bounded, and that the coefficient of $`T^4`$ is bounded away from zero near $`r=\rho `$. Assumption $`๐_{\mathrm{๐๐}}`$ implies that $`T(r)\mathrm{}`$ as $`r\rho `$. Hence the right side of (37) diverges as $`r\rho `$. This is a contradiction.
## 5 Reduction to the Case $`\alpha (\rho )=0`$, $`\beta (\rho )0`$
Since $`\omega =0`$, the Dirac equation (12) reduces to
$$\sqrt{A}\mathrm{\Phi }^{}=\left(\begin{array}{cccc}w/r& m& c/r& 0\\ m& w/r& 0& c/r\\ c/r& 0& 0& m\\ 0& c/r& m& 0\end{array}\right)\mathrm{\Phi }M\mathrm{\Phi }.$$
(38)
The following Lemma gives some global information on the behavior of the solutions to (38).
###### Lemma 5.1.
The function $`(\alpha \beta +\gamma \delta )`$ is strictly positive, decreasing, and tends to zero as $`r\mathrm{}`$.
Proof. A straightforward calculation gives
$$\sqrt{A}(\alpha \beta +\gamma \delta )^{}=m|\mathrm{\Phi }|^2,$$
so that $`(\alpha \beta +\gamma \delta )(r)`$ is a strictly decreasing function, and thus has a (possibly infinite) limit as $`r\mathrm{}`$. Since $`|\mathrm{\Phi }|^22|\alpha \beta +\gamma \delta |`$, we see that the normalization condition (28) holds only if this limit is zero. It follows that $`(\alpha \beta +\gamma \delta )`$ is strictly positive.
Next we want to show that the spinors have a (possibly infinite) limit as $`r\rho `$. When $`A^{\frac{1}{2}}`$ is integrable near the event horizon, it is an immediate consequence of Lemma 4.1 that this limit exists and is even finite.
###### Corollary 5.2.
If $`A^{\frac{1}{2}}`$ is integrable near the horizon, then $`\mathrm{\Phi }`$ has a finite limit for $`r\rho `$.
Proof. We can integrate (38) from $`r_1`$ to $`r_2`$, $`\rho <r_1<r_2`$,
$$\mathrm{\Phi }(r_2)\mathrm{\Phi }(r_1)=_{r_1}^{r_2}A^{\frac{1}{2}}(r)M(r)\mathrm{\Phi }(r)๐r.$$
Lemma 4.1 yields that the right side converges as $`r_1\rho `$, and hence $`\mathrm{\Phi }`$ has a finite limit.
In the case when $`A^{\frac{1}{2}}`$ is not integrable near the horizon, we argue as follows. According to the power ansatz (30), the matrix in (38) has a finite limit on the horizon. Exactly as shown in \[7, Section 5\] using the stable manifold theorem, there are fundamental solutions of the Dirac equation (that is, a basis of solutions of the ODE (38)) which behave near the event horizon exponentially like $`\mathrm{exp}(\lambda _jA^{\frac{1}{2}})`$, where $`\lambda _j\mathrm{IR}`$ are the eigenvalues for $`r\rho `$ of the matrix in (38) (notice that the $`\lambda _j`$ are real since they are the eigenvalues of a symmetric matrix). Thus for any linear combination of these fundamental solutions, the spinor functions are monotone in a neighborhood of the event horizon, and hence as $`r\rho `$, $`\mathrm{\Phi }`$ has a limit in $`\mathrm{IR}\{\pm \mathrm{}\}`$. We set
$$\mathrm{\Phi }(\rho )=\underset{r\rho }{lim}\mathrm{\Phi }(r),(\alpha \beta )(\rho )=\underset{r\rho }{lim}(\alpha \beta )(r).$$
###### Proposition 5.3.
$`(\alpha \beta )(\rho )=0`$.
Proof. We consider the $`(Aw^{})^{}`$ equation (16) with $`\omega =0`$,
$$r^2(Aw^{})^{}=w(1w^2)+e^2r(\sqrt{A}T)\frac{\alpha \beta }{\sqrt{A}}+\frac{r^2(AT^2)^{}}{2AT^2}Aw^{}.$$
(39)
Suppose that
$$(\alpha \beta )(\rho )>\mathrm{\hspace{0.33em}0}.$$
(40)
From hypotheses $`๐_๐`$ and and $`๐_{\mathrm{๐๐}}`$, we se that the coefficient of $`\alpha \beta A^{\frac{1}{2}}`$ is positive near $`r=\rho `$, and the other terms on the right side of (39) are bounded. Thus we may write (39) in the form
$$(Aw^{})^{}=\mathrm{\Phi }(r)+\frac{\mathrm{\Psi }((r)}{\sqrt{A(r)}},$$
(41)
where $`\mathrm{\Phi }`$ is bounded and $`\mathrm{\Psi }>0`$ near $`\rho `$. Thus we can find constants $`\mathrm{\Phi }_0`$, $`\mathrm{\Psi }_0`$ satisfying
$$(Aw^{})^{}>\mathrm{\Phi }_0+\frac{\mathrm{\Psi }_0}{\sqrt{A(r)}},\mathrm{\Psi }_0>0,$$
(42)
for $`r`$ near $`\rho `$. Then exactly as in \[2, Section 3\], it follows that the spinors must vanish in $`r>\rho `$.
If on the other hand
$$(\alpha \beta )(\rho )<\mathrm{\hspace{0.33em}0}.$$
(43)
then (41) holds with $`\mathrm{\Psi }(r)<0`$ near $`\rho `$. Thus
$$(Aw^{})^{}=\mathrm{\Phi }(r)\frac{\mathrm{\Psi }((r)}{\sqrt{A(r)}}.$$
(44)
Setting $`\stackrel{~}{w}=w`$, (44) becomes
$$(A\stackrel{~}{w}^{})^{}=\mathrm{\Phi }(r)\frac{\mathrm{\Psi }((r)}{\sqrt{A(r)}},$$
where $`\mathrm{\Psi }(r)>0`$ for $`r`$ near $`\rho `$. Thus we see that (41) holds for $`w`$ replaced by $`\stackrel{~}{w}`$. This again leads to a contradiction.
The next proposition rules out the case that both $`\alpha `$ and $`\beta `$ vanish on the event horizon.
###### Proposition 5.4.
Either $`\alpha (\rho )=0`$, $`\beta (\rho )0`$ or $`\alpha (\rho )0`$, $`\beta (\rho )=0`$.
Proof. Suppose that
$$\alpha (\rho )=\mathrm{\hspace{0.33em}0}=\beta (\rho ).$$
(45)
According to Lemma 5.1, $`\gamma `$ and $`\delta `$ cannot both vanish on the event horizon. Using (38), we have for $`r`$ near $`\rho `$,
$`\sqrt{A}\alpha ^{}`$ $`=`$ $`{\displaystyle \frac{c}{r}}\gamma +o(1)`$ (46)
$`\sqrt{A}\beta ^{}`$ $`=`$ $`{\displaystyle \frac{c}{r}}\delta +o(1).`$ (47)
If $`A^{\frac{1}{2}}`$ is not integrable near the event horizon, these equations show that $`\gamma (\rho )`$ and $`\delta (\rho )`$ are finite (otherwise multiplying (46) and (47) by $`A^{\frac{1}{2}}`$ and integrating would contradict (45)); if $`A^{\frac{1}{2}}`$ is integrable near $`\rho `$, Corollary 5.2 shows that $`\gamma (\rho )`$ and $`\delta (\rho )`$ are again finite.
From (17) with $`\omega =0`$ we have
$`r(AT^2)^{}`$ $`=`$ $`\left[2m(\alpha ^2\beta ^2+\gamma ^2\delta ^2)+{\displaystyle \frac{4w}{r}}\alpha \beta +{\displaystyle \frac{4c}{r}}(\alpha \delta +\beta \gamma )\right]T^3`$ (48)
$`{\displaystyle \frac{4}{e^2}}(Aw_{}^{}{}_{}{}^{2})T^2.`$
Since the coefficients of $`T^3`$ and $`T^2`$ are bounded, as is the left-hand side, we conclude that, since $`T(r)\mathrm{}`$ as $`r\rho `$, the coefficient of $`T^3`$ must vanish on the horizon,
$$\left[2m(\alpha ^2\beta ^2+\gamma ^2\delta ^2)+\frac{4w}{r}\alpha \beta +\frac{4c}{r}(\alpha \delta +\beta \gamma )\right]_{r=\rho }=\mathrm{\hspace{0.33em}0}.$$
(49)
As a consequence, $`\gamma (\rho )^2=\delta (\rho )^2`$, and Lemma 5.1 yields that
$$\gamma (\rho )=\delta (\rho )\mathrm{\hspace{0.33em}0}.$$
(50)
Furthermore from (46) and (47), for $`r`$ near $`\rho `$,
$$\text{sgn }\alpha (r)=\text{sgn }\gamma (r)\text{and}\text{sgn }\beta (r)=\text{sgn }\delta (r).$$
(51)
From (50) and (51), we see that for $`r`$ near $`\rho `$, the spinors must lie in the shaded areas in one of the two configurations (I) or (II) in Figure 1.
Now we claim that in either configuration (I) or (II), the shaded regions are invariant. For the proof, we consider the Dirac equation (38). One easily checks that the shaded regions in the $`\alpha /\beta `$-plots are invariant, provided that $`\gamma `$ and $`\delta `$ are as depicted in their shaded regions. Similarly, one verifies that the shaded regions in the $`\gamma /\delta `$-plots are invariant, provided that $`\alpha `$ and $`\beta `$ lie in the shaded regions. Moreover, Lemma 5.1 shows that the spinors cannot leave their regions simultaneously (i.e. for the same $`r`$). This proves the claim.
Next we consider the situation for large $`r`$. In the limit $`r\mathrm{}`$, the matrix $`M`$ in (38) goes over to the matrix $`S`$ given by
$$S=\left(\begin{array}{cccc}0& m& 0& 0\\ m& 0& 0& 0\\ 0& 0& 0& m\\ 0& 0& m& 0\end{array}\right).$$
In $`S`$, the non-zero $`2\times 2`$ upper and lower triangular blocks,
$$\left(\begin{array}{cc}0& m\\ m& 0\end{array}\right),$$
have eigenvectors $`(1,1)^t`$ and $`(1,1)^t`$ with corresponding eigenvalues $`m`$ and $`m`$, respectively. Since the system of ODEs
$$\sqrt{A}\mathrm{\Phi }^{}=S\mathrm{\Phi }$$
splits into separate equations for $`(\alpha ,\beta )`$ and $`(\gamma ,\delta )`$, we see that $`(\alpha (r),\beta (r))`$ must be a linear combination of $`e^{c(r)r}(1,1)^t`$ and $`e^{d(r)r}(1,1)^t`$, where the functions $`c`$ and $`d`$ are close to $`m`$. Since the spinors are assumed to be normalizable (i.e. (28) holds), and are non-zero for $`r>\rho `$, it follows that for large $`r`$, the spinors are close to a constant multiple of $`e^{c(r)r}(1,1)^t`$, and thus for large $`r`$, $`\text{sgn }\alpha (r)=\text{sgn }\beta (r)`$. Similarly, for large $`r`$, $`\text{sgn }\gamma (r)=\text{sgn }\delta (r)`$. This is a contradiction to the shaded invariant regions of Figure 1.
The two cases in Proposition 5.4 can be treated very similarly. Therefore we shall in what follows restrict attention to the first case. Furthermore, we know from Lemma 5.1 and Proposition 5.3 that $`(\gamma \delta )(\rho )>0`$. Using linearity of the Dirac equation, we can assume that both $`\gamma (\rho )`$ and $`\delta (\rho )`$ are positive. Hence the remaining problem is to consider the case where
$$\alpha (\rho )=\mathrm{\hspace{0.33em}0},\beta (\rho )\mathrm{\hspace{0.33em}0},\gamma (\rho ),\delta (\rho )>\mathrm{\hspace{0.33em}0}.$$
(52)
## 6 Proof that $`A^{\frac{1}{2}}`$ is Integrable Near the Event Horizon
In this section we shall assume that $`A^{\frac{1}{2}}`$ is not integrable near the event horizon and deduce a contradiction. We work with the power ansatz (29),(30) and thus assume that $`s2`$.
We first consider the case $`w_00`$. The first component of the squared Dirac equation (26) is
$`\sqrt{A}_r(\sqrt{A}_r\alpha )`$ (53)
$`=`$ $`\left[m^2+{\displaystyle \frac{c^2+w^2}{r^2}}+\sqrt{A}\left({\displaystyle \frac{w}{r}}\right)^{}\right]\alpha +{\displaystyle \frac{c(w\sqrt{A})}{r^2}}\gamma .`$
The square bracket is bounded according to $`๐_{\mathrm{๐๐}}`$. Since $`\alpha (\rho )=0`$ and $`\gamma (\rho )>0`$, our assumption $`w_00`$ implies that the right side of (53) is bounded away from zero near the event horizon, i.e. there are constants $`\delta ,\epsilon `$ with
$$\pm \sqrt{A}_r(\sqrt{A}_r\alpha )\delta \text{for }\rho <r<\rho +\epsilon ,$$
where โ$`\pm `$โ corresponds to the two cases $`w_0>1`$ and $`w_0<1`$, respectively. We multiply this inequality by $`A^{\frac{1}{2}}`$ and integrate from $`r_1`$ to $`r_2`$, $`\rho <r_1<r_2`$,
$$\pm \sqrt{A}_r\alpha |_{r_1}^{r_2}\delta _{r_1}^{r_2}A^{\frac{1}{2}}.$$
The right side diverges as $`r_1\rho `$, and thus $`lim_{r\rho }\sqrt{A}_r\alpha =\mathrm{}`$. Hence near the event horizon, $`_r\alpha A^{\frac{1}{2}}`$, and integrating once again yields that $`lim_{r\rho }\alpha =\pm \mathrm{}`$, in contradiction to $`\alpha (\rho )=0`$.
Suppose now that $`w_0=0`$. We first consider the $`A`$-equation (13), which since $`\omega =0`$ becomes
$$rA^{}=\mathrm{\hspace{0.33em}1}A\frac{1}{e^2}\frac{(1w^2)^2}{r^2}\frac{2}{e^2}Aw_{}^{}{}_{}{}^{2}.$$
(54)
Employing the power ansatz (29),(30) gives
$$O(u^{s1})=\mathrm{\hspace{0.33em}1}+O(u^s)\frac{1}{e^2r^2}+O(u^{2\kappa })+O(u^{s+2\kappa 2}).$$
(55)
Here and in what follows,
$$f(u)=O(u^\nu )\text{means that}\underset{r\rho }{lim}u^\nu f(u)\text{ is finite and non-zero},$$
also we omit the expressions โ$`o(u^\nu )`$.โ The constant term in (55) must vanish, and thus $`e^2\rho ^2=1`$. Using also that $`O(u^s)`$ is of higher order, (55) reduces to
$$O(u^{s1})=O(u)+O(u^{2\kappa })+O(u^{s+2\kappa 2}).$$
(56)
Suppose first that $`s>2`$. Then (56) yields that $`\kappa =\frac{1}{2}`$. Substituting our power ansatz into the $`Aw^{}`$-equation (16) gives
$$O(u^{s\frac{3}{2}})=O(u^{\frac{1}{2}})+e^2rT\alpha \beta $$
and thus $`\alpha \beta =O(u^{\frac{1+s}{2}})`$. Since $`\beta (\rho )0`$, we conclude that there are constants $`c_1,\delta >0`$ with
$$|\alpha |c_1u^{\frac{1+s}{2}}\text{for }\rho <r<\rho +\delta .$$
(57)
From this one sees that the first summand on the right side of (53) is of higher order; more precisely,
$$\sqrt{A}_r(\sqrt{A}_r\alpha )=O(u^{\frac{1}{2}}).$$
Multiplying by $`A^{\frac{1}{2}}`$ and integrating twice, we conclude that
$$\alpha =O(u^{\frac{5}{2}s}),$$
and this contradicts (57).
The final case to consider is $`w_0=0`$ and $`s=2`$. Now the $`Aw^{}`$-equation (16) gives
$$O(u^\kappa )=O(u^\kappa )+e^2T\alpha \beta $$
and thus $`\alpha =o(u)`$. This gives a contradiction in (53) unless $`w\sqrt{A}=o(u)`$, and we conclude that $`\kappa =1`$. Now consider the Dirac equation (38). Since $`w(\rho )=0`$, the eigenvalues of the matrix in (38) on the horizon are $`\lambda =\pm \sqrt{m^2+c^2/\rho ^2}`$. As a consequence, the fundamental solutions behave near the horizon $`u^{\pm \sqrt{m^2+c^2/\rho ^2}}`$. The boundary conditions (52) imply that $`\alpha u^{+\sqrt{m^2+c^2/\rho ^2}}`$, whereas $`\beta ,\gamma ,\delta u^{\sqrt{m^2+c^2/\rho ^2}}`$, and we conclude that
$$(\alpha \delta +\beta \gamma )(\rho )>\mathrm{\hspace{0.33em}0}.$$
(58)
Next we consider the $`AT^2`$-equation (17), which for $`\omega =0`$ takes the form (48). It is convenient to introduce for the square bracket the short notation
$$[]=\mathrm{\hspace{0.33em}2}m(\alpha ^2\beta ^2+\gamma ^2\delta ^2)+\frac{4w}{r}\alpha \beta +\frac{4c}{r}(\alpha \delta +\beta \gamma ).$$
(59)
We define the matrix $`B`$ by
$$B=\left(\begin{array}{cccc}m& w/r& 0& c/r\\ w/r& m& c/r& 0\\ 0& c/r& m& 0\\ c/r& 0& 0& m\end{array}\right).$$
A short calculation shows that
$$[]=\mathrm{\hspace{0.33em}2}<\text{ }\mathrm{\Psi },B\mathrm{\Psi }>,$$
and furthermore, using the Dirac equation (38),
$`[]^{}`$ $`=`$ $`2<\text{ }\mathrm{\Psi },B^{}\mathrm{\Psi }>=\mathrm{\hspace{0.33em}4}\alpha \beta \left({\displaystyle \frac{w}{r}}\right)^{}{\displaystyle \frac{4c}{r^2}}(\alpha \delta +\beta \gamma )`$ (60)
$`=`$ $`{\displaystyle \frac{4w^{}}{r}}\alpha \beta {\displaystyle \frac{4w}{r^2}}\alpha \beta {\displaystyle \frac{4c}{r^2}}(\alpha \delta +\beta \gamma ).`$
Since $`(\alpha \beta )(\rho )=0`$ and $`(\alpha \delta +\beta \gamma )(\rho )>0`$ according to (58),
$$[]^{}c_2\text{for }\rho <r<\rho +\delta $$
and a constant $`c_2>0`$ . Integrating on both sides shows that
$$\left|[]\right|c_3u\text{for }\rho <r<\rho +\delta $$
with $`c_3>0`$. As a consequence, the first summand in (48) diverges for $`r\rho `$, whereas the left side and the second summand on the right are bounded in this limit. This is a contradiction.
We conclude that $`A^{\frac{1}{2}}`$ must be integrable near the event horizon, and so $`s<2`$.
## 7 Proof of the Main Theorem
In this section we shall analyze the EDYM equations with the power ansatz (29),(30) near the event horizon. We will derive restrictions for the powers $`s`$ and $`k`$ until only the exceptional case (32) of Theorem 3.1 remains. So far, we know from Section 6 that $`s<2`$. A simple lower bound follows from the $`A`$-equation (13) which for $`\omega =0`$ simplifies to (54). Namely in view of hypothesis $`๐_{\mathrm{๐๐}}`$, the right-hand side of (54) is bounded, and thus $`s1`$. The case $`s=1`$ is excluded just as in by matching the spinors across the horizon and applying a radial flux argument. Thus it remains only to consider $`s`$ in the range
$$1<s<\mathrm{\hspace{0.33em}2}.$$
(61)
We begin by deriving a power expansion for $`\alpha `$ near the event horizon.
###### Lemma 7.1.
Suppose that $`w_00`$ or $`\kappa s/2`$. Then the function $`\alpha `$ behaves near the horizon as
$$\alpha =\alpha _0u^\sigma +o(u^\sigma ),\alpha _00,$$
(62)
where the power $`\sigma `$ is either
$$\sigma =\mathrm{\hspace{0.33em}1}\frac{s}{2}$$
(63)
or
$$\sigma =\{\begin{array}{cc}2s& \text{if }w_00\\ 2s+\mathrm{min}(\kappa ,s/2)& \text{if }w_0=0\text{.}\end{array}$$
(64)
Proof. We set
$$\sigma =sup\{p:\underset{r\rho }{lim\; sup}|u^p\alpha (r)|<\mathrm{}\}\mathrm{}.$$
(65)
Suppose first that $`\sigma <\mathrm{}`$. Then for every $`\nu <\sigma `$ there are constants $`c>0`$ and $`\epsilon >0`$ with
$$|\alpha (r)|<cu^\nu \text{for }\rho <r<\rho +\epsilon .$$
(66)
We consider the first component of the squared Dirac equation (53) and write it in the form
$$\sqrt{A}_r(\sqrt{A}_r\alpha )=f\alpha +g,$$
(67)
where $`f`$ stands for the square bracket and $`g`$ for the last summand in (53), respectively. Multiplying by $`A^{\frac{1}{2}}`$ and integrating gives
$$\sqrt{A}_r\alpha (r)=_\rho ^rA^{\frac{1}{2}}(f\alpha +g)+C$$
with an integration constant $`C`$. We again multiply by $`A^{\frac{1}{2}}`$ and integrate. Since $`\alpha (\rho )=0`$, we obtain
$$\alpha (r)=_\rho ^rA^{\frac{1}{2}}(s)๐s_\rho ^sA^{\frac{1}{2}}(f\alpha +g)+C_\rho ^rA^{\frac{1}{2}}.$$
(68)
Note that the function $`f`$, introduced as an abbreviation for the square bracket in (53), is bounded near the horizon. Hence (66) yields a polynomial bound for $`|f\alpha |`$. Each multiplication with $`A^{\frac{1}{2}}`$ and integration increases the power by $`1\frac{s}{2}`$, and thus there is a constant $`c_1`$ with
$$_\rho ^rA^{\frac{1}{2}}(s)๐s_\rho ^sA^{\frac{1}{2}}|f\alpha |c_1u^{2s+\nu }\text{for }\rho <r<\rho +\epsilon .$$
(69)
Since $`2s>0`$, (69) is of the order $`o(u^{1\frac{s}{2}+\nu })`$, and thus (68) can be written as
$$\alpha (r)=_\rho ^rA^{\frac{1}{2}}(s)๐s_\rho ^sA^{\frac{1}{2}}g+C_\rho ^rA^{\frac{1}{2}}+o(u^{1\frac{s}{2}+\nu }).$$
(70)
Consider the behavior of the first two summands in (70). The function $`g`$ stands for the last summand in (53). If $`w_01`$, it has a non-zero limit on the horizon. If on the other hand $`w_0=1`$, then $`gu^\kappa `$. Substituting into (70) and integrating, one sees that the first summand in (70) is $`u^\sigma `$ with $`\sigma `$ given by (64). The second summand in (70) vanishes if $`C=0`$, and is $`u^\sigma `$ with $`\sigma `$ as in (63). According to (61), $`1\frac{s}{2}<2s<2s+\mathrm{min}(\kappa ,s/2)`$. Thus the values of $`\sigma `$ in (63) and (64) are different, and so the first two summands in (70) cannot cancel each other. If we choose $`\nu `$ so large that $`1\frac{s}{2}+\nu \sigma `$, (70) yields the Lemma.
Suppose now that $`\sigma `$ given by (65) is infinite. Then choosing
$$\nu =\mathrm{max}[1\frac{s}{2},\mathrm{\hspace{0.25em}2}s+\mathrm{min}(\kappa ,s/2)],$$
we see that the first two summands in (70) are of the order $`O(u^s)`$ with $`s`$ according to (63) and (64), respectively, and the last summand is of higher order. Thus (70) implies that $`\sigma `$ as defined by (65) is finite (namely, equal to the minimum of (63) and (64)), giving a contradiction. Thus $`\sigma `$ is indeed finite.
In the proof of Proposition 5.4, we already observed that the square bracket in the $`AT^2`$-equation (17) vanishes on the horizon (49). Let us now analyze this square bracket in more detail, where we use again the notation (59).
###### Proposition 7.2.
$`\kappa <1`$ and
$$[]=O(u^{\kappa +\sigma })+O(u)$$
(71)
with $`\sigma `$ as in Lemma 7.1.
Proof. The derivative of the square bracket is again given by (60). Now $`\alpha _0=0`$, $`\beta _00`$ and from Lemma 5.1, $`\alpha _0\delta _0+\beta _0\gamma _00`$; thus using (29), (30), and (62), we get, for $`r`$ near $`\rho `$,
$$[]^{}=O(u^{\kappa 1+\sigma })+u^{\sigma +(\kappa )}+O(1),$$
(72)
where we again omitted the expressions โ$`o(u^.)`$โ and we use the notation
$$(\kappa )=\{\begin{array}{cc}\kappa \hfill & \text{if }w_0=0\hfill \\ 0\hfill & \text{if }w_00\hfill \end{array}.$$
Integrating (72) and using that $`[]_{r=\rho }=0`$ according to (49), we obtain that
$$[]=O(u^{\kappa +\sigma })+u^{\sigma +(\kappa )+1}+O(u).$$
(73)
Suppose $`\kappa 1`$. Then $`\kappa +\sigma >1`$ and $`\sigma +(\kappa )+1>1`$, and (73) becomes
$$[]=O(u).$$
We write the $`AT^2`$-equation (48) as
$$r(AT^2)^{}=T^3[]\frac{4AT^2w_{}^{}{}_{}{}^{2}}{e^2}.$$
(74)
Since $`(AT^2)^{}`$ is bounded and $`T^3=O(u^{\frac{3s}{2}})`$ (by virtue of hypothesis $`๐_๐`$), (74) behaves near the event horizon like
$$u^0=O(u^{1\frac{3s}{2}})+O(u^{2\kappa 2}).$$
(75)
Since $`2\kappa 20`$ and $`1\frac{3s}{2}<0`$, the right side of (75) is unbounded as $`r\rho `$, giving a contradiction. We conclude that $`\kappa <1`$.
For $`\kappa <1`$, the second summand in (73) is of higher order than the first, and we get (71).
In the remainder of this section, we shall substitute the power expansions (29)โ(31) and (62) into the EDYM equations and evaluate the leading terms (i.e. the lowest powers in $`u`$). This will amount to a rather lengthy consideration of several cases, each of which has several subcases. We begin with the case $`w_00,\pm 1`$. The $`A`$-equation (13) simplifies to (54). The $`AT^2`$-equation (17) for $`\omega =0`$ takes the form (48), and we can for the square bracket use the expansion of Proposition 7.2. Finally, we also consider the $`Aw^{}`$-equation (16). Using the regularity assumption $`๐_๐`$, we obtain
$`A`$-eqn: $`O(u^{s1})=\mathrm{\hspace{0.33em}1}{\displaystyle \frac{(1w_0^2)^2}{e^2\rho ^2}}+O(u^\kappa )+O(u^{s+2\kappa 2})`$ (76)
$`AT^2`$-eqn: $`u^0=O(u^{\kappa +\sigma \frac{3s}{2}})+O(u^{1\frac{3s}{2}})+O(u^{2\kappa 2})`$ (77)
$`Aw^{}`$-eqn: $`O(u^{s+\kappa 2})=w_0(1w_0^2)+O(u^\kappa )+O(u^{\sigma \frac{s}{2}}).`$ (78)
First consider (76). According to $`๐_{\mathrm{๐๐}}`$, $`s+2\kappa 20`$, and so all powers in (76) are positive. We distinguish between the cases where the power $`s+2\kappa 2`$ is larger, smaller, or equal to the other powers on the right of (76). Making sure in each case that the terms of leading powers may cancel each other, we obtain the cases and conditions
(a) $`\kappa <s+2\kappa 2w_0^2=1\pm e\rho ,\kappa =s1,s{\displaystyle \frac{3}{2}}`$ (79)
(b) $`\kappa =s+2\kappa 2w_0^2=1\pm e\rho ,\kappa =2s,s{\displaystyle \frac{3}{2}}`$ (80)
(c) $`\kappa >s+2\kappa 2>0w_0^2=1\pm e\rho ,\kappa ={\displaystyle \frac{1}{2}},s<{\displaystyle \frac{3}{2}}`$ (81)
(d) $`s+2\kappa 2=0\kappa =1{\displaystyle \frac{s}{2}}.`$ (82)
In Case (a), the relations in (79) imply that
$$1\frac{3s}{2}<\mathrm{\hspace{0.33em}2}s4=\mathrm{\hspace{0.33em}2}\kappa 2.$$
Hence (77) yields $`13s/2=\kappa +\sigma 3s/2`$, so
$$\sigma =\mathrm{\hspace{0.33em}1}\kappa =\mathrm{\hspace{0.33em}2}s.$$
(83)
This is consistent with Lemma 7.1. But we get a contradiction in (78) as follows. Since $`\kappa =s1`$, we have $`s+\kappa 2=2s3>0`$; on the other hand,
$$\sigma \frac{s}{2}=\mathrm{\hspace{0.33em}2}s\frac{s}{2}=\mathrm{\hspace{0.33em}2}\frac{3s}{2}<\mathrm{\hspace{0.33em}0}.$$
Thus the left-hand side of (78) is bounded, but the right-hand side is unbounded as $`r\rho `$. This completes the proof in Case (a).
In Case (b), (78) yields that
$$\sigma \frac{s}{2}.$$
(84)
We consider the two cases (63) and (64) in Lemma 7.1. In the first case, (84) yields that $`s1`$, contradicting (61). In the second case, (84) implies that $`s\frac{4}{3}`$. This contradicts the inequality in (80), and thus completes the proof in Case (b).
In Case (c), the relations in (81) give $`s+\kappa 2=s\frac{3}{2}<0`$, and thus (78) implies that $`s+\kappa 2=\sigma \frac{s}{2}`$, so $`\sigma =\frac{3}{2}(s1)`$. According to Lemma 7.1, $`\sigma =2s`$ or $`\sigma =1\frac{s}{2}`$. In the first of these cases, we conclude that $`s=\frac{7}{5}`$ and $`\sigma =\frac{3}{5}`$. Substituting these powers into (77), we get
$$u^0=O(u^1)+O(u^{\frac{11}{10}})+O(u^1),$$
which clearly yields a contradiction. Thus $`\sigma =1\frac{s}{2}`$, giving
$$s=\frac{5}{4},\kappa =\frac{1}{2},\sigma =\frac{3}{8}.$$
(85)
This case is ruled out in Lemma 7.3 below.
In Case (d), we consider (78). Since $`s+\kappa 2=\kappa <0`$, we obtain that $`s+\kappa 2=\sigma \frac{\sigma }{2}`$ and thus $`\sigma =s1`$. Lemma 7.1 yields the two cases
$`s`$ $`=`$ $`{\displaystyle \frac{4}{3}},\kappa ={\displaystyle \frac{1}{3}},\sigma ={\displaystyle \frac{1}{3}}\text{and}`$ (86)
$`s`$ $`=`$ $`{\displaystyle \frac{3}{2}},\kappa ={\displaystyle \frac{1}{4}},\sigma ={\displaystyle \frac{1}{2}}.`$ (87)
The first of these cases is the exceptional case of Theorem 3.1, and the second case is excluded in Lemma 7.4 below. This concludes the proof of Theorem 3.1 in the case $`w_00,\pm 1`$.
We next consider the case $`w_0=\pm 1`$. Then the expansions (76)โ(78) must be modified to
$`A`$-eqn: $`O(u^{s1})=\mathrm{\hspace{0.33em}1}+O(u^{2\kappa })+O(u^{s+2\kappa 2})`$ (88)
$`AT^2`$-eqn: $`u^0=O(u^{\kappa +\sigma \frac{3s}{2}})+O(u^{1\frac{3s}{2}})+O(u^{2\kappa 2})`$ (89)
$`Aw^{}`$-eqn: $`O(u^{s+\kappa 2})=O(u^\kappa )+O(u^{\sigma \frac{s}{2}}).`$ (90)
One sees immediately that, in order to compensate the constant term in (88), $`s+2\kappa 2`$ must be zero. Hence $`s+\kappa 2=\kappa <0`$, and (90) yields that $`s+\kappa 2=\sigma \frac{s}{2}`$ and thus $`\sigma =s1`$. Now consider Lemma 7.1. In case (63), we get the exceptional case of Theorem 3.1, whereas case (64) yields that
$$s=\frac{3}{2},\kappa =\frac{1}{4},\sigma =\frac{1}{2}.$$
This case is ruled out in Lemma 7.4 below, concluding the proof of Theorem 3.1 in the case $`w_0=\pm 1`$.
The final case to consider is $`w_0=0`$. In this case, the expansions corresponding to (76)โ(78) are
$`A`$-eqn: $`O(u^{s1})=\mathrm{\hspace{0.33em}1}{\displaystyle \frac{1}{e^2\rho ^2}}+O(u^{2\kappa })+O(u^{s+2\kappa 2})`$ (91)
$`AT^2`$-eqn: $`u^0=O(u^{\kappa +\sigma \frac{3s}{2}})+O(u^{1\frac{3s}{2}})+O(u^{2\kappa 2})`$ (92)
$`Aw^{}`$-eqn: $`O(u^{s+\kappa 2})=O(u^\kappa )+O(u^{\sigma \frac{s}{2}}).`$ (93)
If $`s+2\kappa 2=0`$, we obtain exactly as in the case $`w_00,\pm 1`$ above that $`\sigma =s1`$. It follows that $`\kappa >\frac{s}{2}`$, and Lemma 7.1 yields either the exceptional case of Theorem 3.1, or $`s=2`$, contradicting (61). If on the other hand $`s+2\kappa 2>0`$, we can in (91) use the inequality $`s+2\kappa 2<2\kappa `$ to conclude that $`s1=s+2\kappa 2`$ and thus $`\kappa =\frac{1}{2}`$. Now $`\kappa <\frac{s}{2}`$, and Lemma 7.1 together with (93) yields the two cases
$`s`$ $`=`$ $`{\displaystyle \frac{5}{4}},\kappa ={\displaystyle \frac{1}{2}},\sigma ={\displaystyle \frac{3}{8}}\text{and}`$
$`s`$ $`=`$ $`{\displaystyle \frac{8}{5}},\kappa ={\displaystyle \frac{1}{2}},\sigma ={\displaystyle \frac{9}{10}}.`$
The first case is ruled out in Lemma 7.3 below, whereas the second case leads to a contradiction in (92). This concludes the proof of Theorem 3.1, except for the special cases treated in the following two lemmas.
###### Lemma 7.3.
There is no solution of the EDYM equations satisfying the power ansatz (29), (30), and (62) with
$$s=\frac{5}{4},\kappa =\frac{1}{2},\sigma =\frac{3}{8}.$$
Proof. Suppose that there is a solution of the EDYM equations with
$`A(r)`$ $`=`$ $`A_0u^{\frac{5}{4}}+o(u^{\frac{5}{4}})`$
$`w(r)`$ $`=`$ $`w_1u^{\frac{1}{2}}+o(u^{\frac{1}{2}})`$
with parameters $`A_0,w_10`$. Consider the $`A`$-equation (54). The left side is of the order $`(r\rho )^{\frac{1}{4}}`$. Thus the constant terms on the right side must cancel each other. Then the right side is also of the order $`u^{\frac{1}{4}}`$. Comparing the coefficients gives
$$\frac{5}{4}\rho A_0=\frac{1}{2e^2}A_0w_1^2.$$
This equation yields a contradiction because both sides have opposite sign.
###### Lemma 7.4.
There is no solution of the EDYM equations satisfying the power ansatz (29), (30), and (62) with
$$s=\frac{3}{2},\kappa =\frac{1}{4},\sigma =\frac{1}{2}.$$
(94)
Proof. According to (29), we can write the function $`\sqrt{A}`$ as
$$\sqrt{A}=u^{\frac{3}{4}}+f(u)\text{with}f=o(u^{\frac{3}{4}}).$$
(95)
Employing the ansatz (29),(30) into the A-equation (54), one sees that
$$Aw_{}^{}{}_{}{}^{2}=u^0+u^{\frac{1}{4}}+o(u^{\frac{1}{4}}).$$
(96)
We solve for $`(w^{})^1`$ and substitute (95) to obtain
$`{\displaystyle \frac{1}{w^{}}}=u^{\frac{3}{4}}+u+c_1f+o(u)`$ (97)
with a real constant $`c_1`$. Now consider the $`AT^2`$-equation (48), which we write again in the form (74) and multiply by $`A`$,
$$rA(AT^2)^{}=(AT^2)^{\frac{3}{2}}A^{\frac{1}{2}}[]\frac{4}{e^2}(AT^2)(Aw_{}^{}{}_{}{}^{2}).$$
(98)
As in the proof of Proposition 7.2, a good expansion for the square bracket is obtained by integrating its derivative. Namely, according to (60),
$$[]^{}=\frac{4}{r}w^{}\alpha \beta +u^0+o(u^0)$$
and thus
$$A^{\frac{1}{2}}(r)[]=\frac{4}{r}A^{\frac{1}{2}}_\rho ^r(w^{}\alpha \beta )(s)๐s+u^{\frac{1}{4}}+o(u^{\frac{1}{4}}).$$
Substituting into (98) and using $`๐_๐`$ and (96), we obtain
$$A^{\frac{1}{2}}(r)_\rho ^rw^{}\alpha \beta =u^0+u^{\frac{1}{4}}+o(u^{\frac{1}{4}}).$$
We multiply by $`\sqrt{A}`$, substitute (95) and differentiate,
$$w^{}\alpha \beta =u^{\frac{1}{4}}+u^0+c_2f^{}+o(u^0).$$
Multiplying by $`1/w^{}`$ and using (97) gives the following expansion for $`\alpha \beta `$,
$$\alpha \beta =u^{\frac{1}{2}}+u^{\frac{3}{4}}+c_3u^{\frac{3}{4}}f^{}+c_4u^{\frac{1}{4}}f+o(u^{\frac{3}{4}}).$$
(99)
Next we multiply the $`Aw^{}`$-equation (16) by $`\sqrt{A}`$ and write it as
$$r^2\sqrt{A}(\sqrt{A}(\sqrt{A}w^{}))^{}=e^2r(\sqrt{A}T)(\alpha \beta )+u^{\frac{3}{4}}+o(u^{\frac{3}{4}}).$$
We apply $`๐_๐`$ and substitute (95), (96), and (99). This gives an equation of the form (modulo higher order terms),
$$u^{\frac{1}{2}}+u^{\frac{3}{4}}+u^{\frac{3}{4}}f^{}+uf^{}+u^{\frac{1}{4}}f+u^0=u^{\frac{1}{2}}+u^{\frac{3}{4}}+u^{\frac{3}{4}}f^{}+u^{\frac{1}{4}}f.$$
The constant term $`u^0`$ must vanish since all the other terms tend to zero as $`u0`$. Furthermore, the $`u^{\frac{1}{2}}`$ terms must cancel because all the other terms are $`o(u^{\frac{1}{2}})`$. We thus obtain
$$u^{\frac{3}{4}}f^{}+u^{\frac{1}{4}}f=u^{\frac{3}{4}}+uf^{}+f,$$
so that
$$uf^{}+f=u+u^{\frac{5}{4}}f^{}+u^{\frac{1}{4}}f=u+o(u),$$
and we find that $`f`$ satisfies an equation of the form
$$d_1uf^{}+d_2f=d_3u+o(u).$$
A straightforward but tedious calculation yields that the coefficients $`d_1`$ and $`d_2`$ both vanish, and that $`d_3`$ is non-zero. This is a contradiction.
## 8 The Exceptional Case
In this section, we consider the exceptional case
$$s=\frac{4}{3},\kappa =\frac{1}{3},\sigma =\frac{1}{3}.$$
By employing the power ansatz (29), (30), and (62) into the EDYM equations and comparing coefficients (using Mathematica), we find that the solution near the event horizon is determined by the five free parameters $`(\beta _0,\gamma _0,m,c,\rho )`$. The remaining parameters are given by
$`\delta _0`$ $`=`$ $`\sqrt{\gamma _0^2\beta _0^2+{\displaystyle \frac{2c}{rm}}\beta _0\gamma _0}`$
$`w_0`$ $`=`$ $`{\displaystyle \frac{c\delta _0}{\beta _0}}`$
$`A_0`$ $`=`$ $`{\displaystyle \frac{9\beta _0}{r}}\sqrt{{\displaystyle \frac{m^2r^2\beta _0^2\mathrm{\hspace{0.25em}2}cmr\beta _0\gamma _0+c^2\gamma _0^2}{r^2(1w_0^2)^2}}}`$
$`\alpha _0`$ $`=`$ $`3{\displaystyle \frac{mr\beta _0c\gamma _0}{r\sqrt{A_0}}}`$
$`w_1`$ $`=`$ $`{\displaystyle \frac{9\alpha _0\beta _0}{2rA_0}}`$
Expanding to higher order, we obtain after an arduous calculation two further constraints on the free parameters, thus reducing the problem to one involving only three parameters. We investigated this three-parameter space numerically starting at $`r_0=\rho +\epsilon `$ and found strong evidence that no global black hole solutions exist. Indeed, either the power ansatz was inconsistent near the event horizon (that is, for $`r`$ close to $`r_0+\epsilon `$ the numerical solution deviated from the power ansatz, and became singular as $`\epsilon 0`$), or else the solutions developed a singularity for finite $`r`$.
Acknowledgments: JS would like to thank the Max Planck Institute for Mathematics in the Sciences, Leipzig, and in particular Professor E. Zeidler, for their hospitality and generous support. FF and JS are grateful to the Harvard University Mathematics Department for support.
The research of JS and STY was supported in part by the NSF, Grant No. DMS-G-9802370 and 33-585-7510-2-30.
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# Higher regularity properties of mappings and morphisms
## 1 Noninvertible mappings and morphisms
Here we introduce general mappings (and later morphisms) without usual requirement of โinvertibilityโ. We start from some standard well-known facts (see e.g. ). Let $`X`$ and $`Y`$ be two arbitrary sets. A mapping $`f`$ from $`X`$ to $`Y`$ is defined by a prescription which assigns an element of $`Y`$ to each element of $`X`$, i.e. $`f:XY`$. Injective mapping (injection) assigns different images to different elements, and in surjective mapping (surjection) every image has at least one pre-image. Bijection has both properties. Usually inverse mapping $`f^1`$ is defined as a new mapping $`g:YX`$ which assigns to each $`yY`$ such $`xX`$ that $`f\left(x\right)=y`$ and so $`f^1=g`$. For injective $`f`$ and any $`AX`$ it is imposed the following โinvertibilityโ condition
$$f^1\left(f\left(A\right)\right)=A.$$
(1)
For surjective $`f`$ and $`BY`$ the standard โinvertibilityโ condition is
$$f\left(f^1\left(B\right)\right)=B.$$
(2)
These conditions are very strong, because they imply possibility to solve the equation $`f\left(x\right)=y`$ for all elements. This situation is very much artificial and takes place in only obvious special cases, when functions are defined by invertible operations. Therefore abstract consideration of quantum field theory constructions is usually restricted by group theory methods. But in many cases, especially while considering supersymmetric theories, there naturally appear noninvertible morphisms and semigroups . That obviously needs extending some first principles and assumptions.
We propose to extend the โinvertibilityโ conditions (1)โ(2) in the following way (which comes from analogy of regularity in semigroup theory ). Instead of $`f^1`$ mapping we introduce less restricted โregularโ $`f^{}`$ mapping by extending โinvertibilityโ to โregularityโ in following way
$$f\left(f^{}\left(f\left(A\right)\right)\right)=f\left(A\right).$$
(3)
For the second equation (2) we have the โreflexive regularityโ condition
$$f^{}\left(f\left(f^{}\left(B\right)\right)\right)=f^{}\left(B\right).$$
(4)
## 2 Invertibility and regularity of morphisms
Among all mappings we distinguish morphisms satisfying closure and associativity, which defines a category $``$ with objects $`Ob`$ as sets $`X,Y,Z`$ and morphisms $`Mor`$ as mappings $`f:XY`$ between them (or $`f=Mor(X,Y)`$). For composition $`h:XZ`$ of morphisms $`f:XY`$ and $`g:YZ`$ instead of $`h\left(x\right)=g\left(f\left(x\right)\right)`$ for mappings we use the notation $`h=gf`$. Associativity implies that $`h\left(gf\right)=\left(hg\right)f=hgf`$. Right cancellative morphisms are epimorphisms which satisfy $`g_1f=g_2f`$ $`g_1=g_2`$, where $`g_{1,.2}:YZ`$ and left cancellative morphisms are monomorphisms which satisfy $`fh_1=fh_2`$ $`h_1=h_2`$, where $`h_{1,.2}:ZX`$.
Let us consider โinvertibilityโ properties of morphisms in general. If $`f`$ satisfies the โright invertibilityโ condition $`ff^1=Id_Y`$ for some $`f^1:YX`$ then $`f`$ is called a retraction, and if $`f`$ satisfies the โleft invertibilityโ condition $`f^1f=Id_X`$ , then it is called a coretraction, where $`Id_X`$ and $`Id_Y`$ are identity mappings $`Id_X:XX`$ and $`Id_Y:YY`$ for which $`xX,`$ $`Id_X\left(x\right)=x`$ and $`yY,`$ $`Id_Y\left(y\right)=y`$. These requirements sometimes are very strong to be fulfilled (see e.g. ). To obtain more weak conditions one has to introduced the following โregularityโ conditions
$$ff_{in}^{}f=f,$$
(5)
where $`f_{in}^{}`$ is called an inner inverse , and such $`f`$ is called regular. Similar โreflexive regularityโ conditions
$$f_{out}^{}ff_{out}^{}=f_{out}^{}$$
(6)
defines an outer inverse $`f_{out}^{}`$. Notice that in general $`f_{in}^{}f_{out}^{}f^1`$ or it can be that $`f^1`$ does not exist at all. If $`f_{in}^{}`$ is an inner inverse, then
$$f^{}=f_{in}^{}ff_{in}^{}$$
(7)
is always both inner and outer inverse or generalized inverse (quasi-inverse) , and so for any regular $`f`$ there exists (need not be unique) $`f^{}`$ from (7) for which both regularity conditions (5) and (6) hold.
Let us consider a composition of two morphisms and its โinvertibilityโ properties. It can be shown, that a retraction is an epimorphism, and a regular monomorphism is a coretraction . If $`f`$ and $`g`$ are regular, it is not necessary that their product $`fg`$ is also regular, but, if $`fg`$ is regular and $`g`$ is an epimorphism, then $`f`$ is also regular, and opposite, if $`fg`$ is and $`f`$ is a monomorphism, then $`g`$ is regular.
If the composition $`h=gf`$ belongs to the same class of functions (closure), then all such morphisms form a semigroup of such functions . If for any $`f:XY`$ there will be a unique $`f^{}:YX`$ satisfying (5)โ(6), this semigroup is called an inverse semigroup which we denote $`๐`$.
Let us define two idempotent โprojection operatorsโ $`๐ซ_f=ff^{}`$, $`๐ซ_f:YY`$ and $`๐ซ_f^{}=f^{}f`$, $`๐ซ_f^{}:XX`$ satisfying $`๐ซ_f๐ซ_f=๐ซ_f`$, $`๐ซ_ff=f๐ซ_f^{}=f`$ and $`๐ซ_f^{}๐ซ_f^{}=๐ซ_f^{}`$, $`๐ซ_f^{}f^{}=f^{}๐ซ_f=f^{}`$. If we introduce the $``$-operation $`\left(f\right)^{}=f^{}`$ by formulas (5 )โ(6) and assume that this operation acts on the product of morphisms $`f:XY`$ and $`g:YZ`$ in the following way $`\left(gf\right)^{}=f^{}g^{}`$, then commutativity of projectors $`๐ซ_f๐ซ_g^{}=๐ซ_g^{}๐ซ_f`$ leads to closure of the semigroup, i.e. the product $`gf`$ also satisfies both regularity conditions (5)โ(6).
## 3 Higher analogs of regularity
Here we try to construct higher analogs of regularity conditions (5 )โ(6). They were proposed for some particular case (noninvertible analogs of supermanifolds) in (for other generalizations of regularity see e.g. ).
Let we have two elements $`f`$ and its regular $`f^{}`$ (in sense of (5)) of semigroup $`๐`$. Consider a third morphism $`f^{}:XY`$ and analyze the action $`ff^{}f^{}:XY`$. This means the composition $`ff^{}f^{}`$ cannot be equal to identity $`Id_X`$. Therefore it is possible to โregularizeโ $`ff^{}f^{}`$ in the following way
$$ff^{}f^{}f^{}=f^{}.$$
(8)
This formula can be called as 2-regularity condition and be considered as a definition of $``$-operation. For 3-regularity and $`f^{}:YX`$ we can obtain an analog of (5) in the form
$$ff^{}f^{}f^{}f=f.$$
(9)
By recursive considerations we can propose the following formula of $`n`$ -regularity
$`ff^{}f^{}\mathrm{}f^{\stackrel{2k}{\stackrel{}{\mathrm{}}}}f^{}`$ $`=f^{}`$ (10)
$`ff^{}f^{}\mathrm{}f^{\stackrel{2k+1}{\stackrel{}{\mathrm{}}}}f`$ $`=f.`$ (11)
Note that for even number of stars $`f^{\stackrel{2k}{\stackrel{}{\mathrm{}}}}:XY`$ and for odd number of stars $`f^{\stackrel{2k+1}{\stackrel{}{\mathrm{}}}}:YX`$. In general case all $`f^\stackrel{๐}{\stackrel{}{\mathrm{}}}`$ are different, and, for instance, $`\left(f^{}\right)^{}f^{}`$. The action of $`\stackrel{๐}{\stackrel{}{\mathrm{}}}`$-operation on product depends from number of stars and is the following
$`\left(gf\right)^{\stackrel{2k+1}{\stackrel{}{\mathrm{}}}}`$ $`=f^{\stackrel{2k+1}{\stackrel{}{\mathrm{}}}}g^{\stackrel{2k+1}{\stackrel{}{\mathrm{}}}},`$ (12)
$`\left(gf\right)^{\stackrel{2k}{\stackrel{}{\mathrm{}}}}`$ $`=g^{\stackrel{2k}{\stackrel{}{\mathrm{}}}}f^{\stackrel{2k}{\stackrel{}{\mathrm{}}}}.`$ (13)
We can introduce โhigher projectorโ by the formula
$$๐ซ_f^{\left(n\right)}=ff^{}f^{}\mathrm{}f^\stackrel{๐}{\stackrel{}{\mathrm{}}}.$$
(14)
It is easy to check the following properties
$`๐ซ_f^{\left(2k\right)}f^{}`$ $`=f^{},`$ (15)
$`๐ซ_f^{\left(2k+1\right)}f`$ $`=f.`$ (16)
and idempotence $`๐ซ_f^{\left(n\right)}๐ซ_f^{\left(n\right)}=๐ซ_f^{\left(n\right)}`$.
## 4 Semicommutative diagrams and obstruction
In previous subsection we considered morphisms and regularity properties for two given objects $`X`$ and $`Y`$, because we studied various types of inverses. Now we will extend these consideration to any number of objects and introduce semicommutative diagrams (firstly considered in ).
Obviously, that for two morphisms $`f:XY`$ and $`g:YX`$ instead of โinvertibilityโ $`gf=Id_X`$ we have the same generalization as regularity (5), i.e. $`fgf=f`$, where $`g`$ plays the role of an inner inverse .
Usually, for 3 objects $`X,Y,Z`$ and 3 morphisms $`f:XY`$ and $`g:YZ`$ and $`h:ZX`$ one can have the โinvertibleโ triangle commutative diagram $`hgf=Id_X`$. Its regular extension has the form
$$fhgf=f.$$
(17)
Such a diagram (from the right)
can be called a semicommutative diagram. This triangle case can be expanded on any number of objects and morphisms.
To measure difference between semicommutative and commutative cases let us introduce self-mappings $`e_X^{\left(n\right)}:XX`$ which are defined by
$`e_X^{\left(1\right)}`$ $`=Id_X,`$ (18)
$`e_X^{\left(2\right)}`$ $`=gf,`$ (19)
$`e_X^{\left(3\right)}`$ $`=hgf,`$ (20)
$`\mathrm{}`$
It is obvious that for commutative diagrams all $`e_X^{\left(n\right)}`$ are equal to identity $`e_X^{\left(n\right)}=Id_X`$. The deviation of $`e_X^{\left(n\right)}`$ from identity will give us measure of obstruction of commutativity, and therefore we call $`e_X^{\left(n\right)}`$ obstructors. The minimum number $`n=n_{obstr}`$ for which $`e_X^{\left(n\right)}Id_X`$ occurs will define a quantitative measure of obstruction $`n_{obstr}`$.
In terms of obstructors $`e_X^{\left(n\right)}`$ the $`n`$-regularity condition can be written in the short form
$$fe_X^{\left(n\right)}=f.$$
(21)
From definitions (18)โ(20) and (21) it simply follows that obstructors $`e_X^{\left(n\right)}`$ are idempotents.
It can be noted that for a given $`n`$ additional arrows do not lead to new commutativity relations due to (21). Therefore only one additional arrow can be important for extension of the commutativity to semicommutativity.
## 5 Monoidal categories, functors and regularity
The morphisms $`e_X^{\left(n\right)}`$ can be used to extend the notion of a functor $`๐ฅ:_1_2`$. All the standard definitions of a functor (as a mapping of one category to another with preserving composition of morphisms ) do not changed, but preservation of identity $`๐ฅ\left(Id_{X_1}\right)=Id_{X_2}`$, where $`X_2=๐ฅX_1`$, $`X_1\mathrm{Ob}_1`$, $`X_2\mathrm{Ob}_2`$, can be replaced by requirement of preservation of morphisms $`e_X^{\left(n\right)}`$ as
$$๐ฅ^{\left(n\right)}\left(e_{X_1}^{\left(n\right)}\right)=e_{X_2}^{\left(n\right)},$$
(22)
where $`e_{X_1}^{\left(n\right)}\mathrm{Mor}_1`$, $`e_{X_2}^{\left(n\right)}\mathrm{Mor}_2`$ defined in (18 )โ(20) for two categories. Then the generalized functor $`๐ฅ^{\left(n\right)}`$ becomes $`n`$-dependent. From (18) it follows that $`n=1`$ corresponds to the standard functor, i.e. $`๐ฅ^{\left(1\right)}=๐ฅ`$. In the same manner we can โregularizeโ natural transformations .
Let $``$ be a monoidal category equipped with a monoidal operation $`:\times `$. A triple of objects $`X,Y,Z`$ is said to be a regular $`3`$-cycle if and only if every sequence of morphisms $`X\stackrel{๐}{}Y\stackrel{๐}{}Z\stackrel{}{}X`$ define uniquely the morphism $`e_X^{\left(3\right)}:XX`$ by the following relation $`e_X^{\left(3\right)}:=hgf`$ and subjects to the relation $`fhgf=f`$. The object $`Y`$ is said to be a (first) regular dual of $`X`$, and the object $`Z`$ is called the second regular dual of $`X`$. We denote by $`C_3\left(\right)`$ the collection of all regular $`3`$-cycles on $``$ . This collection is said to be regularity in $``$. The generalization to arbitrary $`n4`$ is obvious. Let $`X,Y,Z`$ and $`X^{},Y^{},Z^{}`$ be two regular $`3`$-cycles in $``$. Then the morphism $`f:XX^{}`$ such that $`fe_X^{\left(3\right)}=e_X^{}^{\left(3\right)}f`$ is said to be a 3-cycle morphism. If $`f:XX^{}`$ and $`g:X^{}X^{\prime \prime }`$ are two 3-cycle morphisms, then the composition $`gf:XX^{\prime \prime }`$ is also a 3-cycle morphism. Moreover the regularity $`C_3\left(\right)`$ forms a monoidal category with 3-cycles as objects and 3-cycle morphisms. The monoidal product of two regular 3-triples $`X,Y,Z`$ and $`X^{},Y^{},Z^{}`$ is the triple $`XX^{},YY^{},ZZ^{}`$ which is also a regular $`3`$-cycle. The category $`C_3\left(\right)`$ is said to be a regularization of $``$ .
Let us consider a symmetric monoidal category $``$ playing an important role in topological QFT , supersymmetry , quantum groups and quantum statistics .In $``$ for any two objects $`X`$ and $`Y`$ and the operation $`XY`$ one usually defines a natural isomorphism (โbraidingโ ) by $`\mathrm{B}_{X,Y}:XYYX`$ satisfying the symmetry condition (โinvertibilityโ)
$$\mathrm{B}_{Y,X}\mathrm{B}_{X,Y}=Id_{XY}$$
(23)
which formally defines $`\mathrm{B}_{Y,X}=\mathrm{B}_{X,Y}^1:YXXY`$. Note that possible nonsymmetric braiding in context of the noncommutative geometry was considered in (see also ). Here we introduce a โregularโ extension of the symmetry condition (23) in the form
$$\mathrm{B}_{X,Y}\mathrm{B}_{X,Y}^{}\mathrm{B}_{X,Y}=\mathrm{B}_{X,Y},$$
(24)
where in general $`\mathrm{B}_{X,Y}^{}\mathrm{B}_{X,Y}^1`$. Such a category can be called a โregularโ category to distinct from symmetric and โbraidedโ categories .
In categorical sense the prebraiding relations usually are defined as
$`\mathrm{B}_{XY,Z}`$ $`=๐_{X,Z,Y}^R๐_{X,Y,Z}^L,`$ (25)
$`\mathrm{B}_{Z,XY}`$ $`=๐_{X,Z,Y}^L๐_{X,Y,Z}^R,`$ (26)
$`๐_{X,Y,Z}^L`$ $`=Id_X\mathrm{B}_{Y,Z},`$ (27)
$`๐_{X,Y,Z}^R`$ $`=\mathrm{B}_{X,Y}Id_Z,`$ (28)
and prebraidings $`\mathrm{B}_{XY,Z}`$ and $`\mathrm{B}_{Z,XY}`$ satisfy (for symmetric case) the โinvertibilityโ property $`\mathrm{B}_{XY,Z}^1\mathrm{B}_{XY,Z}=Id_{XYZ}`$ , where $`\mathrm{B}_{XY,Z}^1=\mathrm{B}_{Z,XY}`$. In this notations the standard โinvertibleโ Yang-Baxter equation is
$$๐_{Y,Z,X}^R๐_{Y,X,Z}^L๐_{X,Y,Z}^R=๐_{Z,X,Y}^L๐_{X,Z,Y}^R๐_{X,Y,Z}^L.$$
(29)
Possible โnoninvertibleโ (endomorphism semigroup) solutions of this equation without introduction of $`e_X^{\left(n\right)}`$ were studied in . For โnoninvertibleโ braidings satisfying regularity (24) it is naturally to exploit the obstructors $`e_X^{\left(n\right)}`$ instead of identity $`Id_X`$ as
$`๐_{X,Y,Z}^{L\left(n\right)}`$ $`=e_X^{\left(n\right)}\mathrm{B}_{Y,Z},`$ (30)
$`๐_{X,Y,Z}^{R\left(n\right)}`$ $`=\mathrm{B}_{X,Y}e_Z^{\left(n\right)},`$ (31)
to weaken prebraiding construction in the following way
$`\mathrm{B}_{XY,Z}^{\left(n\right)}`$ $`=๐_{X,Z,Y}^{R\left(n\right)}๐_{X,Y,Z}^{L\left(n\right)},`$ (32)
$`\mathrm{B}_{Z,XY}^{\left(n\right)}`$ $`=๐_{X,Z,Y}^{L\left(n\right)}๐_{X,Y,Z}^{R\left(n\right)},`$ (33)
Then their โinvertibilityโ can be also โregularizedโ as follows
$$\mathrm{B}_{XY,Z}^{\left(n\right)}\mathrm{B}_{XY,Z}^{\left(n\right)}\mathrm{B}_{XY,Z}^{\left(n\right)}=\mathrm{B}_{XY,Z}^{\left(n\right)},$$
(34)
where in general case $`\mathrm{B}_{XY,Z}^{\left(n\right)}\mathrm{B}_{XY,Z}^1`$. Thus the corresponding $`n`$ -โnoninvertibleโ analog of the Yang-Baxter equation (33) is
$$๐_{Y,Z,X}^{R\left(n\right)}๐_{Y,X,Z}^{L\left(n\right)}๐_{X,Y,Z}^{R\left(n\right)}=๐_{Z,X,Y}^{L\left(n\right)}๐_{X,Z,Y}^{R\left(n\right)}๐_{X,Y,Z}^{L\left(n\right)}.$$
(35)
Its solutions can be found by application of the semigroup methods (see e.g. ).
The introduced formalism can be used in analysis of categories with some weaken invertibility conditions, which can appear in nontrivial supersymmetric or noncommutative geometry constructions beyond the group theory.
Acknowledgments. One of the authors (S.D.) would like to thank Jerzy Lukierski for kind hospitality at the University of Wrocลaw, where this work was done, and to Andrzej Borowiec, Andrzej Frydryszak and Cezary Juszczak for valuable discussions and help during his stay in Wrocลaw.
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# Gap and subgap tunnelling in cuprates
## Abstract
We describe strongly attractive carriers in cuprates in the framework of a simple quasi-one dimensional Hamiltonian with a local attraction. In contrast with the conventional BCS theory there are two energy scales, a temperature independent incoherent gap $`\mathrm{\Delta }_p`$ and a temperature dependent coherent gap $`\mathrm{\Delta }_c(T)`$ combining into one temperature dependent global gap $`\mathrm{\Delta }=(\mathrm{\Delta }_p^2+\mathrm{\Delta }_c^2)^{1/2}`$. The temperature dependence of the gap and single particle (Giaver) tunnelling spectra in cuprates are quantitatively described. A framework for understanding of two distinct energy scales observed in Giaver tunnelling and electron-hole reflection experiments is provided.
There is convincing experimental evidence that the pairing of carriers takes place well above T<sub>c</sub> in underdoped cuprates (for a review see Ref. ). If carriers are paired their magnetic moments compensate each other so one could expect that the normal state uniform magnetization should fall with decreasing temperature because more and more holes are bound into singlet pairs. This unexpected drop of the normal state magnetic susceptibility was experimentally observed and explained in the framework of the bipolaron theory of cuprates . There is also a gap in the tunnelling and photoemission, which is almost temperature independent below T<sub>c</sub> and exists well above T<sub>c</sub> , so that some segments of a โlarge Fermi surfaceโ are actually missing . Kinetic and thermodynamic data suggest that the gap opens in both charge and spin channels and exists at any relevant temperature in a wide range of doping. A plausible explanation is that the normal gap is half of the bipolaron binding energy , although alternative explanations have also been proposed. The temperature and doping dependence of the gap still remains a subject of controversy. Moreover, reflection experiments, in which an incoming electron from the normal side of a normal/superconducting contact is reflected as a hole along the same trajectory , revealed a much smaller gap edge than the bias at the tunnelling conductance maxima in a few underdoped cuprates . Recent intrinsic tunnelling measurements on a series of Bi โ2212โ single crystals showed distinctly different behaviour of the superconducting and normal state gaps with the magnetic field. Such coexistance of two distinct gaps in cuprates is not well understood .
In this letter we propose a model, which describes the temperature dependence of the gap, tunnelling spectra and electron-hole reflection in cuprates. The assumption is that the attraction potential in cuprates is large compared with the Fermi energy. The main point of our letter is independent of the microscopic nature of the attraction. Real-space pairs might be lattice and/or spin bipolarons , or any other preformed pairs.
We start with a generic one-dimensional Hamiltonian including the kinetic energy of carriers in the effective mass ($`m`$) approximation and a local attraction potential, $`V(xx^{})=U\delta (xx^{})`$ as
$`H`$ $`=`$ $`{\displaystyle \underset{s}{}}{\displaystyle ๐x\psi _s^{}(x)\left(\frac{1}{2m}\frac{d^2}{dx^2}\mu \right)\psi _s(x)}`$ (1)
$``$ $`U{\displaystyle ๐x\psi _{}^{}(x)\psi _{}^{}(x)\psi _{}(x)\psi _{}(x)},`$ (2)
where $`s=,`$ is the spin ($`\mathrm{}=k_B=1`$). The first band to be doped in cuprates is the oxygen band inside the Hubbard gap as established in polarised photoemission . This band is almost one dimensional as discussed in Ref. , so that our (quasi) one-dimensional approximation is a realistic starting point.
Solving a two-particle problem with the $`\delta `$-function potential one obtains a bound state with the binding energy
$$2\mathrm{\Delta }_p=\frac{1}{4}mU^2,$$
(3)
and with the radius of the bound state $`r=2/mU`$. We assume that this radius is less than the inter-carrier distance in cuprates, which puts a constraint on the doping level, $`E_F<2\mathrm{\Delta }_p`$, where $`E_F`$ is the free-carrier Fermi energy. Then real-space pairs are formed. If three-dimensional corrections to the energy spectrum of pairs are taken into account (see, for example, Ref. ) the ground state of the system is the Bose-Einstein condensate. The chemical potential is pinned below the band edge by about $`\mathrm{\Delta }_p`$ both in the superconducting and normal state , so that the normal state single-particle gap is $`\mathrm{\Delta }_p`$. The binding energy $`2\mathrm{\Delta }_p`$ might change due to the same corrections. However, this change does not affect our further results as soon as they are expressed in terms of $`\mathrm{\Delta }_p`$ rather than $`U`$.
Now we take into account that in the superconducting state ($`T<T_c`$) the single-particle excitations interact with the condensate via the same potential $`U`$. Applying the Bogoliubov approximation we reduce the Hamiltonian, Eq.(1) to a quadratic form as
$`H`$ $`=`$ $`{\displaystyle \underset{s}{}}{\displaystyle ๐x\psi _s^{}(x)\left(\frac{1}{2m}\frac{d^2}{dx^2}\mu \right)\psi _s(x)}`$ (4)
$`+`$ $`{\displaystyle }dx[\mathrm{\Delta }_c\psi _{}^{}(x)\psi _{}^{}(x)+H.c.],`$ (5)
where the coherent pairing potential
$$\mathrm{\Delta }_c=U\psi _{}(x)\psi _{}(x)$$
(6)
is proportional to the square root of the condensate density, $`\mathrm{\Delta }_c=constant\times n_0(T)^{1/2}`$. The single-particle excitation energy spectrum $`E(k)`$ is found using the Bogoliubov transformation as
$$E(k)=\left[(k^2/2m+\mathrm{\Delta }_p)^2+\mathrm{\Delta }_c^2\right]^{1/2},$$
(7)
if one assumes that the condensate density does not depend on position. This spectrum is quite different from the BCS quasiparticles because the chemical potential is negative, $`\mu =\mathrm{\Delta }_p`$. The single particle gap, $`\mathrm{\Delta }`$, defined as the minimum of $`E(k)`$, is given by
$$\mathrm{\Delta }=\left[\mathrm{\Delta }_p^2+\mathrm{\Delta }_c^2\right]^{1/2}.$$
(8)
It varies with temperature from $`\mathrm{\Delta }(0)=\left[\mathrm{\Delta }_p^2+\mathrm{\Delta }_c(0)^2\right]^{1/2}`$ at zero temperature down to the temperature independent $`\mathrm{\Delta }_p`$ above $`T_c`$. The condensate density depends on temperature as $`1(T/T_c)^{d/2}`$ in the ideal three ($`d=3`$) and (quasi) two-dimensional ($`d=2`$) Bose-gas. In the three-dimensional $`charged`$ Bose-gas it has an exponential temperature dependence at low temperatures due to a plasma gap in the Bogoliubov collective excitation spectrum , which might be highly anisotropic in cuprates . Near T<sub>c</sub> one can expect a power law dependence, $`n_0(T)1(T/T_c)^n`$ with $`n>d/2`$ because the condensate plasmon depends on temperature. The theoretical temperature dependence, Eq.(6) describes well the pioneering experimental observation of the anomalous gap in $`YBa_2Cu_3O_{7\delta }`$ in the electron-energy-loss spectra by Demuth $`et`$ $`al`$ , Fig.1, with $`\mathrm{\Delta }_c(T)^2=\mathrm{\Delta }_c(0)^2\times [1(T/T_c)^n]`$ below $`T_c`$ and zero above $`T_c`$, and $`n=4`$.
The normal metal-superconductor (SIN) tunnelling conductance via a dielectric contact, $`dI/dV`$ is proportional to the density of states, $`\rho (E)`$ of the spectrum Eq.(5). Taking also into account the scattering of single-particle excitations by a random potential, thermal lattice and spin fluctuations one finds at $`T=0`$
$$dI/dV=constant\times [\rho \left(\frac{2eV2\mathrm{\Delta }}{ฯต_0}\right)+A\rho \left(\frac{2eV2\mathrm{\Delta }}{ฯต_0}\right)],$$
(9)
with
$$\rho (\xi )=\frac{4}{\pi ^2}\times \frac{Ai(2\xi )Ai^{}(2\xi )+Bi(2\xi )Bi^{}(2\xi )}{[Ai(2\xi )^2+Bi(2\xi )^2]^2},$$
(10)
$`A`$ is the asymmetry coefficient, $`Ai(x),Bi(x)`$ the Airy functions, and $`ฯต_0`$ is the scattering rate. We compare the conductance, Eq.(7) with one of the best STM spectra measured in $`Ni`$-substituted $`Bi_2Sr_2CaCu_2O_{8+x}`$ single crystals by Hancottee $`et`$ $`al`$, Fig.2. This experiment showed anomalously large $`2\mathrm{\Delta }/T_c>12`$ with the temperature dependence of the gap similar to that in Fig.1. The theoretical conductance, Eq.(7) describes well the anomalous $`gap/T_c`$ ratio, injection/emission assymmetry, zero-bias conductance at zero temperature, and the spectral shape inside and outside the gap region. There is no doubt that the gap, Fig.2 is s-like, which is compatible with the phase-sensitive experiments in the framework of the bipolaron theory. Within the theory the single-particle gap might be almost $`k`$ independent while the symmetry of the Bose-Einstein condensate wave-function (i.e. of the order parameter) is $`dwave`$.
Finally, we propose a simple theory of the tunnelling into bosonic (bipolaronic) superconductor in the metallic (no-barrier) regime. As in the canonical BCS approach applied to the normal metal-superconductor tunnelling by Blonder, Tinkham and Klapwijk and to the normal-superconductor boundary in the intermediate type I state by one of us , the incoming electron produces only outgoing particles in the superconductor ($`x>l`$), allowing for a reflected electron and (Andreev) hole in the normal metal ($`x<0`$). There is also a buffer layer of the thickness $`l`$ at the normal metal-superconductor boundary ( $`x=0`$), where the chemical potential with respect to the bottom of the conduction band changes gradually from a positive large value $`\mu `$ in the metal to a negative value $`\mathrm{\Delta }_p`$ in the bosonic superconductor. We approximate this buffer layer by a layer with a constant chemical potential $`\mu _b`$ ($`\mathrm{\Delta }_p<\mu _b<\mu `$) and with the same strength of the pairing potential $`\mathrm{\Delta }_c`$ as in the bulk superconductor. The Bogoliubov-de Gennes equations may be written as usual , with the only difference that the chemical potential with respect to the bottom of the band is a function of the coordinate $`x`$,
$`\left(\begin{array}{cc}(1/2m)d^2/dx^2\mu (x)& \mathrm{\Delta }_c\\ \mathrm{\Delta }_c& (1/2m)d^2/dx^2+\mu (x)\end{array}\right)\psi (x)`$ (11)
$`=E\psi (x).`$ (12)
Thus the two-componet wave function in the normal metal is given by
$$\psi _n(x<0)=\left(\begin{array}{c}1\\ 0\end{array}\right)e^{iq^+x}+b\left(\begin{array}{c}1\\ 0\end{array}\right)e^{iq^+x}+a\left(\begin{array}{c}0\\ 1\end{array}\right)e^{iq^{}x},$$
(13)
while in the buffer layer it has the form
$`\psi _b(0<x<l)`$ $`=`$ $`\alpha \left(\begin{array}{c}1\\ \frac{\mathrm{\Delta }_c}{E+\xi }\end{array}\right)e^{ip^+x}+\beta \left(\begin{array}{c}1\\ \frac{\mathrm{\Delta }_c}{E\xi }\end{array}\right)e^{ip^{}x}`$ (14)
$`+`$ $`\gamma \left(\begin{array}{c}1\\ \frac{\mathrm{\Delta }_c}{E+\xi }\end{array}\right)e^{ip^+x}+\delta \left(\begin{array}{c}1\\ \frac{\mathrm{\Delta }_c}{E\xi }\end{array}\right)e^{ip^{}x},`$ (15)
where the momenta associated with the energy $`E`$ are
$$q^\pm =[2m(\mu \pm E)]^{1/2}$$
(16)
and
$$p^\pm =[2m(\mu _b\pm \xi )]^{1/2}$$
(17)
with $`\xi =(E^2\mathrm{\Delta }_c^2)^{1/2}`$. The well-behaved solution in the superconductor with negative chemical potential is given by
$$\psi _s(x>l)=c\left(\begin{array}{c}1\\ \frac{\mathrm{\Delta }_c}{E+\xi }\end{array}\right)e^{ik^+x}+d\left(\begin{array}{c}1\\ \frac{\mathrm{\Delta }_c}{E\xi }\end{array}\right)e^{ik^{}x},$$
(18)
where the momenta associated with the energy $`E`$ are
$$k^\pm =[2m(\mathrm{\Delta }_p\pm \xi )]^{1/2}.$$
(19)
The coefficients $`a,b,c,d,\alpha ,\beta ,\gamma ,\delta `$ are determined from the boundary conditions, which are continuity of $`\psi (x)`$ and its derivatives at $`x=0`$ and $`x=l`$. Applying the boundary conditions, and carrying out an algebraic reduction, we find
$$a=2\mathrm{\Delta }_cq^+(p^+f^{}g^+p^{}f^+g^{})/D,$$
(20)
$`b`$ $`=`$ $`1+2q^+[(E+\xi )f^+(q^{}f^{}p^{}g^{})`$ (21)
$``$ $`(E\xi )f^{}(q^{}f^+p^+g^+)]/D,`$ (22)
with
$`D`$ $`=`$ $`(E+\xi )(q^+f^++p^+g^+)(q^{}f^{}p^{}g^{})`$ (23)
$``$ $`(E\xi )(q^+f^{}+p^{}g^{})(q^{}f^+p^+g^+),`$ (24)
and $`f^\pm =p^\pm \mathrm{cos}(p^\pm l)ik^\pm \mathrm{sin}(p^\pm l)`$, $`g^\pm =k^\pm \mathrm{cos}(p^\pm l)ip^\pm \mathrm{sin}(p^\pm l)`$.
The transmisson coefficient for electrical current, $`1+|a|^2|b|^2`$ is shown in Fig.3 for different values of $`l`$ when the coherent gap $`\mathrm{\Delta }_c`$ is smaller than the pair-breaking gap $`\mathrm{\Delta }_p`$, and in Fig.4 for the opposite case, $`\mathrm{\Delta }_p<\mathrm{\Delta }_c`$. In the first case, Fig.3, we find two distinct energy scales, one is $`\mathrm{\Delta }_c`$ in the subgap region due to electron-hole reflection and the other one is $`\mathrm{\Delta }`$, which is the single-particle band edge. On the other hand, there is only one gap $`\mathrm{\Delta }_c`$, which can be seen in the second case, Fig.4. We notice that the transmission has no subgap structure if the buffer layer is absent ($`l=0`$) in both cases. In the extreme case of a wide buffer layer, $`l>>(2m\mathrm{\Delta }_p)^{1/2}`$, Fig.3, or $`l>>(2m\mathrm{\Delta }_c)^{1/2}`$, Fig.4, there are some oscillations of the transmission due to the bound states inside the buffer layer. It was shown in Ref. that the pair-breaking gap $`\mathrm{\Delta }_p`$ is inverse proportional to the doping level. On the other hand, the coherent gap $`\mathrm{\Delta }_c`$ scales with the condensate density, and therefore with the critical temperature, determined as the Bose-Einstein condensation temperature of strongly anisotropic 3D bosons . Therefore we expect that $`\mathrm{\Delta }_p>>\mathrm{\Delta }_c`$ in the underdoped cuprates, Fig.3, while $`\mathrm{\Delta }_p\mathrm{\Delta }_c`$ in the optimally doped cuprates, Fig.4. Thus the model accounts for the two different gaps experimentally observed in Giaver tunnelling and electron-hole reflection in the underdoped cuprates and for a single gap in the optimally doped samples . The transmission, Fig.3 and Fig.4, is entirely due to the coherent tunnelling into the condensate and (or) into the single-particle band of the bosonic superconductor. There is also an incoherent transmission into localised single-particle impurity states and into incoherent (โsupracondensateโ) bound pair states, which might explain a significant featureless background in the subgap region .
In conclusion, we have proposed a simple general model, which provides an explanation of the temperature dependence of the gap and of the single-particle tunnelling spectra in cuprates. The main assumption is that the attractive potential is large compared with the Fermi energy, so that the ground state is the Bose-Einstein condensate of tightly bound pairs. We have developed a theory of tunnelling in the metallic regime with no barrier and found two different energy scales in the transmission as observed experimentally.
We acknowledge support of this work by the Leverhulme Trust (London), grant VP/261.
Figure Captures
Fig.1. Temperature dependence of the gap, Eq.(6) (line) compared with the experiment (dots) for $`\mathrm{\Delta }_p=0.7\mathrm{\Delta }(0)`$ .
Fig.2. Theoretical tunnelling conductance, Eq.(7) (line) compared with STM conductance in Ni-doped $`Bi_2Sr_2CaCu_2O_{8+x}`$ ) (dots) for $`2\mathrm{\Delta }=90`$ meV, $`A=1.05`$, $`ฯต_0=40`$ meV.
Fig.3. Transmission versus voltage (measured in units of $`\mathrm{\Delta }_p/e`$) for $`\mathrm{\Delta }_c=0.2\mathrm{\Delta }_p`$, $`\mu =10\mathrm{\Delta }_p`$, $`\mu _b=2\mathrm{\Delta }_p`$ and $`l=0`$ (thick line), $`l=1`$ (thick dashed line), $`l=4`$ (thin line), and $`l=8`$ (thin dashed line) (in units of $`1/(2m\mathrm{\Delta }_p)^{1/2}`$).
Fig.4. Transmission versus voltage (measured in units of $`\mathrm{\Delta }_c/e`$) for $`\mathrm{\Delta }_p=0.2\mathrm{\Delta }_c`$, $`\mu =10\mathrm{\Delta }_c`$, $`\mu _b=2\mathrm{\Delta }_c`$ and $`l=0`$ (thick line), $`l=1`$ (thick dashed line), $`l=4`$ (thin line), and $`l=8`$ (thin dashed line) (in units of $`1/(2m\mathrm{\Delta }_c)^{1/2}`$).
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# Geometry of River Networks I: Scaling, Fluctuations, and Deviations
## I Introduction
Networks are intrinsic to a vast number of complex forms observed in the natural and man-made world. Networks repeatedly arise in the distribution and sharing of information, stresses and materials. Complex networks give rise to interesting mathematical and physical properties as observed in the Internet reka99 , the โsmall-worldโ phenomenon watts98 , the cardiovascular system zamir99 , force chains in granular media coppersmith96 , and the wiring of the brain cherniak99 .
Branching, hierarchical geometries make up an important subclass of all networks. Our present investigations concern the paradigmatic example of river networks. The study of river networks, though more general in application, is an integral part of geomorphology, the theory of earth surface processes and form. Furthermore, river networks are held to be natural exemplars of allometry, i.e., how the dimensions of different parts of a structure scale or grow with respect to each other mandelbrot83 ; rodriguez-iturbe97 ; rinaldo98 ; ball99 ; dodds2000pa . The shapes of drainage basins, for example, are reported to elongate with increasing basin size hack57 ; maritan96a ; rigon96 .
At present, there is no generally accepted theory explaining the origin of this allometric scaling. The fundamental problem is that an equation of motion for erosion, formulated from first principles, is lacking. The situation is somewhat analogous to issues surrounding the description of the dynamics of granular media jaeger96a ; kadanoff99 , noting that erosion is arguably far more complex. Nevertheless, a number of erosion equations have been proposed ranging from deterministic smith72 ; kramer92 ; izumi95 ; sinclair96 to stochastic theories banavar97 ; somfai97 ; pastor-satorras98 ; pastor98b ; cieplak98a ; giacometti2000u . Each of these models attempts to describe how eroding surfaces evolve dynamically. In addition, various heuristic models of both surface and network evolution also exist. Examples include simple lattice-based models of erosion kramer92 ; takayasu92 ; leheny95 ; caldarelli97 , an analogy to invasion percolation stark91 , the use of optimality principles and self-organized criticality rodriguez-iturbe97 ; sun94 ; sun95 , and even uncorrelated random networks dodds2000pa ; leopold62 ; scheidegger67 . Since river networks are an essential feature of eroding landscapes, any appropriate theory of erosoion must yield surfaces with network structures comparable to that of the real world. However, no model of eroding landscapes or even simply of network evolution unambiguously reproduces the wide range of scaling behavior reported for real river networks.
A considerable problem facing these theories and models is that the values of scaling exponents for river network scaling laws are not precisely known. One of the issues we address in this work is universality dodds2000pa ; maritan96a . Do the scaling exponents of all river networks belong to a unique universality class or are there a set of classes obtained for various geomorphological conditions? For example, theoretical models suggest a variety of exponent values for networks that are directed versus non-directed, created on landscapes with heterogeneous versus homogeneous erosivity and so on dodds2000pa ; manna96 ; maritan96b . Clearly, refined measurements of scaling exponents are imperative if we are to be sure of any network belonging to a particular universality class. Moreover, given that there is no accepted theory derivable from simple physics, more detailed phenomenological studies are required.
Motivated by this situation, we perform here a detailed investigation of the scaling properties of river networks. We analytically characterize fluctuations about scaling showing that they grow with system size. We also report significant and ubiquitous deviations from scaling in real river networks. This implies surprisingly strong restrictions on the parameter regimes where scaling holds and cautions against measurements of exponents that ignore such limitations. In the case of the Mississippi basin, for example, we find that although our study region span four orders of magnitude in length, scaling may be deemed valid over no more than 1.5 orders of magnitude. Furthermore, we repeatedly find the scaling within these bounds to be only approximate and that no exact, single exponent can be deduced. We show that scaling breaks down at small scales due to the presence of linear basins and at large scales due to the inherent discreteness of network structure and correlations with overall basin shape. Significantly, this latter correlation imprints upon river network structure the effects and history of geology.
This paper is the first of a series of three on river-network geometry. Having addressed scaling laws in the present work, we proceed in second and third articles dodds2000ub ; dodds2000uc to consider river network structure at a more detailed level. In dodds2000ub we examine the statistics of the โbuilding blocksโ of river networks, i.e., segments of streams and sub-networks. In particular, we analytically connect distributions of various kinds of stream length. Part of this material is employed in the present article and is a direct generalization of Hortonโs laws horton45 ; dodds99pa . In the third article dodds2000uc , we proceed from the findings of dodds2000ub to characterize how these building blocks fit together. Central to this last work is the study of the frequency and spatial distributions of tributary branches along the length of a stream and is itself a generalization of the descriptive picture of Tokunaga tokunaga66 ; tokunaga78 ; tokunaga84 .
## II Basin allometry
### II.1 Hackโs law
In addressing these broader issues of scaling in branching networks, we set as our goal to understand the river network scaling relationship between basin area $`a`$ and the length $`l`$ of a basinโs main stream:
$$la^h.$$
(1)
Known as Hackโs law hack57 , this relation is central to the study of scaling in river networks maritan96a ; dodds99pa . Hackโs exponent $`h`$ is empirically found to lie in the range from 0.5 to 0.7 hack57 ; maritan96a ; rigon96 ; gray61 ; mueller72 ; mosley73 ; mueller73 ; montgomery92 ; rigon98 . Here, we postulate a generalized form of Hackโs law that shows good agreement with data from real world networks.
We focus on Hackโs law because of its intrinsic interest and also because many interrelationships between a large number of scaling laws are known and only a small subset are understood to be independent maritan96a ; dodds99pa . Thus, our results for Hackโs law will be in principle extendable to other scaling laws. With this in mind, we will also discuss probability densities of stream length and drainage area.
Hackโs law is stated rather loosely in equation (1) and implicitly involves some type of averaging which needs to be made explicit. It is most usually considered to be the relationship between mean main stream length and drainage area, i.e.,
$$la^h.$$
(2)
Here, $``$ denotes ensemble average and $`l=l(a)`$ is the mean main stream length of all basins of area $`a`$. Typically, one performs regression analysis on $`\mathrm{log}l`$ against $`\mathrm{log}a`$ to obtain the exponent $`h`$.
### II.2 Fluctuations and deviations
In seeking to understand Hackโs law, we are naturally led to wonder about the underlying distribution that gives rise to this mean relationship. By considering fluctuations, we begin to see Hackโs law as an expression of basin morphology. What shapes of basins characterized by $`(a,l)`$ are possible and with what probability do they occur?
An important point here is that Hackโs law does not exactly specify basin shapes. An additional connection to Euclidean dimensions of the basin is required. We may think of a basinโs longitudinal length $`L_{}`$ and its width $`L_{}`$. The main stream length $`l`$ is reported to scale with $`L_{}`$ as
$$lL_{}^d,$$
(3)
where typically $`1.0d1.1`$, maritan96a ; tarboton90 . Hence, we have $`al^{1/h}L_{}^{d/h}`$. All other relevant scaling laws exponents can be related to the pair of exponents $`(d,h)`$ which therefore characterize the universality class of a river network dodds2000pa ; dodds99pa . If $`d/h=2`$ we have that basins are self-similar whereas if $`d/h<2`$, we have that basins are elongating. So, while Hackโs law gives a sense of basin allometry, the fractal properties of main stream lengths need also be known in order to properly quantify the scaling of basin shape.
In addition to fluctuations, complementary insights are provided by the observation and understanding of deviations from scaling. We are thus interested in discerning the regularities and quirks of the joint probability distribution $`P(a,l)`$. We will refer to $`P(a,l)`$ as the Hack distribution.
Hack distributions for the Kansas river basin and the Mississippi river basin are given in Figures 1(a) and 1(b). Fluctuations about and deviations from scaling are immediately evident for the Kansas and to a lesser extent for the Mississippi. The first section of the paper will propose and derive analytic forms for the Hack distribution under the assumption of uniform scaling with no deviations. Here, as well as in the following two papers of this series dodds2000ub ; dodds2000uc , we will motivate our results with a random network model originally due to Scheidegger scheidegger67 .
We then expand our discussion to consider deviations from exact scaling. In the case of the Kansas river, a striking example of deviations from scaling is the linear branch separated from the body of the main distribution shown in Figure 1(a). This feature is less prominent in the lower resolution Mississippi data. Note that this linear branch is not an artifact of the measurement technique or data set used. This will be explained in our discussion of deviations at small scales in the paperโs second section.
We then consider the more subtle deviations associated with intermediate scales. At first inspection, the scaling appears to be robust. However, we find gradual drifts in โexponentsโ that prevent us from identifying a precise value of $`h`$ and hence a corresponding universality class.
Both distributions also show breakdowns in scaling for large areas and stream lengths and this is addressed in the final part of our section on deviations. The reason for such deviations is partly due to the decrease in number of samples and hence self-averaging, as area and stream lengths are increased. However, we will show that the direction of the deviations depends on the overall basin shape. We will quantify the extent to which such deviations can occur and the effect that they have on measurements of Hackโs exponent $`h`$.
Throughout the paper, we will return to the Hack distributions for the Kansas and the Mississippi rivers as well as data obtained for the Amazon, the Nile and the Congo rivers.
## III Fluctuations: an analytic form for the Hack distribution
To provide some insight into the nature of the underlying Hack distribution, we present a line of reasoning that will build up from Hackโs law to a scaling form of $`P(a,l)`$. First let us assume for the present discussion of fluctuations that an exact form of Hackโs law holds:
$$l=\theta a^h$$
(4)
where we have introduced the coefficient $`\theta `$ which we discuss fully later on. Now, since Hackโs law is a power law, it is reasonable to postulate a generalization of the form
$$P(l|a)=\frac{1}{a^h}F_l\left(\frac{l}{a^h}\right).$$
(5)
The prefactor $`1/a^h`$ provides the correct normalization and $`F_l`$ is the โscaling functionโ we hope to understand. The above will be our notation for all conditional probabilities. Implicit in equation (5) is the assumption that all moments and the distribution itself also scale. For example, the $`q`$th moment of $`P(l|a)`$ is
$$l^q(a)a^{qh},$$
(6)
which implies
$$l^q(a)=k^{qh}l^q(ak).$$
(7)
where $`k`$. Also, for the distribution $`P(l|a)`$ it follows from equation (5) that
$$k^hP(lk^h|ak)=\frac{k^h}{a^hk^h}F_l\left(\frac{lk^h}{a^hk^h}\right)=P(l|a).$$
(8)
We note that previous investigations of Hackโs law maritan96a ; rigon96 consider the generalization in equation (5). Rigon et al. rigon98 also examine the behavior of the moments of the distribution $`P(l|a)`$ for real networks. Here, we will go further to characterize the full distribution $`P(a,l)`$ as well as both $`P(l|a)`$ and $`P(a|l)`$. Along these lines, Rigon et al. rigon96 suggest that the function $`F_l(x)`$ is a โfinite-sizeโ scaling function analogous to those found in statistical mechanics, i.e.: $`F_l(x)0`$ as $`x\mathrm{}`$ and $`F_l(x)c`$ as $`x0`$. However, as we will detail below, the restrictions on $`F_l(x)`$ can be made stronger and we will postulate a simple Gaussian form. More generally, $`F_l(x)`$ should be a unimodal distribution that is non-zero for an interval $`[x_1,x_2]`$ where $`x_1>0`$. This is so because for any given fixed basin area $`a`$, there is a minimum and maximum $`l`$ beyond which no basin exists. This is also clear upon inspection of Figures 1(a) and 1(b).
We observe that neither drainage area nor main stream length possess any obvious features so as to be deemed the independent variable. Hence, we can also view Hackโs law as its inversion $`al^{1/h}`$. Note that the constant of proportionality is not necessarily $`\theta ^{1/h}`$ and is dependent on the nature of the full Hack distribution. We thus have another scaling ansatz as per equation (5)
$$P(a|l)=1/l^{1/h}F_a(a/l^{1/h}).$$
(9)
The conditional probabilities $`P(l|a)`$ and $`P(a|l)`$ are related to the joint probability distribution as
$$P(a,l)=P(a)P(l|a)=P(l)P(a|l),$$
(10)
where $`P(l)`$ and $`P(a)`$ are the probability densities of main stream length and area. These distributions are in turn observed to be power laws both in real world networks and models rodriguez-iturbe97 ; takayasu88 ; meakin91 :
$$P(a)N_aa^\tau \text{and}P(l)N_ll^\gamma .$$
(11)
where $`N_a`$ and $`N_l`$ are appropriate prefactors and the tilde indicates asymptotic agreement between both sides for large values of the argument. Furthermore, the exponents $`\tau `$ and $`\gamma `$ are related to Hackโs exponent $`h`$ via the scaling relations maritan96a ; dodds99pa
$$\tau =2h\text{and}\gamma =1/h.$$
(12)
Equations (5), (9), (10), (11), and (12) combine to give us two forms for $`P(a,l)`$,
$$P(a,l)=1/a^2F(l/a^h)=1/l^{2/h}G(a/l^{1/h}),$$
(13)
where $`x^2F(x^h)=G(x)`$ and, equivalently, $`F(y)=y^2G(y^{1/h})`$.
## IV Random directed networks
We will use results from the Scheidegger model scheidegger67 to motivate the forms of these distributions. In doing so, we will also connect with some problems in the theory of random walks.
Scheideggerโs model of river networks is defined on a triangular lattice as indicated by Figure 2. Flow in the figure is directed down the page. At each site, the stream flow direction is randomly chosen between the two diagonal directions shown. Periodic boundary conditions are applied in all of our simulations. Each site locally drains an area of $`\alpha ^2`$, where the lattice unit $`\alpha `$ is the distance between neighboring sites, and each segment of stream has a length $`\alpha `$. For simplicity, we will take $`\alpha `$ to be unity. We note that connections exist between the Scheidegger model and models of particle aggregation takayasu88 ; huber91 , Abelian sandpiles dhar90 ; dhar92 ; dhar99 and limiting cases of force chain models in granular media coppersmith96 .
Since Scheideggerโs model is based on random flow directions, the Hack distributions have simple interpretations. The boundaries of drainage basins in the model are random walks. Understanding Hackโs law therefore amounts to understanding the first collision time of two random walks that share the same origin in one dimension. If we subtract the graph of one walk from the other, we see that the latter problem is itself equivalent to the first return problem of random walks feller68I .
Many facets of the first return problem are well understood. In particular, the probability of $`n`$, the number of steps taken by a random walk until it first returns to the origin, is asymptotically given by
$$P(n)\frac{1}{\sqrt{2}\pi }n^{3/2}.$$
(14)
But this number of steps is also the length of the basin $`l`$. Therefore, we have
$$P(l)\frac{2}{\pi }l^{3/2},$$
(15)
where because we are considering the difference of two walks, we use $`P(l)=P(n/2)|_{n=l}`$. Also, we have found the prefactor $`N_l=2/\pi `$.
We thus have that $`\gamma =3/2`$ for the Scheidegger model. The scaling relations of equation (12) then give $`h=2/3`$ and $`\tau =4/3`$. The value of $`h`$ is also readily obtained by noting that the typical area of a basin of length $`l`$ is $`all^{1/2}=l^{3/2}=l^{1/h}`$ since the boundaries are random walks.
## V Area-length distribution for random, directed networks
Something that is less well studied is the joint distribution of the area enclosed by a random walk and the number of steps to its first return. In terms of the Scheidegger model, this is precisely the Hack distribution.
We motivate some general results based on observations of the Scheidegger model. Figure 3 shows the normalized distributions $`P(al^{3/2})`$ and $`P(la^{2/3})`$ as derived from simulations of the model. Given the scaling ansatzes for $`P(l|a)`$ and $`P(a|l)`$ in equations (5) and (9), we see that $`P(y=la^{2/3})=F_l(y)`$ and $`P(x=al^{3/2})=F_a(x)`$.
Note that we have already used Hackโs law for the Scheidegger model with $`h=2/3`$ to obtain these distributions. The results are for ten realizations of the model on a $`10^4`$ by $`10^4`$ lattice, taking $`10^7`$ samples from each of the ten instances. For $`P(a|l)`$, only sites where $`l100`$ were taken, and similarly, for $`P(l|a)`$, only sites where $`a500`$ were included in the histogram.
We postulate that the distribution $`P(y=la^{2/3})`$ is a Gaussian having the form
$$P(l|a)=\frac{1}{\sqrt{2\pi }a^{2/3}\eta }\mathrm{exp}\{(la^{2/3}\theta )^2/2\eta ^2\}$$
(16)
We estimate the mean of $`F_l`$ to be $`\theta 1.675`$ (this is the same $`\theta `$ as found in equation (4)) and the standard deviation to be $`\eta 0.321`$. The fit is shown in Figure 3 as a solid line. The above equation agrees with the form of the scaling ansatz of equation (5) and we now have the assertion that $`F_l`$ is a Gaussian defined by the two parameters $`\theta `$ and $`\eta `$.
Note that the $`\theta `$ and $`\eta `$ are coefficients for the actual mean and standard deviation. In other words, for fixed $`a`$, the mean of $`P(l|a)`$ is $`\theta a^{2/3}`$ and its standard deviation is $`\eta a^{2/3}`$. Having observed their context, we will refer to $`\theta `$ and $`\eta `$ as the Hack mean coefficient and Hack standard deviation coefficient.
From this starting point we can create $`P(a,l)`$ and $`P(a|l)`$, the latter providing a useful test. Since $`P(a)N_aa^\tau =N_aa^{4/3}`$, as per equation (11), we have
$`P(a,l)`$ $`=`$ $`{\displaystyle \frac{N_a}{a^{4/3}}}{\displaystyle \frac{1}{\sqrt{2\pi }a^{2/3}\eta }}\mathrm{exp}\{(la^{2/3}\theta )^2/2\eta ^2\},`$ (17)
$`=`$ $`{\displaystyle \frac{N_a}{\sqrt{2\pi }a^2\eta }}\mathrm{exp}\{(la^{2/3}\theta )^2/2\eta ^2\}.`$
As expected, we observe the form of equation (17) to be in accordance with that of equation (13). Note that the scaling function $`F`$ (and equivalently $`G`$) is defined by the three parameters $`\theta `$, $`\eta `$ and $`N_a`$, the latter of which may be determined in terms of the former as we will show below. Also, since we expect all scaling functions to be only asymptotically correct, we cannot use equation (17) to find an expression for the normalization $`N_a`$. Equation (17) ceases to be valid for small $`a`$ and $`l`$. However, we will be able to do so once we have $`P(a|l)`$ since we are able to presume $`l`$ is large and therefore that the scaling form is exact. Using equation (15) and the fact that $`P(a|l)=P(a,l)/P(l)`$ from equation (10) we then have
$`P(a|l)={\displaystyle \frac{\pi l^{3/2}}{2}}{\displaystyle \frac{N_a}{\sqrt{2\pi }a^2\eta }}\mathrm{exp}\{(l/a^{2/3}\theta )^2/2\eta ^2\},`$ (18)
$`=`$ $`{\displaystyle \frac{N_a\sqrt{\pi }l^{3/2}}{2^{3/2}a^2\eta }}\mathrm{exp}\{(l/a^{2/3}\theta )^2/2\eta ^2\},`$
$`=`$ $`{\displaystyle \frac{1}{l^{3/2}}}{\displaystyle \frac{N_a\sqrt{\pi }}{2^{3/2}\eta }}(a/l^{3/2})^2`$
$`\times \mathrm{exp}\{((a/l^{3/2})^{2/3}\theta )^2/2\eta ^2\}.`$
In rearranging the expression of $`P(a|l)`$, we have made clear that its form matches that of equation (5).
A closed form expression for the normalization factor $`N_a`$ may now be determined by employing the fact that $`_{a=0}^{\mathrm{}}\text{d}aP(a|l)=1`$.
$`1={\displaystyle _{a=0}^{\mathrm{}}}\text{d}aP(a|l),`$ (19)
$`=`$ $`{\displaystyle _{a=0}^{\mathrm{}}}{\displaystyle \frac{\text{d}a}{l^{3/2}}}{\displaystyle \frac{N_a\sqrt{\pi }}{2^{3/2}\eta }}(a/l^{3/2})^2`$
$`\times \mathrm{exp}\{((a/l^{3/2})^{2/3}\theta )^2/2\eta ^2\},`$
$`=`$ $`{\displaystyle \frac{N_a\sqrt{3\pi }}{2^{5/2}\eta }}{\displaystyle _{u=0}^{\mathrm{}}}\text{d}uu^{1/2}\mathrm{exp}\{(u\theta )^2/2\eta ^2\},`$
where we have used the substitution $`a/l^{3/2}=u^{3/2}`$ and hence also $`l^{3/2}\text{d}a=(3/2)u^{5/2}\text{d}u`$. We therefore have
$$N_a=\frac{2^{5/2}\eta }{\sqrt{3\pi }}\left[_{u=0}^{\mathrm{}}\text{d}uu^{1/2}\mathrm{exp}\{(u\theta )^2/2\eta ^2\}\right]^1.$$
(20)
We may thus write down all of the scaling functions $`F_l`$, $`F_a`$, $`F`$ and $`G`$ for the Scheidegger model:
$`F_l(z)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi }\eta }}\mathrm{exp}\{(z\theta )^2/2\eta ^2\},`$ (21)
$`F_a(z)`$ $`=`$ $`{\displaystyle \frac{N_a\sqrt{\pi }}{2^{3/2}\eta }}z^2\mathrm{exp}\{(z^{2/3}\theta )^2/2\eta ^2\},`$ (22)
$`F(z)`$ $`=`$ $`{\displaystyle \frac{N_a}{\sqrt{2\pi }\eta }}\mathrm{exp}\{(z\theta )^2/2\eta ^2\},\text{and}`$ (23)
$`G(z)`$ $`=`$ $`{\displaystyle \frac{2N_a}{\sqrt{\pi }2^{3/2}\eta }}z^2\mathrm{exp}\{(z^{2/3}\theta )^2/2\eta ^2\}.`$ (24)
Recall that all of these forms rest on the assumption that $`F_l(z)`$ is a Gaussian. In order to check this assumption, we return to Figure 3. The empirical distribution $`P(z=a/l^{3/2})`$ is shown on the left marked with circles. The solid line through these points is $`F_a(z)`$ as given above in equation (21). There is an excellent match so we may be confident about our proposed form for $`F_a(z)`$. We note that the function $`F_a(z)`$ may be thought of as a fractional inverse Gaussian distribution, the inverse Gaussian being a well known distribution arising in the study of first passage times for random walks feller68I . It is worth contemplating the peculiar form of $`P(a|l)`$ in terms of first return random walks. Here, we have been able to postulate the functional form of the distribution of areas bound by random walks that first return after $`n`$ steps. If one could understand the origin of the Gaussian and find analytic expressions for $`\theta `$ and $`\eta `$, then the problem would be fully solved.
## VI Area-length distribution extended to real networks
We now seek to extend these results for Scheideggerโs model to real world networks. We will look for the same functional forms for the Hack distributions that we have found above. The conditional probability distributions pertaining to Hackโs law take on the forms
$$P(l|a)=1/a^hF_l(la^h)=\frac{a^h}{\sqrt{2\pi }\eta }\mathrm{exp}\{(la^h\theta )^2/2\eta ^2\},$$
(25)
and
$`P(a|l)=l^{1/h}F_a(al^{1/h})`$
$`=`$ $`l^{1/h}{\displaystyle \frac{N_a}{\sqrt{2\pi }N_l\eta }}(al^{1/h})^2\mathrm{exp}\{((al^{1/h})^h\theta )^2/2\eta ^2\},`$
(26)
and the full Hack distribution is given by
$`P(a,l)`$ $`=`$ $`a^2F(la^h)=l^{2/h}G(al^{1/h})`$ (27)
$`=`$ $`{\displaystyle \frac{N_a}{\sqrt{2\pi }\eta a^2}}\mathrm{exp}\{(la^h\theta )^2/2\eta ^2\}.`$
with $`N_a`$ determined by equation (20). The three parameters $`h`$, the Hack exponent, $`\theta `$, the Hack mean coefficient, and $`\eta `$, the Hack standard deviation coefficient, are in principle landscape dependent. Furthermore, in $`\eta `$ we have a basic measure of fluctuations in the morphology of basins.
Figures 4 and 5 present Hack scaling functions for the Mississippi and Nile river basins. These Figures is to be compared with the results for the Scheidegger model in Figure 3.
For both rivers, Hackโs exponent $`h`$ was determined first from a stream ordering analysis (we discuss stream ordering later in Section X). Estimates of the parameters $`\theta `$ and $`\eta `$ were then made using the scaling function $`F_l`$ presuming a Gaussian form.
We observe the Gaussian fit for the Mississippi is more satisfactory than that for the Nile. These fits are not rigorously made because even though we have chosen data ranges where deviations (which we address in the following section) are minimal, deviations from scaling do still skew the distributions. The specific ranges used to obtain $`F_l`$ and $`F_a`$ respectively are for the Mississippi: $`8.5<\mathrm{log}_{10}a<9.5`$ and $`4.75<\mathrm{log}_{10}l<6`$, and for the Nile: $`9<\mathrm{log}_{10}a<11`$ and $`5.5<\mathrm{log}_{10}l<6`$ (areas are in km<sup>2</sup> and lengths km). Furthermore, we observe that the estimate of $`h`$ has an effect in the resulting forms of $`F_l`$ and $`F_a`$. Nevertheless, here we are attempting to capture the essence of the generalized form of Hackโs law in real networks.
We then use the parameters $`h`$, $`\theta `$ and $`\eta `$ and equation (26) to construct our theoretical $`F_a`$, the smooth curves in Figures 4 and 5. As for the Scheidegger model data in Figure 3, we see in both examples approximate agreement between the measured $`F_a`$ and the one predicted from the form of $`F_l`$. Table 1 shows estimates of $`h`$, $`\theta `$ and $`\eta `$ for the five major river basins studied.
Given our reservations about the precision of these values of $`\theta `$ and $`\eta `$, we are nevertheless able to make qualitative distinctions. Recalling that $`l=\theta a^h`$, we see that, for fixed $`h`$, higher values of $`\theta `$ indicate relatively longer stream lengths for a given area and hence longer and thinner basins. The results therefore suggest the Nile, and to a lesser degree the Amazon, have basins with thinner profiles than the Congo and, in particular, the Mississippi and Kansas. This seems not unreasonable since the Nile is a strongly directed network constrained within a relatively narrow overall shape. This is somewhat in spite of the fact that the shape of an overall river basin is not necessarily related to its internal basin morphology, an observation we will address later in Section X.
The importance of $`\theta `$ is tempered by the value of $`h`$. Hackโs exponent affects not only the absolute measure of stream length for a given area but also how basin shapes change with increasing area. So, in the case of the Kansas the higher value of $`h`$ suggests basin profiles thin with increasing size. This is in keeping with overall directedness of the network. Note that our measurements of the fractal dimension $`d`$ of stream lengths for the Kansas place it to be $`d=1.04\pm 0.02`$. Therefore, $`d/h1.9<2`$ and elongation is still expected when we factor in the scaling of $`l`$ with $`L_{}`$.
The Nile and Amazon also all have relatively high $`\eta `$ indicating greater fluctuations in basin shape. In comparison, the Mississippi, Kansas and Congo appear to have less variation. Note that the variability of the Kansas is in reasonable agreement with that of the whole Mississippi river network for which it is a sub-basin.
Finally, regardless of the actual form of the distribution underlying Hackโs law, fluctuations are always present and an estimate of their extent is an important measurement. Thus, the Hack mean and standard deviation coefficients, $`\theta `$ and $`\eta `$, are suggested to be of sufficient worth so as to be included with any measurement of the Hack exponent $`h`$.
## VII Deviations from scaling
In generalizing Hackโs law, we have sought out regions of robust scaling, discarding ranges where deviations become prominent. We now bring our attention to the nature of the deviations themselves.
We observe three major classes of deviations which we will define by the scales at which they occur: small, intermediate and large. Throughout the following sections we primarily consider deviations from the mean version of Hackโs law, $`l=\theta a^h`$, given in equation (4). Much of the understanding we gain from this will be extendable to deviations for higher moments.
To provide an overview of what follows, examples of mean-Hack distributions for the Kansas river and the Mississippi are shown in Figures 6(a) and 6(b). Hackโs law for the Kansas river exhibits a marked deviation for small areas, starting with a near linear relationship between stream length and area. A long crossover region of several orders of magnitude in area then leads to an intermediate scaling regime wherein we attempt to determine the Hack exponent $`h`$. In doing so, we show that such regions of robust scaling are surprisingly limited for river network quantities. Moreover, we observe that, where present, scaling is only approximate and that no exact exponents can be ascribed to the networks we study here. It follows that the identification of universality classes based on empirical evidence is a hazardous step.
Finally, the approximate scaling of this intermediate region then gives way to a break down in scaling at larger scales due to low sampling and correlations with basin shape. The same deviations are present in the relatively coarse-grained Mississippi data but are less pronounced.
## VIII Deviations at small scales
At small scales, we find the mean-Hack distribution to follow a linear relationship, i.e., $`la^1`$. This feature is most evident for the Kansas river as shown in Figure 6(a) and in more detail in Figure 7. The linear regime persists for nearly 1.5 orders of magnitude in basin area. To a lesser extent, the same trend is apparent in the Mississippi data, Figure 6(b).
Returning to the full Hack distribution of Figures 1(a) and 1(b), we begin to see the origin of this linear regime. In both instances, a linear branch separates from the body of the main distribution. Since $`l`$ cannot grow faster than $`a^1`$, the linear branch marks an upper bound on the extent of the distribution in $`(a,l)`$ coordinates. When averaged to give the mean-Hack distribution, this linear data dominates the result for small scales.
We find this branch evident in all Hack distributions. It is not an artifact of resolution and in fact becomes more pronounced with increased map precision. The origin of this linear branch is simple: data points along the branch correspond to positions in narrow sub-networks, i.e., long, thin โvalleys.โ To understand the separation of this linear branch from the main body of the distribution, consider Figure 8 which depicts a stream draining such a valley with length and area $`(l_1,a_1)`$ that meets a stream from a basin with characteristics $`(l_2,a_2)`$. The area and length of the basin formed at this junction is thus $`(a_3,l_3)=(a_1+a_2,\mathrm{max}(l_1,l_2))`$ (in the Figure, $`\mathrm{max}(l_1,l_2)=l_2`$). The greater jump in area moves the point across into the main body of the Hack distribution, creating the separation of the linear branch.
In fact, the full Hack distribution is itself comprised of many such linear segments. As in the above example, until a stream does not meet any streams of comparable size, then its area and length will roughly increase in linear fashion. When it does meet such a stream, there is a jump in area and the trace of a new linear segment is started in the distribution. We will see this most clearly later on when we study deviations at large scales.
For very fine scale maps, on the order of meters, we might expect to pick up the scale of the unchannelized, convex regions of a landscape, i.e., โhillslopesโ dietrich98 . This length scale represents the typical separation of branches at a networkโs finest scale. The computation of stream networks for these hillslope regions would result in largely non-convergent (divergent or parallel) flow. Therefore, we would have linear โbasinsโ that would in theory contribute to the linear branch we observe dodds2000pa . Potentially, the crossover in Hackโs law could be used as a determinant of hillslope scale, a crucial parameter in geomorphology dietrich98 ; dietrich93a . However, when long, thin network structures are present in a network, this hillslope scale is masked by their contribution.
Whether because of the hillslope scale or linear network structure, we see that at small scales, Hackโs law will show a crossover in scaling from $`h=1`$ to a lower exponent. The crossoverโs position depends on the extent of linear basins in the network. For example, in the Kansas River basin, the crossover occurs when $`(l,a)=(4\times 10^3\text{m},10^7\text{m}^2)`$. Since increased map resolution can only increase measures of length, the crossoverโs position must occur at least at such a length scale which may be many orders of magnitude greater than the scale of the map.
However, the measurement of the area of such linear basins will potentially grow with coarse-graining. Note that for the Mississippi mean-Hack distribution, the crossover begins around $`a=10^{6.5}`$ m<sup>2</sup> whereas for the Kansas, the crossover initiates near $`a=10^{5.5}`$ m<sup>2</sup> but the ends of the crossovers in both cases appear to agree, occurring at around $`a=10^7`$ m<sup>2</sup>. Continued coarse-graining will of course eventually destroy all statistics and introduce spurious deviations. Nevertheless, we see here that the deviation which would only be suggested in the Mississippi data is well confirmed in the finer-grain Kansas data.
## IX Deviations at intermediate scales
As basin area increases, we move out of the linear regime, observing a crossover to what would be considered the normal scaling region of Hackโs law. We detail our attempts to measure the Hack exponents for the Kansas and Mississippi examples. Rather than relying solely on a single regression on a mean-Hack distribution, we employ a more precise technique that examines the distributionโs derivative. As we will show, we will not be able to find a definite value for the Hack exponent in either case, an important result in our efforts to determine whether or not river networks belong to specific universality classes.
To determine $`h`$, we consider Hackโs law (equation (4)) explicitly in logarithmic coordinates,
$$\mathrm{log}_{10}l(a)=\mathrm{log}_{10}\theta +h\mathrm{log}_{10}a.$$
(28)
The derivative of this equation with respect to $`\mathrm{log}_{10}a`$ then gives Hackโs exponent as a function of area,
$$h(a)=\frac{\text{d}}{\text{d}\mathrm{log}_{10}a}\mathrm{log}_{10}l(a).$$
(29)
We may think of $`h(a)`$ as a โlocal Hack exponent.โ Note that non-constant trends in $`h(a)`$ indicate scaling does not hold. We calculate the discrete derivative as above for the Kansas and Mississippi. We smooth the data by taking running averages with varying window sizes of $`n`$ samples, the results for $`n=32`$ being shown in Figures 9(a) and 9(b) where the spacing of $`\mathrm{log}a`$ is $`0.02`$ orders of magnitude. Thus, the running averages for the figures are taken over corresponding area ranges 0.64 orders of magnitude.
Now, if the scaling law in question is truly a scaling law, the above type of derivative will fluctuate around a constant value of exponent over several orders of magnitude. With increasing $`n`$, these fluctuations will necessarily decrease and we should see the derivative holding steady around the exponentโs value.
At first glance, we notice considerable variation in $`h(a)`$ for both data sets with the Kansas standing out. Fluctuations are reduced with increasing $`n`$ but we observe continuous variation of the local Hack exponent with area. For the example of the Kansas, the linear regime and ensuing crossover appear as a steep rise followed by a drop and then another rise during all of which the local Hack exponent moves well below $`1/2`$.
It is after these small scale fluctuations that we would expect to find Hackโs exponent. For the Kansas river data, we see the derivative gradually climbs for all values of $`n`$ before reaching the end of the intermediate regime where the putative scaling breaks down altogether.
For the Kansas show in Figure 9(a), the dashed line represents $`h=0.57`$, our estimate of Hackโs exponent from simple regression on the mean-Hack distribution of Figure 6(a). For the regression calculation, the intermediate region was identified from the figure to be $`10^7<a<10^{10}`$ m<sup>2</sup>. We see from the smoothed derivative in Figure 9(a) that the value $`h=0.57`$ is not precise. After the crossover from the linear region has been completed, we observe a slow rise from $`h0.54`$ to $`h0.63`$. Thus, the local Hack exponent $`h(a)`$ gradually climbs above $`h=0.57`$ rather than fluctuate around it.
A similar slow change in $`h(a)`$ is observed for the Mississippi data. We see in Figure 9(b) a gradual rise and then fall in $`h(a)`$. The dashed line here represents $`h=0.55`$, the value of which was determined from Figure 6(b) using regression on the range $`10^9<a<10^{11.5}`$ m<sup>2</sup>. The range of $`h(a)`$ is roughly $`[0.52,0.58]`$. Again, while $`h=0.55`$ approximates the derivative throughout this intermediate range of Hackโs law, we cannot claim it to be a precise value.
We observe the same drifts in $`h(a)`$ in other datasets and for varying window size $`n`$ of the running average. The results suggest that we cannot assign specific Hack exponents to these river networks and are therefore unable to even consider what might be an appropriate universality class. The value of $`h`$ obtained by regression analysis is clearly sensitive to the the range of $`a`$ used. Furthermore, these results indicate that we should maintain healthy reservations about the exact values of other reported exponents.
## X Deviations at large scales
We turn now to deviations from Hackโs law at large scales. As we move beyond the intermediate region of approximate scaling, fluctuations in $`h(a)`$ begin to grow rapidly. This is clear on inspection of the derivatives of Hackโs law in Figures 9(a) and 9(b). There are two main factors conspiring to drive these fluctuations up. The first is that the number of samples of sub-basins with area $`a`$ decays algebraically in $`a`$. This is just the observation that $`P(a)a^\tau `$ as per equation (11). The second factor is that fluctuations in $`l`$ and $`a`$ are on the order of the parameters themselves. This follows from our generalization of Hackโs law which shows, for example, that the moments $`l^q`$ of $`P(l|a)`$ grow like $`a^{qh}`$. Thus, the standard deviation grows like the mean: $`\sigma (l)=(l^2l^2)^{1/2}a^hl`$.
### X.1 Stream ordering and Hortonโs laws
So as to understand these large scale deviations from Hackโs law, we need to examine network structure in depth. One way to do this is by using Horton-Strahler stream ordering horton45 ; strahler57 and a generalization of the well-known Hortonโs laws horton45 ; schumm56a ; peckham99 ; dodds99pa ; dodds2000ub . This will naturally allow us to deal with the discrete nature of a network that is most apparent at large scales.
Stream ordering discretizes a network into a set of stream segments (or, equivalently, a set of nested basins) by an iterative pruning. Source streams (i.e., those without tributaries) are designated as stream segments of order $`\omega =1`$. These are removed from the network and the new source streams are then labelled as order $`\omega =2`$ stream segments. The process is repeated until a single stream segment of order $`\omega =\mathrm{\Omega }`$ is left and the basin itself is defined to be of order $`\mathrm{\Omega }`$.
Natural metrics for an ordered river network are $`n_\omega `$, the number of order $`\omega `$ stream segments (or basins), $`\overline{a}_\omega `$, the average area of order $`\omega `$ basins, $`\overline{l}_\omega `$, the average main stream length of order $`\omega `$ basins, and $`\overline{l}_\omega ^{\text{ (s)}}`$, the average length of order $`\omega `$ stream segments. Hortonโs laws state that these quantities change regularly from order to order, i.e.,
$$\frac{n_\omega }{n_{\omega +1}}=R_n\text{and}\frac{\overline{X}_{\omega +1}}{\overline{X}_\omega }=R_X,$$
(30)
where $`X=a`$, $`l`$ or $`l^{\text{(s)}}`$. Note that all ratios are defined to be greater than unity since areas and lengths increase but number decreases. Also, there are only two independent ratios since $`R_aR_n`$ and $`R_lR_{l^{\text{(s)}}}`$ dodds99pa . Hortonโs laws mean that stream-order quantities change exponentially with order. For example, (30) gives that $`l_\omega (R_l)^\omega `$.
### X.2 Discrete version of Hackโs law
Returning to Hackโs law, we examine its large scale fluctuations with the help of stream ordering. We are interested in the size of these fluctuations and also how they might correlate with the overall shape of a basin. First, we note that the structure of the network at large scales is explicitly discrete. Figure 10 demonstrates this by plotting the distribution of $`(a,l)`$ without the usual logarithmic transformation. Hackโs law is seen to be composed of linear fragments. As explained above in Figure 8, areas and length increase in proportion to each other along streams where no major tributaries enter. As soon as a stream does combine with a comparable one, a jump in drainage area occurs. Thus, we see in Figure 10 isolated linear segments which upon ending at a point $`(a_1,l_1)`$ begin again at $`(a_1+a_2,l_1)`$, i.e., the main stream length stays the same but the area is shifted.
We consider a stream ordering version of Hackโs law given by the points $`(\overline{a}_\omega ,\overline{l}_\omega )`$. The scaling of these data points is equivalent to scaling in the usual Hackโs law. Also, given Hortonโs laws, it follows that $`h=\mathrm{ln}R_l/\mathrm{ln}R_n`$ (using $`R_aR_n`$). Along the lines of the derivative we introduced to study intermediate scale fluctuations in equation (29), we have here an order-based difference:
$$h_{\omega ,\omega 1}=\frac{\mathrm{log}\overline{l}_\omega /\overline{l}_{\omega 1}}{\mathrm{log}\overline{a}_\omega /\overline{a}_{\omega 1}}.$$
(31)
We can further extend this definition to differences between non-adjacent orders:
$$h_{\omega ,\omega ^{}}=\frac{\mathrm{log}\overline{l}_\omega /\overline{l}_\omega ^{}}{\mathrm{log}\overline{a}_\omega /\overline{a}_\omega ^{}}.$$
(32)
This type of difference, where $`\omega ^{}<\omega `$, may be best thought of as a measure of trends rather than an approximate discrete derivative.
Using these discrete differences, we examine two features of the order-based versions of Hackโs law. First we consider correlations between large scale deviations within an individual basin and second, correlations between overall deviations and basin shape. For the latter, we will also consider deviations as they move back into the intermediate scale. This will help to explain the gradual deviations from scaling we have observed at intermediate scales.
Since deviations at large scales are reflective of only a few basins, we require an ensemble of basins to provide sufficient statistics. As an example of such an ensemble, we take the set of order $`\mathrm{\Omega }=7`$ basins of the Mississippi basin. For the dataset used here where the overall basin itself is of order $`\mathrm{\Omega }=11`$, we have 104 order $`\mathrm{\Omega }=7`$ sub-basins. The Horton averages for these basins are $`\overline{a}_716600`$ km<sup>2</sup>, $`\overline{l}_7350`$ km, and $`\overline{L}_7210`$ km.
For each basin, we first calculate the Horton averages $`(\overline{a}_\omega ,\overline{l}_\omega )`$. We then compute $`h_{\omega ,\omega 1}`$, the Hack difference given in equation (31). To give a rough picture of what is observed, Figure 11 shows a scatter plot of $`h_{\omega ,\omega 1}`$ for all order $`\mathrm{\Omega }=7`$ basins. Note the increase in fluctuations with increasing $`\omega `$. This increase is qualitatively consistent with the smooth versions found in the single basin examples of Figures 9(a) and 9(b). In part, less self-averaging for larger $`\omega `$ results in a greater spread in this discrete derivative. However, as we will show, these fluctuations are also correlated with fluctuations in basin shape.
### X.3 Effect of basin shape on Hackโs law
In what follows, we extract two statistical measures of correlations between deviations in Hackโs law and overall basin shape. These are $`r`$, the standard linear correlation coefficient and $`r_s`$, the Spearman rank-order correlation coefficient press92 ; lehmann75 ; sprent93 . For $`N`$ observations of data pairs $`(u_i,v_i)`$, $`r`$ is defined to be
$$r=\frac{_{i=1}^N(u_i\mu _u)(v_i\mu _v)}{_{i=1}^N(x_i\mu _u)^2_{i=1}^N(y_i\mu _v)^2}=\frac{C(u,v)}{\sigma _u\sigma _v},$$
(33)
where $`C(u,v)`$ is the covariance of the $`u_i`$โs and $`v_i`$โs, $`\mu _u`$ and $`\mu _v`$ their means, and $`\sigma _u`$ and $`\sigma _v`$ their standard deviations. The value of Spearmanโs $`r_s`$ is determined in the same way but for the $`u_i`$ and $`v_i`$ replaced by their ranks. From $`r_s`$, we determine a two-sided significance $`p_s`$ via Studentโs t-distribution press92 .
We define $`\kappa `$, a measure of basin aspect ratio, as
$$\kappa =L^2/a.$$
(34)
Long and narrow basins correspond to $`\kappa 1`$ while for short and wide basins, we have $`\kappa 1`$
We now examine the discrete derivatives of Hackโs law in more detail. In order to discern correlations between large scale fluctuations within individual basins, we specifically look at the last two differences in a basin: $`h_{\mathrm{\Omega },\mathrm{\Omega }1}`$ and $`h_{\mathrm{\Omega }1,\mathrm{\Omega }2}`$. For each of the Mississippiโs 104 order $`\mathrm{\Omega }=7`$ basins, these values are plotted against each other in Figure 12. Both our correlation measurements strongly suggest these differences are uncorrelated. The linear correlation coefficient is $`r=0.060`$ and, similarly, we have $`r_s=0.080`$. The significance $`p_s=0.43`$ implies that the null hypothesis of uncorrelated data cannot be rejected.
Thus, for Hackโs law in an individual basin, large scale fluctuations are seen to be uncorrelated. However, correlations between these fluctuations and other factors may still exist. This leads us to our second test which concerns the relationship between trends in Hackโs law and overall basin shape.
Figure 13 shows a comparison of the aspect ratio $`\kappa `$ and $`h_{7,5}`$ for the order $`\mathrm{\Omega }=7`$ basins of the Mississippi. The measured correlation coefficients are $`r=0.50`$ and $`r_s=0.53`$, giving a significance of $`p_s<10^8`$. Furthermore, we find the differences $`h_{7,6}`$ ($`r=0.34`$, $`r_s=0.39`$ and $`p_s<10^4`$) and $`h_{6,5}`$ ($`r=0.35`$, $`r_s=0.34`$ and $`p_s<10^3`$) are individually correlated with basin shape. We observe this correlation between basin shape and trends in Hackโs law at large scales, namely $`h_{\mathrm{\Omega },\mathrm{\Omega }1}`$, $`h_{\mathrm{\Omega }1,\mathrm{\Omega }2}`$ and $`h_{\mathrm{\Omega },\mathrm{\Omega }2}`$, repeatedly in our other data sets. In some cases, correlations extend further to $`h_{\mathrm{\Omega }2,\mathrm{\Omega }3}`$.
Since the area ratio $`R_a`$ is typically in the range 4โ5, Hackโs law is affected by boundary conditions set by the geometry of the overall basin down to sub-basins one to two orders of magnitude smaller in area than the overall basin. These deviations are present regardless of the absolute size of the overall basin. Furthermore, the origin of the basin boundaries being geologic or chance or both is irrelevantโlarge scale deviations will still occur. However, it is reasonable to suggest that particularly strong deviations are more likely the result of geologic structure rather than simple fluctuations.
## XI Conclusion
Hackโs law is a central relation in the study of river networks and branching networks in general. We have shown Hackโs law to have a more complicated structure than is typically given attention. The starting generalization is to consider fluctuations around scaling. Using the directed, random network model, a form for the Hack distribution underlying Hackโs law may be postulated and reasonable agreement with real networks is observed. Questions of the validity of the distribution aside, the Hack mean coefficient $`\theta `$ and the Hack standard deviation coefficient $`\eta `$ should be standard measurements because they provide further points of comparison between theory and other basins.
With the idealized Hack distribution proposed, we may begin to understand deviations from its form. As with any scaling law pertaining to a physical system, cutoffs in scaling must exist and need to be understood. For small scales, we have identified the presence of linear sub-basins as the source of an initial linear relation between area and stream length. At large scales, statistical fluctuations and geologic boundaries give rise to basins whose overall shape produces deviations in Hackโs laws. Both deviations extend over a considerable range of areas as do the crossovers which link them to the region of intermediate scales, particularly the crossover from small scales.
Finally, by focusing in detail on a few large-scale examples networks, we have found evidence that river networks do not belong to well defined universality classes. The relationship between basin area and stream length may be approximately, and in some cases very well, described by scaling laws but not exactly so. The gradual drift in exponents we observe suggests a more complicated picture, one where subtle correlations between basin shape and geologic features are intrinsic to river network structure.
## Acknowledgements
This work was supported in part by NSF grant EAR-9706220 and the Department of Energy grant DE FG02-99ER 15004.
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# hep-th/0005236 Anisotropic Four-Dimensional NS-NS String Cosmology
## Abstract
An anisotropic (Bianchi type I) cosmology is considered in the four-dimensional NS-NS sector of low-energy effective string theory coupled to a dilaton and an axion-like $`H`$-field within a de Sitter-Einstein frame background. The time evolution of this Universe is discussed in both the Einstein and string frames.
PACS number(s): 04.20.Jb, 04.65.+e, 98.80.-k
Pre-Big Bang cosmological models , based on the low energy limit of the string theory, have been intensively investigated in the recent physics literature -. Generically, in these type of models the dynamics of the Universe is dominated by massless bosonic fields. In the string frame, the four-dimensional NS-NS effective action, which is common to both the heterotic and type II string theories, is given by -
$$\widehat{๐ฎ}=d^4x\sqrt{\widehat{g}}e^{2\varphi }\left\{\widehat{R}+\widehat{\kappa }(\varphi )^2\frac{1}{12}\widehat{H}_{[3]}^2\widehat{U}\right\},$$
(1)
where $`H_{\mu \nu \lambda }=_{[\mu }B_{\nu \lambda ]}`$ is an antisymmetric tensor field, $`\widehat{\kappa }`$ is a generalized dilaton coupling constant and $`\widehat{U}=\widehat{U}(\varphi )`$ is a dilaton potential. In addition, $`\widehat{H}_{[3]}^2`$ means the square of the $`H`$-field with respect to the metric $`\widehat{g}_{\mu \nu }`$. The low-energy string action posseses a symmetry property, called scale factor duality, which lets us expect that the present phase of the Universe is preceded by an inflationary pre-Big Bang phase. Explicit dual solutions can be constructed for each Bianchi space-time, except the Bianchi class A types VIII and IX models .
By means of the conformal rescaling
$$g_{\mu \nu }=e^{2\varphi }\widehat{g}_{\mu \nu },$$
(2)
the action (1) can be transformed to the so-called Einstein frame as
$$๐ฎ=d^4x\sqrt{g}\left\{R\kappa (\varphi )^2\frac{1}{12}e^{4\varphi }H_{[3]}^2U\right\},$$
(3)
where $`\kappa =6\widehat{\kappa }`$, $`U=e^{2\varphi }\widehat{U}`$ and $`H_{[3]}^2`$ denotes the square of the antisymmetric field by $`g_{\mu \nu }`$.
Analytic biaxial (two scale factors only) Bianchi type I geometry has been previously considered in for the case with nonvanishing $`H_{[3]}`$ but without a dilaton field potential, i.e. $`U0`$. Triaxial models with the central deficit charge constrained to zero in the presence of a modulus field (representing the evolution of compact extra dimensions) have been analyzed in . Recently, a study of spatially flat and homogeneous string cosmologies, considering the combined effects of the dilaton, modulus, two-form potential and central charge deficit, and using methods from the qualitative theory of differential equations (phase portrait analysis) has been presented in ,
The general Bianchi type I space-time for arbitrary dimensional dilaton gravities, with vanishing antisymmetric tensor $`H_{\mu \nu \lambda }`$ and in the presence of an exponential type dilaton field potential, have been obtained in both the Einstein and string frames .
It is the purpose of the present letter to consider, in the framework of a four-dimensional Bianchi type I geometry, the effects on the dynamics and evolution of the early Universe of a non-vanishing antisymmetric field and of a string frame exponential type dilaton field potential.
In the Einstein frame the field equations, which follow from variation of (3), are given by
$`R_{\mu \nu }\kappa _\mu \varphi _\nu \varphi {\displaystyle \frac{1}{2}}g_{\mu \nu }U`$ (4)
$`{\displaystyle \frac{1}{4}}e^{4\varphi }\left(H_{\mu \alpha \beta }H_\nu {}_{}{}^{\alpha \beta }{\displaystyle \frac{1}{3}}g_{\mu \nu }H^2\right)`$ $`=`$ $`0,`$ (5)
$`_\mu \left(e^{4\varphi }H^{\mu \nu \lambda }\right)`$ $`=`$ $`0,`$ (6)
$`^2\varphi +{\displaystyle \frac{1}{6\kappa }}e^{4\varphi }H^2{\displaystyle \frac{1}{2\kappa }}{\displaystyle \frac{U}{\varphi }}`$ $`=`$ $`0.`$ (7)
Moreover, the $`H`$-field must satisfy the integrability condition (Bianchi identity) $`_{[\mu }H_{\nu \lambda \rho ]}=0`$.
In four dimensions, every three-form field can be dualized to a pseudoscalar. Thus, an appropriate ansatz for the $`H`$-field is
$$H^{\mu \nu \lambda }=\frac{1}{\sqrt{g}}e^{4\varphi }ฯต^{\mu \nu \lambda \rho }_\rho h,$$
(8)
where $`ฯต^{\mu \nu \lambda \rho }=\delta _{[0}^\mu \delta _1^\nu \delta _2^\lambda \delta _{3]}^\rho `$ is the total antisymmetric tensor and $`h=h(t)`$ is the Kalb-Ramond axion field. Then the field equation (6) is satisfied automatically and the Bianchi identity becomes
$$_\mu \left(\sqrt{g}e^{4\varphi }^\mu h\right)=0.$$
(9)
Moreover, we shall assume that in the string frame the dilaton field potential is of exponential type
$$\widehat{U}(\varphi )=\mathrm{\Lambda }e^{2\varphi },$$
(10)
with $`\mathrm{\Lambda }`$ a non-negative constant (de Sitter space-time). Therefore in the Einstein frame the effect of the potential is identical to that of a cosmological constant, $`U(\varphi )=\mathrm{\Lambda }`$.
For the Bianchi type I space-time, in the Einstein frame,
$$ds^2=dt^2+\underset{i=1}{\overset{3}{}}a_i^2(t)\left(dx^i\right)^2,$$
(11)
and the ansatz (8,10), the field equations (5,7,9) take the form
$`3\dot{\theta }+{\displaystyle \underset{i=1}{\overset{3}{}}}\theta _i^2+\kappa \dot{\varphi }^2+{\displaystyle \frac{1}{2}}e^{4\varphi }\dot{h}^2{\displaystyle \frac{1}{2}}\mathrm{\Lambda }`$ $`=`$ $`0,`$ (12)
$`{\displaystyle \frac{1}{V}}{\displaystyle \frac{d}{dt}}(V\theta _i){\displaystyle \frac{1}{2}}\mathrm{\Lambda }`$ $`=`$ $`0,i=1,2,3,`$ (13)
$`\ddot{h}+3\theta \dot{h}+4\dot{\varphi }\dot{h}`$ $`=`$ $`0,`$ (14)
$`{\displaystyle \frac{1}{V}}{\displaystyle \frac{d}{dt}}(V\dot{\varphi }){\displaystyle \frac{1}{\kappa }}e^{4\varphi }\dot{h}^2`$ $`=`$ $`0,`$ (15)
where we have introduced the volume scale factor, $`V:=_{i=1}^3a_i`$, directional Hubble factors, $`\theta _i:=\dot{a}_i/a_i,i=1,2,3`$, and the mean Hubble factor, $`\theta :=_{i=1}^3\theta _i/3=\dot{V}/3V`$. We shall also introduce two basic physical observational quantities in cosmology: the mean anisotropy parameter, $`A:=_{i=1}^3(\theta _i\theta )^2/3\theta ^2`$, and the deceleration parameter, $`q=\frac{d}{dt}\theta ^11`$.
By summing equations (13) we obtain
$$\frac{1}{V}\frac{d}{dt}(V\theta )=\frac{1}{2}\mathrm{\Lambda },$$
(16)
which, together with (13), leads to
$$\theta _i=\theta +K_iV^1,i=1,2,3,$$
(17)
with $`K_i,i=1,2,3`$ being constants of integration satisfying $`_{i=1}^3K_i=0`$.
It is worth noticing that, in this framework, the geometry of the considered Universe, which is described by $`a_i(t),i=1,2,3`$, is determined only by the existence of the cosmological constant $`\mathrm{\Lambda }`$ and is โdecoupledโ from the matter fields $`\varphi `$ and $`H`$. (The effect of matter fields is presented in the magnitude of the parameters, i.e. constants of integration.)
From equation (16) we obtain the time evolution of the mean Hubble factor,
$$\theta (\tau )=\sqrt{\frac{\mathrm{\Lambda }}{6}}\mathrm{coth}\tau ,$$
(18)
leading to
$`V(\tau )`$ $`=`$ $`V_0\mathrm{sinh}\tau ,`$ (19)
$`a_i(\tau )`$ $`=`$ $`a_{i0}\mathrm{sinh}^{\alpha _i^+}{\displaystyle \frac{\tau }{2}}\mathrm{cosh}^{\alpha _i^{}}{\displaystyle \frac{\tau }{2}},i=1,2,3,`$ (20)
where $`\tau :=\sqrt{3\mathrm{\Lambda }/2}(tt_0)`$ and $`\alpha _i^\pm :=1/3\pm \sqrt{2/3\mathrm{\Lambda }}K_i/V_0`$. The mean anisotropy and the deceleration parameter are given by
$`A(\tau )`$ $`=`$ $`{\displaystyle \frac{2K^2}{\mathrm{\Lambda }V_0^2}}\text{sech}^2\tau ,`$ (21)
$`q(\tau )`$ $`=`$ $`3\text{sech}^2\tau 1,`$ (22)
where $`K^2=_{i=1}^3K_i^2`$.
Equation (14) can be integrated to give
$$\dot{h}=Ce^{4\varphi }V^1,$$
(23)
with $`C`$ a constant of integration. Thus the dynamics of the dilaton field in the Einstein frame is described by the following differential equation
$$\mathrm{sinh}\tau \frac{d}{dt}\left(\mathrm{sinh}\tau \dot{\varphi }\right)=\frac{C^2}{\kappa V_0^2}e^{4\varphi },$$
(24)
with the general solution
$$e^{2\varphi (\tau )}=\phi _0^2\left(\mathrm{tanh}^\omega \frac{\tau }{2}+\mathrm{tanh}^\omega \frac{\tau }{2}\right),$$
(25)
where we denote $`\omega :=\sqrt{8\varphi _0/3\mathrm{\Lambda }}/V_0`$, $`\phi _0^2:=\sqrt{C^2/8\kappa \varphi _0}`$ and $`\varphi _0>0`$ is a constant of integration. The antisymmetric tensor field is given by
$$h(\tau )=h_0+\frac{\kappa \sqrt{\varphi _0}}{C}\frac{\mathrm{tanh}^{2\omega }\frac{\tau }{2}1}{\mathrm{tanh}^{2\omega }\frac{\tau }{2}+1},$$
(26)
with $`h_0`$ an arbitrary constant.
The integration constants must satisfy the consistency condition,
$$K^2=\mathrm{\Lambda }V_0^2\kappa \varphi _0,$$
(27)
which follows from equation (12).
In the case of vanishing cosmological constant, $`\mathrm{\Lambda }0`$, the general solution in the Einstein frame of the gravitational field equations for a Bianchi type I geometry with dilaton and Kalb-Ramond axion fields is given by:
$`\theta (t)`$ $`=`$ $`{\displaystyle \frac{1}{3t}},V(t)=V_0t,`$ (28)
$`a_i(t)`$ $`=`$ $`a_{i0}t^{1/3+K_i/V_0},i=1,2,3,`$ (29)
$`A`$ $`=`$ $`3K^2V_0^2=const.,q=2=const.,`$ (30)
$`e^{2\varphi (t)}`$ $`=`$ $`\phi _0^2\left(t^\alpha +t^\alpha \right),`$ (31)
$`h(t)`$ $`=`$ $`h_0{\displaystyle \frac{2\kappa \sqrt{\varphi _0}}{C}}\left(t^{2\alpha }+1\right)^1,`$ (32)
together with the consistency condition
$$K^2=\frac{2}{3}V_0^2\kappa \varphi _0,$$
(33)
where $`\alpha =2\sqrt{\varphi _0}/V_0`$.
In order to find the general solution of the gravitational field equations in the string frame with the line element
$$d\widehat{s}^2=d\widehat{t}^2+\underset{i=1}{\overset{3}{}}\widehat{a}_i^2(\widehat{t})\left(dx^i\right)^2,$$
(34)
we must perform the conformal transformation (2). To obtain a simpler mathematical form of the equations we shall introduce a new variable $`\eta =\mathrm{tanh}\tau /2`$, $`\eta [0,1]`$ and denote $`\widehat{\phi }_0:=\sqrt{8/3\mathrm{\Lambda }}\phi _0`$. Then the string frame time evolution of the Bianchi type I space-time with dilaton and Kalb-Ramond axion fields and an exponential type dilaton potential can be expressed in the following exact parametric form:
$`\widehat{t}(\eta )`$ $`=`$ $`\widehat{t}_0+\widehat{\phi }_0{\displaystyle \frac{\sqrt{\eta ^\omega +\eta ^\omega }}{1\eta ^2}๐\eta },`$ (35)
$`\widehat{V}(\eta )`$ $`=`$ $`\widehat{V}_0(\eta ^\omega +\eta ^\omega )^{3/2}{\displaystyle \frac{\eta }{1\eta ^2}},`$ (36)
$`\widehat{\theta }(\eta )`$ $`=`$ $`{\displaystyle \frac{1}{3\widehat{\phi }_0}}{\displaystyle \frac{1\eta ^2}{\eta \sqrt{\eta ^\omega +\eta ^\omega }}}\left({\displaystyle \frac{3\omega }{2}}{\displaystyle \frac{\eta ^\omega \eta ^\omega }{\eta ^\omega +\eta ^\omega }}+{\displaystyle \frac{1+\eta ^2}{1\eta ^2}}\right),`$ (37)
$`\widehat{a}_i(\eta )`$ $`=`$ $`\widehat{a}_{i0}{\displaystyle \frac{\eta ^{\alpha _i^+}\sqrt{\eta ^\omega +\eta ^\omega }}{(1\eta ^2)^{1/3}}}i=1,2,3,`$ (38)
$`\widehat{A}(\eta )`$ $`=`$ $`{\displaystyle \frac{2K^2}{\mathrm{\Lambda }V_0^2}}\left({\displaystyle \frac{3\omega }{2}}{\displaystyle \frac{\eta ^\omega \eta ^\omega }{\eta ^\omega +\eta ^\omega }}+{\displaystyle \frac{1+\eta ^2}{1\eta ^2}}\right)^2,`$ (39)
$`\widehat{q}(\eta )`$ $`:=`$ $`{\displaystyle \frac{d}{d\widehat{t}}}\widehat{\theta }^11={\displaystyle \frac{1\eta ^2}{\widehat{\phi }_0\sqrt{\eta ^\omega +\eta ^\omega }}}{\displaystyle \frac{d}{d\eta }}\widehat{\theta }^11.`$ (40)
In the case of a vanishing cosmological constant the string frame solution of the gravitational field equations with dilaton and axion fields is given again in a parametric form by:
$`\widehat{t}(t)`$ $`=`$ $`\widehat{t}_0+\phi _0{\displaystyle \sqrt{t^\alpha +t^\alpha }๐t},`$ (41)
$`\widehat{V}(t)`$ $`=`$ $`\widehat{V}_0t(t^\alpha +t^\alpha )^{3/2},`$ (42)
$`\widehat{\theta }(t)`$ $`=`$ $`\phi _0^1\left({\displaystyle \frac{\alpha }{2}}{\displaystyle \frac{t^\alpha t^\alpha }{t^\alpha +t^\alpha }}+{\displaystyle \frac{1}{3}}\right){\displaystyle \frac{1}{t(t^\alpha +t^\alpha )^{1/2}}},`$ (43)
$`\widehat{a}_i(t)`$ $`=`$ $`\widehat{a}_{i0}t^{1/3+K_i/V_0}\sqrt{t^\alpha +t^\alpha },i=1,2,3,`$ (44)
$`\widehat{A}(t)`$ $`=`$ $`{\displaystyle \frac{K^2}{3V_0^2}}\left({\displaystyle \frac{\alpha }{2}}{\displaystyle \frac{t^\alpha t^\alpha }{t^\alpha +t^\alpha }}+{\displaystyle \frac{1}{3}}\right)^2,`$ (45)
$`\widehat{q}(t)`$ $`=`$ $`2{\displaystyle \frac{\alpha \left(\frac{1}{3}\frac{t^\alpha t^\alpha }{t^\alpha +t^\alpha }+\frac{\alpha }{2}\right)}{\left(\frac{\alpha }{2}\frac{t^\alpha t^\alpha }{t^\alpha +t^\alpha }+\frac{1}{3}\right)^2}}.`$ (46)
In the present letter we have presented the exact solution of the gravitational field equations for a Bianchi type I space-time with dilaton and axion fields in both the Einstein and string frames. In the Einstein frame the evolution of the Bianchi type I Universe in the presence of a cosmological constant starts from a singular state, but with finite values of the mean anisotropy and deceleration parameter. In the large time limit the mean anisotropy tends to zero, $`A0`$, and the Universe will end in an isotropic inflationary de Sitter phase with a negative deceleration parameter, $`q<0`$. In the large time limit the dilaton and axion fields become constants. Moreover, in the Einstein frame, the dynamics and evolution of the Universe is determined only by the presence of a cosmological constant and there is no coupling between the metric and the dilaton and axion fields.
In the string frame the dilaton and axion fields are coupled to the metric. Depending on the values of the constant $`\omega `$ there are two distinct types of behavior. In the first type of evolution, corresponding to $`\omega <2/3`$, the Universe starts from a singular state with zero values of the scale factors, $`\widehat{a}_i(0)=0,i=1,2,3`$ and expands indefinitely. In the second case, when $`\omega >2/3`$, the Bianchi type I Universe starts its evolution with infinite values of the scale factors and collapses to a bounce state, corresponding to minimum finite non-zero values of the scale factors. From this non-singular state the Universe starts to expand, ending in an isotropic and inflationary era. The values of the physical quantities at the bounce correspond to the values of $`\eta `$ satisfying the equation $`d\widehat{V}/d\eta =0`$ or
$$\frac{(\eta ^\omega +\eta ^\omega )^{3/2}}{1\eta ^2}\left(\frac{3\omega }{2}\frac{\eta ^\omega \eta ^\omega }{\eta ^\omega +\eta ^\omega }+\frac{1+\eta ^2}{1\eta ^2}\right)=0.$$
(47)
The string frame time variation of the volume scale factor of the Bianchi type I space-time for different values of $`\omega `$ is presented in Fig.1. Independently of which type of evolution classified by the value of $`\omega `$, in the presence of an exponential type dilaton potential and of an axion field, the Bianchi type I Universe always isotropizes in the large time limit, $`\widehat{A}0`$ for $`\widehat{t}\mathrm{}`$. But the dynamics of the mean anisotropy factor is very different for the two types of evolution. During the collapse to the bounce the mean anisotropy increases to an infinite value and then, during the expansionary period, tends rapidly to zero. Hence in this case the expansionary evolution of the Bianchi type I Universe starts with non-singular scale factors and with maximum anisotropy. The string frame time variation of the anisotropy parameter and of the deceleration parameter are represented in the Figures 2 and 3, respectively. In the string frame and in the presence of a dilaton potential the large time evolution is inflationary for all times and for all $`\omega `$.
In the absence of a cosmological constant or a dilaton field potential the Universe does not isotropize. In this case the Einstein frame mean anisotropy is constant for all times and the evolution is of the Kasner type. In the string frame the mean anisotropy tends, in the large time limit, to a constant non-zero value, hence showing that the Universe will never end in an isotropic flat Robertson-Walker type phase. The deceleration parameter in both frames is positive for all times and an inflationary evolution is also impossible. Therefore string cosmological models involving only pure dilaton and axion fields do not have, at least in the case of Bianchi type I anisotropic geometries, the ability of providing realistic cosmological models. To obtain a transition from an anisotropic state to an isotropic inflationary one the โgood oldโ cosmological constant is still the key ingredient.
One of the authors (CMC) would like to thank prof. J.M. Nester for useful comments. The work of CMC was supported in part by the National Science Council (Taiwan) under grant NSC 89-2112-M-008-016.
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# UPRโ885โT Global Fits to Electroweak Data Using GAPP
## 1 PRECISION TESTS
Precision analysis of electroweak interactions follows three major objectives: high precision tests of the SM; the determination of its fundamental parameters; and studies of indications and constraints of possible new physics beyond the SM, such as supersymmetry or new gauge bosons. Currently, the experimental information comes from the very high precision $`Z`$ boson measurements at LEP 1 and the SLC, direct mass measurements and constraints from the Tevatron and LEP 2, and low energy precision experiments, such as in atomic parity violation, $`\nu `$ scattering, and rare decays. These measurements are compared with the predictions of the SM and its extensions. The level of precision is generally very high. Besides the need for high-order loop calculations, it is important to utilize efficient renormalization schemes and scales to ensure sufficient convergence of the perturbative expansions.
The tasks involved called for the creation of a special purpose FORTRAN package, GAPP, short for the Global Analysis of Particle Properties . It is mainly devoted to the calculation of pseudo-observables, i.e., observables appropriately idealized from the experimental reality. The reduction of raw data to pseudo-observables is performed by the experimenters with available packages (e.g., ZFITTER for $`Z`$ pole physics). For cross section and asymmetry measurements at LEP 2 (not implemented in the current version, GAPP$`\mathrm{\_}`$99.7), however, this reduction is not optimal and convoluted expressions should be used instead. GAPP attempts to gather all available theoretical and experimental information; it allows the addition of extra parameters describing new physics; it treats all relevant SM inputs as global fit parameters; and it can easily be updated with new calculations, data, observables, or fit parameters. For clarity and speed it avoids numerical integrations throughout. It is based on the modified minimal subtraction ($`\overline{\mathrm{MS}}`$) scheme which demonstrably avoids large expansion coefficients.
GAPP is endowed with the option to constrain nonstandard contributions to the oblique parameters defined to affect only the gauge boson self-energies (e.g. $`S`$, $`T`$, and $`U`$); specific anomalous $`Z`$ couplings; the number of active neutrinos (with standard couplings to the $`Z`$ boson); and the masses, mixings, and coupling strengths of extra $`Z`$ bosons appearing in models of new physics. With view on the importance of supersymmetric extensions of the SM on one hand, and upcoming experiments on the other, I also included the $`bs\gamma `$ transition amplitude, and intend to add the muon anomalous magnetic moment. In the latter case, there are theoretical uncertainties from hadronic contributions which are partially correlated with the renormalization group (RG) evolutions of the QED coupling and the weak mixing angle. These correlations will be partially taken into account by including heavy quark effects in analytical form; see Ref. for a first step in this direction. By comparing this scheme with more conventional ones, it will also be possible to isolate a QCD sum rule and to rigorously determine the charm and bottom quark $`\overline{\mathrm{MS}}`$ masses, $`\widehat{m}_c`$ and $`\widehat{m}_b`$, with high precision.
## 2 GAPP
### 2.1 Basic structure
In the default running mode of the current version, GAPP$`\mathrm{\_}`$99.7, a fit is performed to 41 observables, out of which 26 are from $`Z`$ pole measurements at LEP 1 and the SLC. The Fermi constant, $`G_F`$ (from the muon lifetime), the electromagnetic fine structure constant, $`\alpha `$ (from the quantum Hall effect), and the light fermion masses are treated as fixed inputs. The exception is $`\widehat{m}_c`$ which strongly affects the RG running<sup>1</sup><sup>1</sup>1Quantities defined in the $`\overline{\mathrm{MS}}`$ scheme are denoted by a caret. of $`\widehat{\alpha }(\mu )`$ for $`\mu >\widehat{m}_c`$. I therefore treat $`\widehat{m}_c`$ as a fit parameter and include an external constraint with an enhanced error to absorb hadronic threshold uncertainties of other quark flavors, as well as theoretical uncertainties from the application of perturbative QCD at relatively low energies. Other fit parameters are the $`Z`$ boson mass, $`M_Z`$, the Higgs boson mass, $`M_H`$, the top quark mass, $`m_t`$, and the strong coupling constant, $`\alpha _s`$, so that there are 37 effective degrees of freedom. Given current precisions, $`M_Z`$ may alternatively be treated as an additional fixed input.
The file fit.f basically contains a simple call to the minimization program MINUIT (from the CERN program library) which is currently used in data driven mode (see smfit.dat). It in turn calls the core subroutine fcn and the $`\chi ^2`$-function chi2, both contained in chi2.f. Subroutine fcn defines constants and flags; initializes parts of the one-loop package FF ; and makes the final call to subroutine values in main.f which drives the output (written to file smfit.out). In chi2 the user actively changes and updates the data for the central values, errors, and correlation coefficients of the observables, and includes or excludes individual contributions to $`\chi ^2`$ (right after the initialization, chi2 = 0.d0). To each observable (as defined at the beginning of chi2) corresponds an entry in each of the fields value, error, smval, and pull, containing the central observed value, the total (experimental and theoretical) error, the calculated fit value, and the standard deviation, respectively. The function chi2 also contains calls to various other subroutines where the actual observable calculations take place. These are detailed in the following subsections.
Another entry to GAPP is provided through mh.f which computes the probability distribution function of $`M_H`$. The probability distribution function is the quantity of interest within Bayesian data analysis (as opposed to point estimates frequently used in the context of classical methods), and defined as the product of a prior density and the likelihood, $`\mathrm{exp}(\chi ^2/2)`$. If one chooses to disregard any further information on $`M_H`$ (such as from triviality considerations or direct searches) one needs a non-informative prior. It is recommanded to choose a flat prior in a variable defined on the whole real axis, which in the case of $`M_H`$ is achieved by an equidistant scan over $`\mathrm{log}M_H`$. An informative prior is obtained by activating one of the approximate Higgs exclusion curves from LEP 2 near the end of chi2.f. These curves affect values of $`M_H`$ even larger than the corresponding quoted 95% CL lower limit and includes an extrapolation to the kinematic limit; notice that this corresponds to a conservative treatment of the upper $`M_H`$ limit.
Contour plots can be obtained using the routine mncontours from MINUIT. For the cases this fails, some simpler and slower but more robust contour programs are also included in GAPP, but these have to be adapted by the user to the case at hand.
### 2.2 $`\widehat{\alpha }`$, $`\mathrm{sin}^2\widehat{\theta }_W`$, $`M_W`$
At the core of present day electroweak analyses is the interdependence between $`G_F`$, $`M_Z`$, the $`W`$ boson mass, $`M_W`$, and the weak mixing angle, $`\mathrm{sin}^2\theta _W`$. In the $`\overline{\mathrm{MS}}`$ scheme it can be written as ,
$$\widehat{s}^2=\frac{A^2}{M_W^2(1\mathrm{\Delta }\widehat{r}_W)},\widehat{s}^2\widehat{c}^2=\frac{A^2}{M_Z^2(1\mathrm{\Delta }\widehat{r}_Z)},$$
(1)
where,
$$A=\left[\frac{\pi \alpha }{\sqrt{2}G_F}\right]^{1/2}=37.2805(2)\text{ GeV},$$
(2)
$`\widehat{s}^2`$ is the $`\overline{\mathrm{MS}}`$ mixing angle, $`\widehat{c}^2=1\widehat{s}^2`$, and where,
$$\mathrm{\Delta }\widehat{r}_W=\frac{\alpha }{\pi }\widehat{\mathrm{\Delta }}_\gamma +\frac{\widehat{\mathrm{\Pi }}_{WW}(M_W^2)\widehat{\mathrm{\Pi }}_{WW}(0)}{M_W^2}+\text{V + B},$$
(3)
and,
$$\mathrm{\Delta }\widehat{r}_Z=\mathrm{\Delta }\widehat{r}_W+(1\mathrm{\Delta }\widehat{r}_W)\frac{\widehat{\mathrm{\Pi }}_{ZZ}(M_Z^2)\frac{\widehat{\mathrm{\Pi }}_{WW}(M_W^2)}{\widehat{c}^2}}{M_Z^2},$$
(4)
collect the radiative corrections computed in sin2th.f. The $`\widehat{\mathrm{\Pi }}`$ indicate $`\overline{\mathrm{MS}}`$ subtracted self-energies, and V + B denote the vertex and box contributions to $`\mu `$ decay. Although these relations involve the $`\overline{\mathrm{MS}}`$ gauge couplings they employ on-shell gauge boson masses, absorbing a large class of radiative corrections .
$`\mathrm{\Delta }\widehat{r}_W`$ and $`\mathrm{\Delta }\widehat{r}_Z`$ are both dominated by the contribution $`\widehat{\mathrm{\Delta }}_\gamma (M_Z)`$ which is familiar from the RG running of the electromagnetic coupling,
$$\widehat{\alpha }(\mu )=\frac{\alpha }{1\frac{\alpha }{\pi }\widehat{\mathrm{\Delta }}_\gamma (\mu )},$$
(5)
and computed in alfahat.f up to four-loop $`๐ช(\alpha \alpha _s^3)`$. Contributions from $`c`$ and $`b`$ quarks are calculated using an unsubtracted dispersion relation . If $`\mu `$ is equal to the mass of a quark, three-loop matching is performed and the definition of $`\widehat{\alpha }`$ changes accordingly. Pure QED effects are included up to next-to-leading order (NLO) while higher orders are negligible. Precise results can be obtained for $`\mu <2m_\pi `$ and $`\mu m_c`$.
Besides full one-loop electroweak corrections, $`\mathrm{\Delta }\widehat{r}_W`$ and $`\mathrm{\Delta }\widehat{r}_Z`$ include enhanced two-loop contributions of $`๐ช(\alpha ^2m_t^4)`$ (implemented using the analytic expressions of Ref. ) and $`๐ช(\alpha ^2m_t^2)`$ (available as expansions in small and large $`M_H`$); mixed electroweak/QCD corrections of $`๐ช(\alpha \alpha _s)`$ and $`๐ช(\alpha \alpha _s^2m_t^2)`$ ; the analogous mixed electroweak/QED corrections of $`๐ช(\alpha ^2)`$; and fermion mass corrections also including the leading gluonic and photonic corrections.
### 2.3 $`Z`$ decay widths and asymmetries
The partial width for $`Zf\overline{f}`$ decays is given by,
$$\mathrm{\Gamma }_{f\overline{f}}=\frac{N_C^fM_Z\widehat{\alpha }}{24\widehat{s}^2\widehat{c}^2}|\widehat{\rho }_f|\left[14|Q_f|\mathrm{Re}(\widehat{\kappa }_f)\widehat{s}^2+8Q_f^2\widehat{s}^4|\widehat{\kappa }_f|^2\right]$$
$$\times \left[1+\delta _{\mathrm{QED}}+\delta _{\mathrm{QCD}}^{\mathrm{NS}}+\delta _{\mathrm{QCD}}^\mathrm{S}\frac{\widehat{\alpha }\widehat{\alpha }_s}{4\pi ^2}Q_f^2+๐ช(m_f^2)\right].$$
(6)
$`N_C^f`$ is the color factor, $`Q_f`$ is the fermion charge, and $`\widehat{\rho }_f`$ and $`\widehat{\kappa }_f`$ are form factors which differ from unity through one-loop electroweak corrections and are computed in rho.f and kappa.f, respectively. For $`fb`$ there are no corrections of $`๐ช(\alpha ^2m_t^4)`$ and contributions of $`๐ช(\alpha ^2m_t^2)`$ to $`\widehat{\kappa }_f`$ and $`\widehat{\rho }_f`$ are very small and presently neglected. On the other hand, vertex corrections of $`๐ช(\alpha \alpha _s)`$ are important and shift the extracted $`\alpha _s`$ by $`0.0007`$.
The $`Zb\overline{b}`$ vertex receives extra corrections due to heavy top quark loops. They are large and have been implemented in bvertex.f based on Ref. . $`๐ช(\alpha ^2m_t^4)`$ corrections are included, as well, while those of $`๐ช(\alpha ^2m_t^2)`$ are presently unknown. The leading QCD effects of $`๐ช(\alpha \alpha _sm_t^2)`$ and all subleading $`๐ช(\alpha \alpha _s)`$ corrections are incorporated into $`\widehat{\rho }_b`$ and $`\widehat{\kappa }_b`$, but not the $`๐ช(\alpha \alpha _s^2m_t^2)`$ contribution which is presently available only for nonsinglet diagrams .
In Eq. (6), $`\delta _{\mathrm{QED}}`$ are the $`๐ช(\alpha )`$ and $`๐ช(\alpha ^2)`$ QED corrections. $`\delta _{\mathrm{QCD}}^{\mathrm{NS}}`$ are the universal QCD corrections up to $`๐ช(\alpha _s^3)`$ which include quark mass dependent contributions due to double-bubble type diagrams . $`\delta _{\mathrm{QCD}}^\mathrm{S}`$ are the singlet contributions to the axial-vector and vector partial widths which start, respectively, at $`๐ช(\alpha _s^2)`$ and $`๐ช(\alpha _s^3)`$, and induce relatively large family universal but flavor non-universal $`m_t`$ effects . The corrections appearing in the second line of Eq. (6) are evaluated in lep100.f.
The dominant massless contribution to $`\delta _{\mathrm{QCD}}^{\mathrm{NS}}`$ can be obtained by analytical continuation of the Adler $`D`$-function, which (in the $`\overline{\mathrm{MS}}`$ scheme) has a very well behaved perturbative expansion $`1+_{i=0}d_ia_s^{i+1}`$ in $`a_s=\widehat{\alpha }_s(M_Z)/\pi `$ (see the Appendix for details). The process of analytical continuation from the Euclidean to the physical region induces further terms which are proportional to $`\beta `$-function coefficients, enhanced by powers of $`\pi ^2`$, and start at $`๐ช(\alpha _s^3)`$. Fortunately, these terms involve only known coefficients up to $`๐ช(\alpha _s^5)`$, and the only unknown coefficient in $`๐ช(\alpha _s^6)`$ is proportional to the four-loop Adler function coefficient, $`d_3`$. In the massless approximation,
$$\delta _{\mathrm{QCD}}^{\mathrm{NS}}a_s+1.4092a_s^2(0.681+12.086)a_s^3+$$
(7)
$$(d_389.19)a_s^4+(d_4+79.7)a_s^5+(d_5121d_3+3316)a_s^6,$$
and terms of order $`a_s^710^{10}`$ are clearly negligible. Notice, that the $`๐ช(\alpha _s^6)`$ term effectively reduces the sensitivity to $`d_3`$ by about 18%. Eq. (7) amounts to a reorganization of the perturbative series in terms of the $`d_i`$ times some function of $`\alpha _s`$; a similar idea is routinely applied to the perturbative QCD contribution to $`\tau `$ decays .
Final state fermion mass effects of $`๐ช(m_f^2)`$ (and $`๐ช(m_b^4)`$ for $`b`$ quarks) are best evaluated by expanding in $`\widehat{m}_q^2(M_Z)`$ thus avoiding large logarithms in the quark masses. The singlet contribution of $`๐ช(\alpha _s^2m_b^2)`$ is also included.
The dominant theoretical uncertainty in the $`Z`$ lineshape determination of $`\alpha _s`$ originates from the massless quark contribution, and amounts to about $`\pm 0.0004`$ as estimated in the Appendix. There are several further uncertainties, all of $`๐ช(10^4)`$: from the $`๐ช(\alpha _s^4)`$ heavy top quark contribution to the axial-vector part of $`\delta _{\mathrm{QCD}}^\mathrm{S}`$; from the missing $`๐ช(\alpha \alpha _s^2m_t^2)`$ and $`๐ช(\alpha ^2m_t^2)`$ contributions to the $`Zb\overline{b}`$-vertex; from further non-enhanced but cohering $`๐ช(\alpha \alpha _s^2)`$-vertex corrections; and from possible contributions of non-perturbative origin. The total theory uncertainty is therefore,
$$\mathrm{\Delta }\alpha _s(M_Z)=\pm 0.0005,$$
(8)
which can be neglected compared to the current experimental error. If $`\widehat{m}_b`$ is kept fixed in a fit, then its parametric error would add an uncertainty of $`\pm 0.0002`$, but this would not change the total uncertainty (8).
Polarization asymmetries are (in some cases up to a trivial factor 3/4 or a sign) given by the asymmetry parameters,
$$A_f=\frac{14|Q_f|\mathrm{Re}(\widehat{\kappa }_f)\widehat{s}^2}{14|Q_f|\mathrm{Re}(\widehat{\kappa }_f)\widehat{s}^2+8Q_f^2\widehat{s}^4|\widehat{\kappa }_f|^2},$$
(9)
and the forward-backward asymmetries by,
$$A_{FB}(f)=\frac{3}{4}A_eA_f.$$
(10)
The hadronic charge asymmetry, $`Q_{FB}`$, is the linear combination,
$$Q_{FB}=\underset{q=d,s,b}{}R_qA_{FB}(q)\underset{q=u,c}{}R_qA_{FB}(q),$$
(11)
and the hadronic peak cross section, $`\sigma _{\mathrm{had}}`$, is stored in sigmah, and defined by,
$$\sigma _{\mathrm{had}}=\frac{12\pi \mathrm{\Gamma }_{e^+e^{}}\mathrm{\Gamma }_{\mathrm{had}}}{M_Z^2\mathrm{\Gamma }_Z^2}.$$
(12)
Widths and asymmetries are stored in the fields gamma(f), alr(f), and afb(f). The fermion index, f, and the partial width ratios, R(f), are defined in Table 1.
### 2.4 Fermion masses
I use $`\overline{\mathrm{MS}}`$ masses as far as QCD is concerned, but retain on-shell masses for QED since renormalon effects are unimportant in this case. This results in a hybrid definition for quarks. Accordingly, the RG running of the masses to scales $`\mu \widehat{m}_q`$ uses pure QCD anomalous dimensions. The running masses correspond to the functions msrun($`\mu `$), mcrun($`\mu `$), etc. which are calculated in masses.f to three-loop order. Anomalous dimensions are also available at four-loop order , but can safely be neglected. Also needed is the RG evolution of $`\alpha _s`$ which is implemented to four-loop precision in alfas.f.
I avoid pole masses for the five light quarks throughout. Due to renormalon effects, these can be determined only up to $`๐ช(\mathrm{\Lambda }_{\mathrm{QCD}})`$ and would therefore induce an irreducible uncertainty of about 0.5 GeV. In fact, perturbative expansions involving the pole mass show unsatisfactory convergence. In contrast, the $`\overline{\mathrm{MS}}`$ mass is a short distance mass which can, in principle, be determined to arbitrary precision, and perturbative expansions are well behaved with coefficients of order unity (times group theoretical factors which grow only geometrically). Note, however, that the coefficients of expansions involving large powers of the mass, $`\widehat{m}^n`$, are rather expected to be of $`๐ช(n)`$. This applies, e.g., to decays of heavy quarks ($`n=5`$) and to higher orders in light quark mass expansions.
The top quark pole mass enters the analysis when the results on $`m_t`$ from on-shell produced top quarks at the Tevatron are included. In subroutine polemasses(nf,mpole) $`\widehat{m}_q(\widehat{m}_q)`$ is converted to the quark pole mass, mpole, using the two-loop perturbative relation from Ref. . The exact three-loop result has been approximated (for $`m_t`$) by employing the BLM scale for the conversion. Since the pole mass is involved it is not surprising that the coefficients are growing rapidly. The third order contribution is 31%, 75%, and 145% of the second order for $`m_t`$ ($`\mathrm{๐๐}=6`$), $`m_b`$ ($`\mathrm{๐๐}=5`$), and $`m_c`$ ($`\mathrm{๐๐}=4`$), respectively. I take the three-loop contribution to the top quark pole mass of about 0.5 GeV as the theoretical uncertainty, but this is currently negligible relative to the experimental error. At a high energy lepton collider it will be possible to extract the $`\overline{\mathrm{MS}}`$ top quark mass directly and to abandon quark pole masses altogether.
### 2.5 $`\nu `$ scattering
The ratios of neutral-to-charged current cross sections,
$$R_\nu =\frac{\sigma _{\nu N}^{NC}}{\sigma _{\nu N}^{CC}},R_{\overline{\nu }}=\frac{\sigma _{\overline{\nu }N}^{NC}}{\sigma _{\overline{\nu }N}^{CC}},$$
(13)
have been measured precisely in deep inelastic $`\nu `$ ($`\overline{\nu }`$) hadron scattering (DIS) at CERN (CDHS and CHARM) and Fermilab (CCFR). The most precise result was obtained by the NuTeV Collaboration at Fermilab who determined the Paschos-Wolfenstein ratio,
$$R^{}=\frac{\sigma _{\nu N}^{NC}\sigma _{\overline{\nu }N}^{NC}}{\sigma _{\nu N}^{CC}\sigma _{\overline{\nu }N}^{CC}}R_\nu rR_{\overline{\nu }},$$
(14)
with $`r=\sigma _{\overline{\nu }N}^{CC}/\sigma _{\nu N}^{CC}`$. Results on $`R_\nu `$ are frequently quoted in terms of the on-shell weak mixing angle (or $`M_W`$) as this incidentally gives a fair description of the dependences on $`m_t`$ and $`M_H`$. One can write approximately,
$$R_\nu =g_L^2+g_R^2r,R_{\overline{\nu }}=g_L^2+\frac{g_R^2}{r},R^{}=g_L^2g_R^2,$$
(15)
where,
$$g_L^2=\frac{1}{2}\mathrm{sin}^2\theta _W+\frac{5}{9}\mathrm{sin}^4\theta _W,g_R^2=\frac{5}{9}\mathrm{sin}^4\theta _W.$$
(16)
However, the study of new physics requires the implementation of the actual linear combinations of effective four-Fermi operator coefficients, $`ฯต_{L,R}(u)`$ and $`ฯต_{L,R}(d)`$, which have been measured. With the appropriate value for the average momentum transfer, q2, as input, these are computed in the subroutines nuh(q2,epsu\_L,epsd\_L,epsu\_R,epsd\_R) (according to Ref. ), nuhnutev, nuhccfr, and nuhcdhs, all contained in file dis.f. Note, that the CHARM results have been adjusted to CDHS conditions . While the experimental correlations between the various DIS experiments are believed to be negligible, large correlations are introduced by the physics model through charm mass threshold effects, quark sea effects, radiative corrections, etc. I constructed the matrix of correlation coefficients using the analysis in Ref. ,
$$R^{}R_\nu R_\nu R_\nu R_{\overline{\nu }}R_{\overline{\nu }}R_{\overline{\nu }}$$
$$\left(\begin{array}{ccccccc}1.00& 0.10& 0.10& 0.10& 0.00& 0.00& 0.00\\ 0.10& 1.00& 0.40& 0.40& 0.10& 0.10& 0.10\\ 0.10& 0.40& 1.00& 0.40& 0.10& 0.10& 0.10\\ 0.10& 0.40& 0.40& 1.00& 0.10& 0.10& 0.10\\ 0.00& 0.10& 0.10& 0.10& 1.00& 0.15& 0.15\\ 0.00& 0.10& 0.10& 0.10& 0.15& 1.00& 0.15\\ 0.00& 0.10& 0.10& 0.10& 0.15& 0.15& 1.00\end{array}\right).$$
(17)
The effective vector and axial-vector couplings, $`g_V^{\nu e}`$ and $`g_A^{\nu e}`$, from elastic $`\nu e`$ scattering are calculated in subroutine nue(q2,gvnue,ganue) in file nue.f. The momentum transfer, $`\mathrm{๐๐ธ}`$, is currently set to zero . Needed is the low energy $`\rho `$ parameter, rhonc, which describes radiative corrections to the neutral-to-charged current interaction strengths. Together with sin2t0 (described below) it is computed in file lowenergy.f.
### 2.6 Low energy observables
The weak atomic charge, Qw, from atomic parity violation and fixed target $`ep`$ scattering is computed in subroutine apv(Qw,Z,AA,C1u,C1d,C2u,C2d) where Z and AA are, respectively, the atomic number and weight. Also returned are the coefficients from lepton-quark effective four-Fermi interactions which are calculated according to .
These observables are sensitive to the low energy mixing angle, sin2t0, which defines the electroweak counterpart to the fine structure constant and is similar to the one introduced in Ref. . There is significant correlation between the hadronic uncertainties from the RG evolutions of $`\widehat{\alpha }`$ and the weak mixing angle. Presently, this correlation is ignored, but with the recent progress in atomic parity violation experiments it should be accounted for in the future.
An additional source of hadronic uncertainty is introduced by $`\gamma Z`$-box diagrams which are unsuppressed at low energies. At present, this uncertainty can be neglected relative to the experimental precision.
Besides apv, file pnc.f contains in addition the subroutine moller for the anticipated polarized fixed target Mรธller scattering experiment at SLAC. Radiative corrections are included following Ref. .
### 2.7 $`bs\gamma `$
Subroutine bsgamma returns the decay ratio,
$$R=\frac{(bs\gamma )}{(bce\nu )}.$$
(18)
It is given by ,
$$R=\frac{6\alpha }{\pi }\left|\frac{V_{ts}^{}V_{tb}}{V_{cb}}\right|^2\frac{S}{f(z)}\frac{|\overline{D}|^2+A/S+\delta _{NP}+\delta _{EW}}{(1+\delta _{NP}^{SL})(1+\delta _{EW}^{SL})},$$
(19)
where $`|V_{ts}^{}V_{tb}/V_{cb}|^2=0.950`$ is a combination of Cabbibo-Kobayashi-Maskawa matrix elements and $`S`$ is the Sudakov factor . $`\delta _{NP}`$ and $`\delta _{EW}`$ are non-perturbative and NLO electroweak corrections, both for the $`bs\gamma `$ and the semileptonic ($`bce\nu `$) decay rates.
$$\overline{D}=C_7^0+\frac{\widehat{\alpha }_s(\widehat{m}_b)}{4\pi }(C_7^1+V),$$
(20)
is called the reduced amplitude for the process $`bs\gamma `$, and is given in terms of the Wilson coefficient $`C_7`$ at NLO. $`C_7`$ and the other $`C_i`$ appearing below are effective Wilson coefficients with NLO RG evolution from the weak scale to $`\mu =\widehat{m}_b`$ understood. The NLO matching conditions at the weak scale have been calculated in Ref. . $`\overline{D}`$ includes the virtual gluon corrections,
$$V=r_2C_2^0+r_7C_7^0+r_8C_8^0,$$
(21)
so that it squares to a positive definite branching fraction. On the other hand, the amplitude for gluon Bremsstrahlung ($`bs\gamma g`$),
$$\begin{array}{c}A=\frac{\widehat{\alpha }_s(\widehat{m}_b)}{\pi }[C_2^0(C_8^0f_{28}(1)+C_7^0f_{27}(1)+C_2^0f_{22}(1))+\hfill \end{array}$$
$$C_8^0(C_8^0f_{88}(\delta )+C_7^0f_{78}(\delta ))+(C_7^0)^2f_{77}(\delta )],$$
(22)
is added linearly to the cross section. The Wilson coefficient $`C_2^0`$ is defined as in Ref. . It enters only at NLO, is significantly larger than $`C_7^0`$, and dominates the NLO contributions. The parameter $`0\delta 1`$ in the coefficient functions $`f_{ij}`$ characterizes the minimum photon energy and has been set to $`\delta =0.9`$ , except for the first line in Eq. (22) where $`\delta =1.0`$ corresponding to the full cross section. The $`f_{2i}`$ are complicated integrals which can be solved in terms of polylogarithms up to 5th order. In the code I use an expansion in $`z=m_c^2/m_b^2`$ and $`\delta =1.0`$. Once experiments become more precise the correction to $`\delta =0.9`$ should be included.
$`f(z)`$ is the phase space factor for the semileptonic decay rate including NLO corrections . I defined the $`\overline{\mathrm{MS}}`$ mass ratio in $`z=[\widehat{m}_c(\widehat{m}_b)/\widehat{m}_b(\widehat{m}_b)]^2`$ at the common scale, $`\mu =\widehat{m}_b`$, which I also assumed for the factor $`\widehat{m}_b^5`$ multiplying the decay widths. Since I do not reexpand the denominator this effects the phase space function at higher orders. Using the $`๐ช(\alpha _s^2)`$-estimate<sup>2</sup><sup>2</sup>2I computed the $`๐ช(\alpha _s^2)`$ coefficient for comparison only, and did not include it in the code. from Ref. , I obtain for the semileptonic decay width,
$$\mathrm{\Gamma }_{SL}\widehat{m}_b^5f_0(z)\left[1+2.7\frac{\widehat{\alpha }_s(\widehat{m}_b)}{\pi }1.6\left(\frac{\widehat{\alpha }_s}{\pi }\right)^2\right],$$
(23)
where $`f_0(z)`$ is the leading order phase space factor. It is amusing that the coefficients in Eq. (23) are comfortably (and perhaps somewhat fortuitously) small, with the $`๐ช(a_s^2)`$-coefficient even smaller than the one in Ref. where a low scale running mass had been advocated. Moreover, using the prefactor $`\widehat{m}_b^5`$ in the numerator of $`R`$ reduces the size of $`r_7`$ in Eq. (21) and therefore the coefficient $`\kappa (\delta )=f_{77}(\delta )+r_7/2`$ which multiplies the term $`a_s(C_7^{0,\mathrm{eff}})^2`$. I obtain $`2.1<\kappa (\delta )<1.4`$, while with the pole mass prefactor $`M_b^5`$ one would have $`8.7\kappa (\delta )<5.3`$.
## Acknowledgements
It is a pleasure to thank Paul Langacker and Damien Pierce for collaborations on precision analyses. I am grateful to Francesca Borzumati and Paolo Gambino for providing me with parts of their FORTRAN codes.
## Appendix A Uncertainties from perturbative QCD
Writing the perturbative expansion of some quantity in its general form for an arbitrary gauge group, it can easily be decomposed into separately gauge invariant parts. Table 2 shows for some (related) examples that after removing the group theoretical prefactors, all coefficients, $`y_i`$, are strictly of order unity, and that their mean, $`\overline{y}`$, is very close to zero. In particular, there is no sign of factorial growth of coefficients. These observations offer a valuable tool to estimate the uncertainties associated with the truncation of the loop expansion, so I would like to make them more precise.
Assume (for simplicity) that the $`y_i`$ are random draws from some normal distribution with unknown mean, $`\mu `$, and variance, $`\sigma ^2`$. One can show that the marginal distribution of $`\mu `$ follows a Student-t distribution with $`n1`$ degrees of freedom, $`t_{n1}`$, centered about $`\overline{y}`$, and with standard deviation,
$$\mathrm{\Delta }\mu =\sqrt{\frac{\underset{i}{}(y_i\overline{y})^2}{n(n3)}}.$$
(24)
As can be seen from the Table, $`\mu `$ is consistent with zero in all cases, justifying the nullification of the unknown coefficients from higher loops. I next assert that the distribution of $`\sigma `$, conditional on $`\mu =0`$, follows a scaled inverse-$`\chi ^2`$ distribution with $`n`$ degrees of freedom, from which I obtain the estimate,
$$\sigma =\sigma _0\pm \mathrm{\Delta }\sigma =\sqrt{\frac{\underset{i}{}y_i^2}{n2}}\left[1\pm \sqrt{\frac{1}{2(n4)}}\right].$$
(25)
Inspection of the Table shows indeed that $`\sigma `$, as the typical size of a coefficient, is estimated to be $`\stackrel{<}{_{}}๐ช(1)`$.
I now focus on the partial hadronic $`Z`$ decay width. As discussed in Section 2.3, the $`๐ช(\alpha _s^3)`$ term, $`d_2`$, is much smaller than the $`\pi ^2`$ term arising from analytical continuation. This is specifically true for the relevant case of $`n_f=5`$ active flavors, where large cancellations occur between gluonic and fermionic loops. Notice, that the $`D`$-function, in contrast to $`R_{\mathrm{had}}`$, has opposite signs in the leading terms proportional to $`C_A^2C_F`$ and $`C_AC_FT_Fn_F`$. Indeed, the Adler $`D`$-function and the $`\beta `$-function have similar structures regarding the signs and sizes of the various terms (see Table 2), and we do expect large cancellations in the $`\beta `$-function. The reason is that it has to vanish identically in the case of $`N=4`$ supersymmetry. Ignoring scalar contributions this case can be mimicked by setting $`T_Fn_f=2C_A`$ (there are 2 Dirac fermions in the $`N=4`$ gauge multiplet) or $`n_f=12`$ for QCD, which is of the right order. In fact, all known QCD $`\beta `$-function coefficients become very small for some value of $`n_f`$ between 6 and 16. We therefore have a reason to expect that similar cancellations will reoccur in the $`d_i`$ at higher orders. As a $`1\sigma `$ error estimate for $`d_3`$, I suggest to use the largest known coefficient ($`3\times 0.71`$) times the largest group theoretical prefactor in the next order ($`C_A^3C_F`$) which results in
$$d_3=0\pm 77.$$
(26)
With Eq. (7) and $`\widehat{\alpha }_s(M_Z)=0.120`$ one can absorb all higher order effects into the $`๐ช(\alpha _s^4)`$-coefficient of $`R_{\mathrm{had}}`$, $`r_3^{\mathrm{eff}}=81\pm 63`$. This shifts the extracted $`\alpha _s`$ from the $`Z`$ lineshape by +0.0005 and introduces the small uncertainty of $`\pm 0.0004`$.
The argument given above does certainly not apply to the quenched case, $`n_f=0`$, and indeed $`d_2(n_f=0)`$ is about $`73\%`$ of the $`\pi ^2`$ term, i.e., large and positive. In the case of $`n_f=3`$, which is of interest for the precision determination of $`\alpha _s`$ from $`\tau `$ decays, $`d_2`$ is about $`38\%`$ of the $`\pi ^2`$ term. If one assumes that the same is true of $`d_3`$, one would obtain $`d_3(n_f=3)=60`$. Estimates based on the principles of minimal sensitivity, PMS, or fastest apparent convergence, FAC, yield $`d_3(n_f=3)=27.5`$ so there might be some indications for a positive $`d_3(n_f=3)`$. In any case, all these estimates lie within the uncertainty in Eq. (26) and we will have to await the proper calculation of the $`๐ช(\alpha _s^4)`$-coefficient to test these hypotheses. Note, that the current $`\tau `$ decay analysis by the ALEPH Collaboration uses $`d_3=50\pm 50`$ which is more optimistic.
The analogous error estimate for the five-loop $`\beta `$-function coefficient yields,
$$\beta _4=0\pm 579.$$
(27)
To get an estimate for the uncertainty in the RG running of $`\widehat{\alpha }_s`$, I translate Eq. (27) into
$$\beta _3=\beta _3\pm \frac{\widehat{\alpha }_s(\mu _0)}{\pi }\beta _4,$$
(28)
where $`\mu _0`$ is taken to be the lowest scale involved. This overestimates the uncertainty from $`\beta _4`$, thereby compensating for other neglected terms of $`๐ช(\alpha _s^{n+4}\mathrm{ln}^n\mu ^2/\mu _0^2)`$. For the RG evolution from $`\mu =m_\tau `$ to $`\mu =M_Z`$ this yields an uncertainty of $`\mathrm{\Delta }\alpha _s(M_Z)=\pm 0.0005`$. Conversely, for fixed $`\alpha _s(M_Z)=0.120`$, I obtain $`\widehat{\alpha }_s(\widehat{m}_b)=0.2313\pm 0.0006`$, $`\widehat{\alpha }_s(m_\tau )=0.3355\pm 0.0045`$, and $`\widehat{\alpha }_s(\widehat{m}_c)=0.403\pm 0.011`$, where I have used $`\widehat{m}_b=4.24`$ GeV and $`\widehat{m}_c=1.31`$ GeV. For comparison, the ALEPH Collaboration quotes an evolution error of $`\mathrm{\Delta }\alpha _s(M_Z)=\pm 0.0010`$ which is twice as large. I emphasize that it is important to adhere to consistent standards when errors are estimated. This is especially true in the context of a global analysis where the precisions of the observables enter as their relative weights.
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# Topological Inflation in Supergravity
## I Introduction
Superstring theories compactified on $`(3+1)`$-dimensional space-time have many discrete symmetries in the low-energy effective Lagrangian . A spontaneous breakdown of such discrete symmetries creates topological defects, i.e. domain walls, in the early universe . If the vacuum-expectation value (VEV) of a scalar field $`\phi `$ is larger than the gravitational scale $`M_G2\times 10^{18}`$ GeV, the region inside the wall undergoes inflationary expansion and eventually becomes the present whole universe . If the universe is open at the beginning, it expands and the spontaneous breakdown of the symmetries always takes place at some epoch in the early universe. It has been recently argued that the quantum creation of the open universe may take place with appropriate continuation from the Euclidean instanton . Thus, topological inflation is a natural consequence of the dynamics of the system, and it does not require any fine-tuning of initial conditions for the beginning universe. Furthermore, it does not cause the โgraceful exitโ problem and the universe becomes homogeneously radiation dominated after reheating.
A simple and interesting model for topological inflation was proposed in the framework of supergravity .<sup>*</sup><sup>*</sup>*Other topological inflation models were studied in the superstring inspired models . However, it was not explicitly shown whether topological inflation really takes place. In this paper, we perform a numerical analysis on the above model and show that topological inflation indeed occurs in a wide range of parameter space. We also show that the condition for successful inflation depends not only on the superpotential, which determines the vacuum expectation value (VEV) of inflaton $`\phi `$, but also on the form of its Kรคhler potential. We in fact find that the required VEV can be as small as $`\phi 1\times M_G`$, which is far below the lower bound of $`\phi =\eta _{cr}1.7M_G`$ derived in Ref. .
## II Topological Inflation Model
We begin with the topological inflation model proposed in Ref. , which is based on $`R`$-invariant supergravity. The gravitational scale $`M_G`$ is set to be unity below. In this model the superpotential for the inflaton superfield $`\varphi (x,\theta )`$ is given by
$$W=v^2X(1g\varphi ^2).$$
(1)
Here, we have imposed $`U(1)_R\times Z_2`$ symmetry and omitted higher-order terms for simplicity. Under the $`U(1)_R`$ we assume
$$X(\theta )e^{2i\alpha }X(\theta e^{i\alpha }),\varphi (\theta )\varphi (\theta e^{i\alpha }).$$
(2)
We also assume that the superfield $`X`$ is even and $`\varphi `$ is odd under the $`Z_2`$. This discrete $`Z_2`$ symmetry is an essential ingredient for the topological inflation . In the above superpotential (1), we always take $`v^2`$ and $`g`$ to be real constants without loss of generality.
The $`R`$\- and $`Z_2`$-invariant Kรคhler potential is given by
$$K(\varphi ,X)=|X|^2+|\varphi |^2+k_1|X|^2|\varphi |^2+\frac{k_2}{4}|X|^4+\frac{k_3}{4}|\varphi |^4+\mathrm{},$$
(3)
where $`k_1`$, $`k_2`$ and $`k_3`$ are constants of order unity.
The potential of a scalar component of the superfields $`X(x,\theta )`$ and $`\varphi (x,\theta )`$ in supergravity is given by
$$V=e^K\left\{\left(\frac{^2K}{z_iz_j^{}}\right)^1D_{z_i}WD_{z_j^{}}W^{}3|W|^2\right\}(z_i=\varphi ,X),$$
(4)
with
$$D_{z_i}W=\frac{W}{z_i}+\frac{K}{z_i}W.$$
(5)
This potential yields an $`R`$-invariant vacuum
$`X=0,\varphi ={\displaystyle \frac{1}{\sqrt{g}}}{\displaystyle \frac{\eta }{\sqrt{2}}},`$ (6)
at which the potential energy vanishes. Here, the scalar components of the superfields are denoted by the same symbols as the corresponding superfields. If $`\eta `$ is larger than the critical value $`\eta _{\mathrm{cr}}`$ which will be discussed in the next section, the topological inflation occurs.
For $`|X|`$ and $`|\varphi |1`$, we approximately rewrite the potential (4) as
$$Vv^4|1g\varphi ^2|^2+v^4(1k_1)|\varphi |^2k_2v^4|X|^2.$$
(7)
If $`k_21`$, $`X`$ field quickly settles down to the origin and we set $`X=0`$ in our analysis taking $`k_21`$. For $`g>0`$, we can identify the inflaton field $`\phi (x)/\sqrt{2}`$ with the real part of the field $`\varphi (x)`$ since the imaginary part of $`\varphi (x)`$ has a positive mass and the real part has a negative mass. Because the positive mass of the imaginary part is larger than the size of the negative mass, the imaginary part is irrelevant for the inflation dynamics and hence we neglect it. Then, we obtain a potential for the inflaton for $`\phi 1`$:
$$V(\phi )v^4\frac{\kappa }{2}v^4\phi ^2,$$
(8)
where
$$\kappa 2g+k_11.$$
(9)
The slow-roll condition for the inflaton $`\phi `$ is satisfied for $`0<\kappa <1`$ and $`0\phi 1`$.We can always take $`\phi `$ positive since we have the $`Z_2`$ symmetry $`(\varphi \varphi )`$. The Hubble parameter during the inflation is given by $`Hv^2/\sqrt{3}`$. The scale factor of the universe increases by a factor of $`e^N`$ when the inflaton $`\phi `$ rolls slowly down the potential from $`\phi _N`$ to $`1`$. The $`e`$-fold number $`N`$ is given by
$`N{\displaystyle \frac{1}{\kappa }}\mathrm{ln}\phi _N.`$ (10)
The amplitude of primordial density fluctuations $`\delta \rho /\rho `$ due to this inflation is written as
$$\frac{\delta \rho }{\rho }\frac{1}{5\sqrt{3}\pi }\frac{v^2}{\kappa \phi _N}2.0\times 10^5.$$
(11)
The normalization is given by the data of anisotropies of the cosmic microwave background radiation (CMB) by the COBE satellite . Since the $`e`$-fold number $`N`$ corresponding to the COBE scale is about $`60`$, which leads to
$`v2.3\times 10^2\sqrt{\kappa }e^{\frac{\kappa N}{2}}|_{N=60}1.8\times 10^33.6\times 10^4,`$ (12)
for $`0.02\kappa 0.1`$.
The interesting point on the above density fluctuations is that it results in the tilted spectrum whose spectrum index $`n_s`$ is given by
$$n_s12\kappa .$$
(13)
We may expect a possible deviation from the Harrison-Zeldvich scale-invariant spectrum $`n_s=1`$. Observational constraint on $`n_s`$ is $`|n_s1|<0.2`$ , which implies $`0<\kappa <0.1`$.
After inflation ends, the inflaton $`\phi `$ may decay into ordinary particles as discussed in Ref. and the reheating temperature is low enough to avoid overproduction of gravitinos which are thermally produced at the reheating epoch. Recently, nonthermal production at the preheating stage was found to be important in some inflation models . For the present model, non-thermal production of gravitinos at the preheating phase is roughly estimated as
$`\left({\displaystyle \frac{n_{3/2}}{s}}\right)_{\mathrm{non}\mathrm{TH}}{\displaystyle \frac{m_\phi ^3}{v^4/T_R}}10^{14}\left({\displaystyle \frac{T_R}{10^{10}\mathrm{GeV}}}\right),`$ (14)
where $`m_\phi (v^2)`$, $`n_{3/2}`$, $`s`$, and $`T_R`$ are the mass of the inflaton, the number density of gravitinos, entropy density and reheating temperature, respectively. This is much less than the thermal production given by $`(n_{3/2}/s)_{\mathrm{TH}}10^{11}(T_R/10^{10}\mathrm{GeV})`$ and hence we can neglect the nonthermal production of gravitinos.
## III Numerical Simulation
We perform numerical simulations to decide whether topological inflation takes place in supergravity and determine the condition for successful topological inflation. For the purpose, we follow time evolution of the domain wall and investigate whether it inflates or not. Since we consider a planar domain wall whose width is order of the horizon scale, we cannot adopt the Friedmann Robertson-Walker metric. Instead, we assume that the spacetime has a reflection symmetry of the coordinate $`x`$ perpendicular to the wall and adopt the metric given by
$$ds^2=dt^2+A^2(t,|x|)dx^2+B^2(t,|x|)(dy^2+dz^2),$$
(15)
where $`A`$ and $`B`$ correspond to the scale factors in the direction of $`x`$ and $`y`$-$`z`$, respectively. If the inflation occurs and the proper width of the wall becomes much larger than the horizon scale, $`A`$ and $`B`$ expand as $`ABe^{Ht}`$ (as shown later) and the universe approaches the de Sitter spacetime. Thus, we examine whether the proper width of the wall becomes much larger than the horizon scale. Once it is realized, the universe expands exponentially .
We adopt the numerical technique developed in Ref. . The Einstein-Hilbert action is given by
$$S=d^4x\sqrt{g}\left[\frac{1}{2}\frac{1}{2}(_\mu \phi )^2V(\phi )\right].$$
(16)
Variating the above action with respect to the metric $`g_{\mu \nu }`$, we obtain the Einstein equations,
$$G_{\mu \nu }_{\mu \nu }\frac{1}{2}=T_{\mu \nu },$$
(17)
where $`T_{\mu \nu }`$ is the energy-momentum tensor,
$$T_{\mu \nu }=_\mu \phi _\nu \phi g_{\mu \nu }\left[\frac{1}{2}(_\mu \phi )^2+V(\phi )\right].$$
(18)
Variation with respect to the scalar field $`\phi `$ gives the equation of motion for the scalar field $`\phi `$,
$$\mathrm{}\phi =\frac{dV(\phi )}{d\phi }.$$
(19)
In order to make it easier to follow time evolution of the system, we choose certain combinations of Einstein equations, which read
$`G_0^0`$ $`=`$ $`๐ฆ_2^2(2๐ฆ3๐ฆ_2^2){\displaystyle \frac{2B^{\prime \prime }}{A^2B}}{\displaystyle \frac{B^2}{A^2B^2}}+{\displaystyle \frac{2A^{}B^{}}{A^3B}}`$ (20)
$`=`$ $`{\displaystyle \frac{\dot{\phi }^2}{2}}+{\displaystyle \frac{\phi ^2}{2A^2}}+V(\phi ),`$ (21)
$`{\displaystyle \frac{1}{2}}G_{01}`$ $`=`$ $`๐ฆ_{2}^{2}{}_{}{}^{}+{\displaystyle \frac{B^{}}{B}}(3๐ฆ_2^2๐ฆ)`$ (22)
$`=`$ $`{\displaystyle \frac{1}{2}}\dot{\phi }\phi ^{},`$ (23)
$`{\displaystyle \frac{1}{2}}(G_1^1+G_2^2+G_3^3G_0^0)`$ $`=`$ $`\dot{๐ฆ}(๐ฆ_1^1)^22(๐ฆ_2^2)^2`$ (24)
$`=`$ $`\dot{\phi }^2V(\phi ),`$ (25)
$`_2^2{\displaystyle \frac{1}{2}}G_0^0`$ $`=`$ $`\dot{๐ฆ_2^2}+{\displaystyle \frac{B_{}^{}{}_{}{}^{2}}{2A^2B^2}}{\displaystyle \frac{3}{2}}(๐ฆ_2^2)^2`$ (26)
$`=`$ $`{\displaystyle \frac{\dot{\phi }^2}{4}}+{\displaystyle \frac{\phi ^2}{4A^2}}{\displaystyle \frac{V(\phi )}{2}},`$ (27)
where an overdot denotes the time derivative and a dash the spatial derivative. $`๐ฆ_{ij}`$ are the extrinsic curvature tensors of constant time hypersurface, given by
$$๐ฆ_1^1=\frac{\dot{A}}{A},๐ฆ_2^2=๐ฆ_3^3=\frac{\dot{B}}{B},$$
(28)
and $`๐ฆ`$ denotes its trace $`๐ฆ๐ฆ_i^i`$. The equation of motion for the scalar field $`\phi `$ becomes
$$\ddot{\phi }๐ฆ\dot{\phi }\frac{\phi ^{\prime \prime }}{A^2}\left(\frac{A^{}}{A}+\frac{2B^{}}{B}\right)\frac{\phi ^{}}{A^2}+\frac{dV(\phi )}{d\phi }=0.$$
(29)
We set the initial condition for numerical simulations. First, we consider an initial configuration of a domain wall. For the convenience of numerical calculations, we take only the region $`|x|/\delta 2`$ We have confirmed that the results do not change even if we take wider ranges of the direction $`x`$ (e.g. $`|x|/\delta =3`$). where $`\delta (=\frac{\eta }{\sqrt{2}v^2})`$ is the width of the domain wall and $`\eta `$ is the VEV of $`\phi `$. We impose the free boundary condition, that is, $`\phi ^{}=A^{}=B^{}=0`$ at the boundaries $`x/\delta =2`$ and $`x/\delta =2`$. Then, we adopt the following initial configuration for the domain wall so that the gradient of the field disappears at the boundaries,
$$\phi (t=0,x)=\{\begin{array}{cc}\eta \left[\frac{x}{\delta }\frac{5}{4}\left(\frac{8}{15}\frac{x}{\delta }\right)^3+\frac{3}{8}\left(\frac{8}{15}\frac{x}{\delta }\right)^5\right]\left(0\frac{x}{\delta }\frac{15}{8}\right),\hfill & \\ \eta \left(\frac{15}{8}\frac{x}{\delta }2\right),\hfill & \end{array}$$
(30)
with $`\phi (t=0,x)=\phi (t=0,x)`$ for $`2x/\delta 0`$. This is a deformed version of the static domain wall solution in a flat spacetime, $`\phi _{\mathrm{flat}}=\eta \mathrm{tanh}(x/\delta )`$. The function $`\phi (t=0,x)`$ is decided so as to satisfy the following three conditions: (1) it is a fifth-order odd polynomial function of $`x`$, (2) the first term coincides with that of the expansion of $`\phi _{\mathrm{flat}}`$, (3) it is smooth at $`x/\delta =15/8`$, that is, $`\phi ^{}=\phi ^{\prime \prime }=0`$ at $`x/\delta =15/8`$. Also, $`\dot{\phi }`$ is set to be 0.
Next, on the initial hypersurface, we determine $`A,B,๐ฆ_2^2,๐ฆ`$ so as to satisfy the Hamiltonian constraint (21) and the momentum constraint (23). We have freedom for the initial hypersurface to have homogeneous and isotropic curvature, which automatically satisfies the momentum constraint (23). This choice leads to
$$\frac{๐ฆ}{3}=๐ฆ_1^1=๐ฆ_2^2=\text{โnegativeโ const},$$
(31)
where โnegativeโ implies that the universe is in an expanding phase. Furthermore, we can take the conformally flat spatial gauge, $`A=B`$, on the initial hypersurface and set $`A=B=1`$ at $`x=0`$. Finally, we determine the negative value of $`๐ฆ`$. Since we adopt the reflection symmetry of the coordinate $`x`$ and the free boundary condition, the condition $`A^{}=B^{}=0`$ at $`x=0`$ and $`x=\pm 2\delta `$ must be satisfied. $`๐ฆ`$ is determined so that the Hamiltonian constraint (21) satisfies the above conditions.
Now the initial settings are completed and hence we have only to follow the time evolution of five variables, $`A,B,๐ฆ,๐ฆ_2^2`$, and $`\phi `$. Note that we have introduced five variables, $`A,B,๐ฆ,๐ฆ_2^2`$, and $`\phi `$ though only three variables are independent. This is partly because the second order differential equations have been reduced to the first order differential equations. Moreover, as for time evolution of $`๐ฆ_2^2`$, we use Eq.(27) only at $`x=0`$ and acquire the value at $`x0`$ by integrating Eq.(23) in the direction of $`x`$ in order to avoid numerical instability.
When inflation takes place, $`A`$ and $`B`$ grow exponentially so that the proper distance from the domain wall core ($`x=0`$) also increases exponentially. In order to see whether this happens or not, we follow time evolution of the width of the wall for a given potential $`V(\phi )`$.
To fix the potential $`V(\phi )`$, we first consider the Kรคhler potential with only terms up to the fourth order,
$$K(\varphi ,X)=|X|^2+|\varphi |^2+k_1|X|^2|\varphi |^2+\frac{k_2}{4}|X|^4+\frac{k_3}{4}|\varphi |^4.$$
(32)
Then, the Lagrangian density is given by
$`(\varphi )`$ $`=`$ $`_{kin}(\varphi )V(\varphi )`$ (33)
$`=`$ $`(1+k_3|\varphi |^2)_\mu \varphi ^\mu \varphi ^{}v^4|1g\varphi ^2|^2{\displaystyle \frac{\mathrm{exp}\left(|\varphi |^2+\frac{k_3}{4}|\varphi ^4|\right)}{1+k_1|\varphi |^2}}.`$ (34)
Here we have set $`X=0`$. Identifying the inflaton field $`\phi (x)/\sqrt{2}`$ with the real part of the field $`\varphi (x)`$, the Lagrangian density becomes
$`(\phi )`$ $`=`$ $`_{kin}(\phi )V(\phi )`$ (35)
$`=`$ $`{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{1}{2}}k_3\phi ^2\right)(_\mu \phi )^2v^4\left(1{\displaystyle \frac{g}{2}}\phi ^2\right)^2{\displaystyle \frac{\mathrm{exp}\left(\frac{1}{2}\phi ^2+\frac{k_3}{16}\phi ^4\right)}{1+\frac{k_1}{2}\phi ^2}},`$ (36)
with the VEV $`\eta =\sqrt{2/g}`$. In the present model we have four free parameters $`k_1,k_2,k_3`$ and $`g`$. However, $`k_2(1)`$ only works as a stabilizer of the $`X`$ field as explained before and it is not important for the dynamics of topological inflation itself. Once the $`X`$ field is stabilized at $`X=0`$, the potential $`V(\phi )`$ does not depend on $`k_2`$. $`k_3`$ is almost irrelevant for the dynamics of $`\varphi `$ field and only changes its VEV due to the redefinition of $`\phi `$ with a canonical kinetic term. Then, we set $`k_3=0`$ first and later consider the case of nonzero $`k_3`$. Thus, we have only two relevant parameters, $`k_1`$ and $`g`$. The potential $`V(\phi )`$ has a pole at $`\phi =\sqrt{2/|k_1|}`$ for $`k_1<0`$. But, we are only interested in the dynamics of $`\phi `$ up to the VEV $`\eta `$ so that there is no problem if $`|k_1|<g`$ for $`k_1<0`$.
We introduce dimensionless quantities, $`\overline{\phi }=\phi /M_G,\overline{x}=xH(x=0),\overline{t}=tH(x=0),\overline{๐ฆ}_{ij}=๐ฆ_{ij}/H(x=0)`$, and $`\overline{\delta }=\delta H(x=0)=1/(\sqrt{3g})`$ where $`H(x=0)=v^2/\sqrt{3}`$. As the first step, we consider the simplest case of $`k_1=0`$ and $`g=0.5(`$i.e., $`\eta =2.0)`$, which leads to the spectral index $`n=1`$. Time evolution of the domain wall is depicted in Fig. 1. The vertical axis represents the value of the scalar field. The horizontal axis represents the proper distance from the domain wall core. As time elapses, the domain wall (the region for $`\phi 0.8`$) expands and topological inflation really takes place. As shown in Figs. 2 and 3, once the domain wall expands enough, the scale factors $`A`$ and $`B`$ increase at the same rate, that is, $`๐ฆ_1^1๐ฆ_2^2H(x=0)`$ inside the wall. Thus, the universe expands exponentially and approaches the de Sitter spacetime.
We consider the dependence on $`\kappa `$ for the VEV fixed, $`\eta =2.0(g=0.5)`$. The result is depicted in Fig. 4. For large $`\kappa `$, the domain wall cannot inflate enough. $`\kappa `$ controls the mass scale of the scalar field $`\phi `$ near the origin. As $`\kappa `$ becomes larger, the scalar field $`\phi `$ rolls down faster so that only the small region of the original domain wall earns the vacuum energy and cannot overcome the gradient energy.
The dependence on $`v`$ is studied, which determines the energy scale of the domain wall. The result with the same parameters in Fig. 4 except for $`v=3.6\times 10^4`$ is shown in Fig. 5. The results are quite the same and have no dependence on $`v`$. This can be interpreted as follows: First, the dependence on $`v`$ only appears through the Hubble parameter during the inflation given by $`Hv^2/\sqrt{3}`$. But, since the width of the domain wall $`\delta 1/(\sqrt{3g}H)`$, the rough criterion $`\delta >H^1`$ becomes independent of $`H`$, that is, $`v`$. Also, rewriting the potential as $`V(\phi )\lambda (\phi ^2\eta ^2)`$ with $`\lambda =(v/\eta )^4`$, the independence of $`v`$ is equivalent to that of $`\lambda `$. The independence of $`v`$ is also found in Ref. .
We now consider the cases of the nonzero $`k_3`$. Since the kinetic term $`_{kin}`$ is not canonical in these cases, we define the scalar field $`\mathrm{\Phi }`$ with the canonical kinetic term as
$$\mathrm{\Phi }=_0^\phi ๐\phi \sqrt{1+\frac{1}{2}k_3\phi ^2}=\{\begin{array}{cc}\frac{1}{2}\phi \sqrt{1+\frac{1}{2}k_3\phi ^2}+\frac{\mathrm{arcsinh}\left(\sqrt{\frac{k_3}{2}}\phi \right)}{\sqrt{2k_3}}\left(k_30\right),\hfill & \\ \frac{1}{2}\phi \sqrt{1\frac{1}{2}|k_3|\phi ^2}+\frac{\mathrm{arcsin}\left(\sqrt{\frac{|k_3|}{2}}\phi \right)}{\sqrt{2|k_3|}}\left(k_3<0\right).\hfill & \end{array}$$
(37)
The results for the same parameters in Fig. 1 except for $`k_3=\pm 0.3`$ are depicted in Fig. 6. The positive $`k_3`$ encourages the occurrence of topological inflation while the negative $`k_3`$ discourages. This is because the VEV of $`\mathrm{\Phi }`$ is larger than $`\eta `$ for the positive $`k_3`$ while smaller for the negative $`k_3`$.
Finally, we discuss the criterion for the VEV of $`\phi `$ for successful topological inflation within $`0\kappa 0.1`$. In previous analyses, we search for only the parameter region satisfying $`|k_1|<g`$ for negative $`k_1`$ because of the appearance of a pole of the potential. The condition $`|k_1|<g`$ leads to $`g<1+\kappa `$ so that $`\eta >1.41(1.35)`$ for $`\kappa =0.0(0.1)`$, where we have confirmed the occurrence of topological inflation. In order to obtain the lower limit of $`\eta `$, we add to the Kรคhler potential the sixth order terms,
$$\mathrm{\Delta }K=l_1|X|^2|\varphi |^4+l_2|X|^4|\varphi |^2+\frac{l_3}{9}|X|^6+\frac{l_4}{9}|\varphi |^6,$$
(38)
where there are four parameters but only $`l_1`$ is relevant. If we take $`l_1`$ satisfying $`l_1>g(g1)`$, the pole smaller than the VEV does not appear in the whole $`k_1g`$ parameter space. The result for smaller $`\eta `$ with $`\kappa =0.0`$ is depicted in Fig. 7. We find the critical value of the breaking scale, $`\eta _{cr}0.95(1.00)`$ for $`\kappa =0.0(0.1)`$, which corresponds to $`\varphi 1/\sqrt{2}`$.
## IV Conclusions and discussions
We have studied a topological inflation in supergravity. First, we have shown that topological inflation really takes place in supergravity. Also, the criterion of successful topological inflation depends not only on the breaking scale of the discrete symmetry but also on the mass of the inflaton near the origin. This is because the inflaton rolls down rapidly from the origin if its mass is large. For a very flat case favored by the observation of the spectral index, $`n_s10.8`$ (i.e., $`0<\kappa <0.1`$), we have found that the critical breaking scale $`\eta _{cr}`$ becomes as small as $`M_G`$, which is smaller than the critical value, $`\eta _{cr}1.7M_G`$ observed in Ref. . Finally we have discussed the primordial spectrum produced by the topological inflation. In general, the topological inflation predicts the tilted spectrum $`n_s<1`$ depending on $`\kappa `$.<sup>ยง</sup><sup>ยง</sup>ยงIt is possible to produce more exotic spectrum including blue one. This is due to the exponential blow of the potential, which is significant for the region $`\phi \stackrel{>}{}M_G`$ and makes the potential more complex than the simple double-well potential.
The present topological inflation model is free from the thermal and nonthermal overproduction of gravitinos since the reheating temperature can be as low as $`10^8`$ GeV. Furthermore, as pointed out in Ref. , this model is consistent with a leptogenesis scenario in which heavy Majorana neutrinos are produced in the inflaton decay and successive decays of the Majorana neutrinos result in lepton asymmetry enough to explain the observed baryon asymmetry in the present universe.
### ACKNOWLEDGMENTS
M.Y. is grateful to T. Kanazawa and J. Yokoyama for useful discussions. M.K. and T.Y. are supported in part by the Grant-in-Aid, Priority Area โSupersymmetry and Unified Theory of Elementary Particlesโ(No. 707). N.S. and M.Y. are partially supported by the Japanese Society for the Promotion of Science.
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# The Labusch Parameter of a Driven Flux Line Lattice in YBa2Cu3O7 Superconducting Films
## 1 Introduction
The investigation of effects directly related to true phase trasitions of the flux line lattice (FLL) in high-temperature superconductors (HTS) is a current topic of research. These studies are not only difficult to perform experimentally due to the necessary resolution but also their interpretation is sometimes not straightforward due to pinning effects. In the quasi three dimensional (3D) HTS YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> (Y123) an anomalous enhancement of the pinning force appears to be associated with the โpeak effectโ (PE) in the dc critical current measured using the voltage criterium . In contrast to conventional low-$`T_c`$ superconductors, this PE occurs at fields much below the upper critical field, near the thermally activated depinning line . Within the classical argument of Pippard the PE in a 3D superconductor with a high density of pinning centers occurs if at high enough fields or temperatures the shear modulus or rigidity of the FLL strongly drops. This softening of the FLL observed increasing temperature and/or magnetic field should overwhelm the usual decrease of the pinning force in order to have an effective enhancement of the pinning force. Within this picture the correlation between the PE and a reduction of the rigidity of the FLL is appealing.
The experimental evidence for a correlation between the PE and the โmeltingโ of the FLL is, however, far from being clear. As noted in Ref. the sharp enhancement of resistance with temperature in Y123 crystals, interpreted as a first-order melting transition , is precisely opposed to the enhancement of pinning expected within Pippardโs picture. The results of a remarkable experimental study characterising the PE in a clean and anisotropic superconductor suggested that much of the transport data interpreted as melting in HTS may be related to nonequilibrium dynamical steady states of the FLL and not to true phase transitions .
Studies of the PE and its dependence on the magnetic field history in Y123 thin films and crystals using the vibrating reed technique casted some doubts on the interpretation of the PE line as a melting line. It was suggested that the onset of the PE is due to the thermal creation of FLL dislocations and that some of the melting lines may be observations of the PE using different criteria . The strongly nonlinear response of vibrating superconductors near the PE is qualitatively consistent with the dynamic phase diagram measured in .
New features recently found in Y123 crystals by transport experiments reveal a PE at high fields in the vortex solid region of the phase diagram . This PE joins smoothly the โ believed to be โ first-order melting line at high temperatures and suggests a correlation with the previously observed PE at low fields . The rather temperature independence of this new PE at high fields was speculated to be the analog of the second magnetization peak observed in Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub> (Bi2212) HTS . Recently published measurements on Y123 crystals up to 27 T reveal that the temperature dependence of the field at the second magnetization peak is qualitatively similar to that observed in Bi2212 crystals . Depending on the crystal properties, the second magnetization peak can be observed even in the same temperature range, i.e. at $`T<50`$K, as in Bi2212 HTS . Recent experimental studies cast strong doubts about the interpretation of the second magnetization peak in Bi2212 HTS based on an pinning enhancement and related to a phase transition of the FLL . The origin of the so-called โpre-meltingโ peaks in the critical current density obtained by transport measurements is therefore still unclear .
As the vibrating reed results in Y123 crystals and thin films showed, the observed PE depends on the magnetic history of the sample as well as on the driving force provided by the reed amplitude, which in general is $`150`$nm. The observed history effect means that the magnitude of the PE depends if the sample was cooled in or without a field. We note that the driving forces in a typical vibrating reed experiment are extremelly low โ shielding currents lower than $`10^3`$ A/m<sup>2</sup> โ which is an advantage in comparison to other methods used to study pinning properties. Therefore, it could be shown that at low enough vibrating amplitudes (low driving force) the PE vanishes in a Y123 film when measured as a function of temperature at constant field . A strong driving current dependence of the PE has been also observed in 2H-NbSe<sub>2</sub> crystal . In this crystal the dissipation near the upper critical field strongly increases with the transport current, in agreement with the observed behavior in , although the PE still remains .
Current driven order, disorder and depinning of the FLL in 2H-NbSe<sub>2</sub> crystals have been studied in detail in Refs. . The influence of the probing current on the first-order transition has been studied in Y123 crystals . In contrast to those studies, the technique used in this work allows us to study the change of pinning using probing and driving currents substantially smaller than the critical one. The vibrating reed measurements provide the space derivative of the pinning force at small distances of the potential minimum. As we will show below, at low temperatures the measurements as a function of a driving dc current reveal an irreversible enhancement of the pinning of the FLL. Loading a vibrating superconductor with a transport current allows us to resolve very small changes in the elastic coupling as a function of applied current far below the vortex depinning transition $`T<T_D(B_a)`$ and applied forces much less than the critical ones. In this paper we show that a metastable state with a distinctly higher level of pinning than that one obtained after field cooling procedure without a driving force can be reached by applying a sufficiently large driving current. Our work shows that the field-cooled state is not a state of minimum potential energy for the FLL, at least for Y123 HTS, and therefore different metastable states may be obtained depending on the used experimental procedure to study pinning-related properties.
## 2 Experimental Details
The measuring setup is schematically depicted in figure 1. The samples (films with substrate) were glued onto a silicon single crystal host reed. The dimensions of the host reed used in this study were (length $`\times `$ width $`\times `$ thickness) $`l\times w\times d=5\times 1\times 0.1`$ mm<sup>3</sup>. The reed was T-shaped and its back-side was covered with 15 nm gold layer for the capacitive detection of the reed motion. The reed driving force supplied by a sinusoidal voltage was kept constant throughout the measurements. Detailed description of the vibrating reed technique and the typical electronic arrangement are published elsewhere . Our system allows us to measure automatically the resonance frequency $`\nu =\omega /2\pi `$ of the reed as a function of field, temperature or driving current density $`J_T`$.
The YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> films were produced by pulsed laser deposition technique with a rotating substrate . The films were deposited on a thin (0.1 mm) SrTiO<sub>3</sub> substrate to minimize the total sample weight glued to the host reed. The results presented in this work were obtained with a film of dimensions: $`d_p400`$ nm, $`w_p1.2`$ mm, and $`l_p1.6`$ mm.
The virgin field-cooled state was achieved by cooling down in a magnetic field through the superconducting transition temperature to the temperature of the measurements without appling a current load. The maximum current load applied in the measurements was 100 mA. The measurements have been performed with the field applied parallel to the CuO<sub>2</sub> planes of the film ($`\theta =0^{}\pm 1^{}`$) and at $`\theta =30^{}`$. No differences between the results obtained for the two different angles were observed.
## 3 Flux line tension and elastic pinning correction in vibrating superconductors
The resonance frequency enhancement of a vibrating superconductor as a function of applied magnetic field $`B_a`$ is due to the restoring force as a result of a pinned FLL. This behavior is well understood and can be quantitatively accounted for by the theory developed in (for a review see Refs. ). We would like to outline here the main results of the theory and restrict ourselves to the configuration depicted in Fig. 1 having the sample glued on a host reed and the magnetic field applied parallel (or slightly tilted) to the length of the host reed and to the $`ab`$-planes of the films.
A superconducting sample with pinned FLL vibrating in a magnetic field distorts the field arround the superconductor and consequently a positive restoring force increases the resonance frequency of the cantilever. In this case, for a sample with dimension $`l_pw_p>>d_p`$ the magnetic line tension can be approximated by
$$P\frac{\pi w_p}{4d_p}w_pd_p\frac{B_a^2}{\mu _0}\frac{l_p}{l}.$$
(1)
The condition of large enough pinning is satisfied for
$$l_p>>\mathrm{\Lambda }_{44},$$
(2)
where $`\mathrm{\Lambda }_{44}=\lambda _{44}(\pi w_p/4d_p)^{1/2}`$ is the effective Campbell penetration depth for tilt waves and $`\lambda _{44}=(c_{44}/\alpha )^{1/2}`$ with $`c_{44}=B_a^2/\mu _0`$ the FLL tilt modulus. The Labusch parameter or elastic coupling between FLL and the superconducting matrix is
$$\alpha (B_a,T)=\frac{^2}{s^2}U(s,B_a,T)|_{s0}=\left|\frac{F_P}{s_p}\right|,$$
(3)
where $`U`$ is the mean value of the interaction energy per unit length (or pinning potential) between vortices and pinning centers, $`F_p=B_aJ_c`$ is the volume pinning force, $`s`$ is the flux line displacement relative to the superconducting matrix, $`s_p`$ is the characteristic range of the pinning potential and $`J_c`$ the critical current density.
The field dependence of the resonance frequency depends on the ratio of the magnetic $`W_MPl`$ to mechanical $`W_RI\omega ^2(0)`$ energies and the pinning strength of the FLL. For $`X=1.33(Pl/I\omega ^2(0))<1`$ and if (2) holds, the ideal enhancement (infinite pinning case) of the resonance frequency is given by
$$\omega _i^2(B_a)\omega ^2(0)Pl/I,$$
(4)
where $`I`$ is the inertia moment of the cantilever with sample. This is obtained experimentally at low enough temperatures and applied fields where Eq. (4) holds. In the present work we determined $`I5.8\times 10^{11}`$ kg m<sup>2</sup> for the used cantilever.
The field dependence of the resonance frequency is in general influenced by the finite pinning which is taken into account by a correction term $`\omega _{pin}^2`$. Non-local and end effects due to the finite length of the sample are also taken into account through the correction terms $`\omega _{nl}^2`$ and $`\omega _{end}^2`$, along with a correction term due to the finite damping of the reed $`\mathrm{\Gamma }^2`$. The field dependence of the resonance frequency can be then approximated by
$$\omega ^2(B_a)\omega ^2(0)\frac{Pl}{I}\omega _{pin}^2\omega _{nl}^2\omega _{end}^2\mathrm{\Gamma }^2.$$
(5)
In general, for the geometry of the used samples $`\omega _{pin}^2>>\omega _{nl}^2,\omega _{end}^2,\mathrm{\Gamma }^2`$.
The correction term due to the finite pinning is given by
$$\omega _{pin}^20.16\omega ^2(0)X\tau \left(\frac{\mathrm{\Lambda }_{44}}{l_p}\right)^2\left(\frac{F(X^{})}{1+r/2}+\frac{pl_pG(X)}{w_p+p\mathrm{\Lambda }_{44}\tau ^{1/2}}\right),$$
(6)
where $`F(X)`$ and $`G(X)`$ are tabulated numerical functions , $`X^{}=X/(1+4\beta X)`$ with $`\beta w_p^2/2\pi ^2l_p^2`$ ; $`p0.54\mathrm{ln}(0.71w_p/d_p)`$, $`r=(2X)^{1/2}\mathrm{\Lambda }_{44}\tau ^{1/2}/l_p`$, and $`\tau 1.08`$ for the case of not too strong pinning, when compressional waves penetrate the superconductor entirely ($`\mathrm{\Lambda }_{44}>d_p`$).
The deviation of the measured resonance frequency from the infinite pinning value $`\omega _i^2`$ is inversely proportional to $`\alpha `$ via $`\omega _{pin}^2\mathrm{\Lambda }_{44}^2\alpha ^1`$ and therefore it allows us to numerically calculate the Labush parameter as a function of $`B_a`$. We should note, however, that Eq. (6) is only an approximation for the finite pinning correction and it does not provide the correct value of $`\alpha `$ for too large ($`\mathrm{\Lambda }<d_p`$) or too weak ($`\omega _{pin}^20.1Pl/I`$) pinning.
In the temperature region of interest 90 K $`T>50`$K the Labusch parameter increases approximately as $`B_a^2`$ having a value $`\alpha (B_a=1\mathrm{T},T=73\mathrm{K})4\times 10^{17}`$ N/m<sup>4</sup> in good agreement with that obtained in Ref. for YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> films.
## 4 Results and discussion
Figure 2 shows the resonance frequency change $`\nu (I_T)\nu (0)`$ as a function of the applied driven current $`I_T`$ at constant field and temperature. The sample was cooled through the superconducting transition to 70 K in a field of 5 T applied parallel to the CuO<sub>2</sub> planes. The current dependence of the virgin state is given by the close symbols in Fig. 2. We clearly observe that the resonance frequency increases with applied current indicating, as explained in the last section, an increase of the elastic constant. The resonance frequency reaches a maximum at a current $`I_010^2`$A which corresponds to $`10^4J_c(70`$K,5 T). At larger currents the resonance frequency slightly decreases. Decreasing the current, the resonance frequency follows the virgin curve down to the current where the resonance frequency shows the maximum. For lower currents the resonance frequency remains practically current independent. This irreversible behavior is observed in the temperature and field range when the resonance frequency increases with current in the field-cooled virgin state.
Figure 3 shows the normalized elastic coupling as a function of the applied current density at two applied fields and at different constant temperatures, calculated from the measured resonance frequency change using Eqs. (5) and (6). All curves were measured starting from the field-cooled and zero-current virgin state. Different behavior is observed at different temperatures of the measurement. At low enough temperatures (e.g. $`T<60`$K) $`\alpha `$ remains independent of the applied current up to high current densities (a slight decrease of $`\alpha `$ is observed at the highest currents). The current independence of $`\alpha `$ at low $`T`$ is interpreted as due to the relatively large pinning which is not overwhelmed in the available applied current density range.
Increasing temperature (e.g. $`T>60`$K) , a maximum in $`\alpha `$ as a function of applied current clearly develops, see Fig. 3. The increase in $`\alpha `$ reaches a maximum at a certain temperature (e.g., at $`T70`$K at $`B_a=1`$T). For higher temperatures, eventually $`\alpha `$ decreases with applied current. In general, we can assume that the observed temperature dependence of $`\alpha (I_T)`$ results from the competing influence of pinning and shear modulus $`c_{66}(B_a,T)`$ of the FLL at low temperatures and thermally activated flux flow at high temperatures. We note that the absence of a โpeak effectโ in $`\alpha `$ as a function of current at high enough temperatures does not imply the absence of the commonly observed โpeak effectโ in $`\alpha `$ as a function of temperature at a constant driving force.
Figure 4 shows the normalized elastic constant as a function of the normalized applied current for a field of 1 T at different temperatures. This behavior is similar to that observed for other fields, see Fig. 3. It is interesting to note that the peak in $`\alpha `$ is reached at a current of the order of $`10^3`$ of the critical current density. The critical current density is obtained from the $`IV`$ characteristic curves and defined at the voltage $`V=2\mu `$V. We note that the rearrangement of the FLL is clearly observed already at currents several orders of magnitude smaller. We would like also to stress that the changes in $`\alpha `$ are observed at a field (or temperature) much below the corresponding depinning line $`B_a470[`$T\]$`(1t)^{1.4}`$ with $`t=T_D(B_a)/T_D(0)`$ measured at $`I_T=0`$, $`T_D(0)=88`$K . The shift of the depinning line with transport current is less than 2 K at the maximum applied current of 100 mA and at fields $`B_a5`$T .
The maximum increase in $`\alpha `$ with driving current observed in our samples is of the order of $`30\%`$, see Fig. 4, in the measured field range ($`B_a1`$T). This increase agrees with the low-limit increase of pinning for defective FLL obtained by computer simulation studies . The FLL defects produced by the driving current are probably dislocations. These are also created by thermal disorder and may lead to an enhancement in the pinning strength when measured as a function of temperature.
The results clearly indicate that the field cooled FLL is in a metastable, relatively ordered state. Upon field cooling the superconducting sample through the superconducting and (thermally activated) depinning transitions, the flux lines have to accomodate themselves in usually a random pinning centers matrix taking into account the balance between two contributions: vortex-vortex and vortex-pinning centers interactions. For our case of strong applied fields $`B_a>>B_{c1}`$ we can assume that the vortex-vortex interaction overwhelms the pinning contribution at temperatures above and just below the depinning line. If the driving currents used to measure the pinning properties are small enough in order not to perturbate the FLL, regions of the FLL remain frozen in a metastable state. The observation of a peak effect will depend therefore on the magnetic history and the perturbation amplitude used to detect the pinning strength. The clear amplitude \- or driving force - and magnetic history dependence of the peak effect observed near the depinning transition in Ref. points out the importance of defects in the FLL to produce the peak effect. The driving current applied at constant field and temperature causes the arrangement of the FLL into a more disordered state and the elastic coupling increases. For intermediate fields, we expect that the higher the applied field the smaller will be the increase of $`\alpha `$ produced by the driving current since the vortex-vortex interaction increases whereas the current driven production of the defects in the FLL weakens. This behavior is observed experimentally, see Fig. 3. After reaching a maximum coupling at a current $`I_0`$, the decrease of the resonance frequency can be understood as due to the decrease of the current dependent pinning barrier. This pinning barrier is field and temperature dependent and, therefore, the observed decrease in $`\alpha `$ depends also on these parameters.
Our results are in good agreement with the conclusions obtained by transport and decoration experiments in 2H-NbSe<sub>2</sub> crystals . In this work the authors found that the number of defects in the vortex structure increases when the FLL starts to depin. The FLL becomes nearly defect-free for large enough applied currents because the vortex-vortex interaction overwhelms. Our results show, see Fig. 4, that at large driven currents $`J_T>10^2J_c`$ the elastic coupling decreases and therefore the pinning. At high enough driven currents we expect a depinned and obviously more ordered FLL than any lattice in the FC state without driving forces.
At this point we would like to compare the published data on the peak effect measured with the vibrating reed technique in Y123 films and crystals as well as the peak effect and current driven (re)organization of the FLL observed in 2H-NbSe<sub>2</sub> crystals . Vibrating-reed results in Y123 thin films showed that the enhancement of $`\alpha `$ as a function of temperature at constant applied field and in the field-cooled (FC) state is observed only if the driving force (or vibrating amplitude of the superconducting sample) is large enough. In contrast to the FC state, the zero-field cooled (ZFC) state shows a weak or no signature of a pinning enhancement near the depinning line.
For the case the peak effect in $`\alpha `$ is observed near the depinning temperature $`T_D(B_a)`$, the results indicate that $`\alpha `$ in the ZFC state and near $`T_D(B_a)`$ is smaller than $`\alpha `$ in the FC state . This observation agrees qualitatively with the smaller critical current density measured in the ZFC state in 2H-NbSe<sub>2</sub> at temperatures below the peak effect . From this observation, the authors in Refs. proposed that the ZFC state has a more ordered FLL than the FC state, arguing that in the $`I_c`$ measurements in the ZFC state the vortices enter the sample with large velocities which lead to a more ordered FLL state. We note, however, that in the case of the vibrating reed, the perturbation to the FLL is kept much smaller than that used for $`I_c`$ measurements. Therefore, we think that, in general, the ZFC state has a less homogeneous field distribution in the sample than in the FC state. Increasing temperature, an effective smaller $`\alpha (B_a,T)`$ is measured in the ZFC state because: (a) the effective field in the sample is smaller and (b) there is a change of the effective pinning potential due to the field gradient. The differences between the ZFC and FC states vanish as soon as the thermal creation of FLL defects and the flux creep compensate the inhomogeneous field distribution.
## 5 Conclusion
We have shown that the field-cooled state in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> superconducting films can be driven to a lower potential equilibrium state at fields and temperatures far below the depinning transition. This new equilibrium state of the FLL has a larger elastic coupling. Our results indicate that after field cooling without driving forces, the static FLL is in a metastable, relatively ordered state. Its pinning can be enhanced by the current driven creation of defects in the FLL. Our results support the recently published interpretation on the โmeltingโ transition that argues that the FLL does not melt at the depinning line but the opposite, it decouples from the superconducting matrix becoming more ordered . This ordering is partially retained in the FLL after FC through the depinning line without driving forces.
Acknowledgements: We would like to acknowledge M. Lorenz for providing us with the YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> superconducting films and Y. Kopelevich for a careful reading of the manuscript and fruitful discussions. This work is supported by the Deutsche Forschungsgemeinschaft within the โInnovationskolleg Phรคnomene an der Miniaturisierungsgrenzeโ (Project H, DFG IK 24/B1-1) and the German-Israeli Foundation for Scientific Research and Development (Grant G-0553-191.14/97).
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# Non-time-orthogonal reference frames in the theory of relativity
## I. Introduction
Although addressed by Einstein<sup>1</sup><sup>1</sup>1 Albert Einstein, โDie Grundlage der allgemeinen Relativit tstheorie,โ Ann. Phys.(Leipzig) 49, 769-822 (1916). <sup>,</sup><sup>2</sup><sup>2</sup>2 Albert Einstein, The Evolution of Physics, (Simon-Schuster, New York, 1938), pp. 226-234. <sup>,</sup><sup>3</sup><sup>3</sup>3 John Stachel, โEinstein and the Rigidly Rotating Disk,โ Chap. 1 in H. Held, General Relativity and Gravitation (Plenum, New York, 1980), pp. 1-15. and others<sup>4</sup><sup>4</sup>4 P. Ehrenfest, โGleichf mrige Rotation starrer K rper und Relativit theorie,โ Phys. Z. 10, 918-928 (1909). <sup>,</sup><sup>5</sup><sup>5</sup>5 Philip Franklin., โThe meaning of rotation in the special theory of relativityโ, Proc. Nat. Acad. Sci. USA 8(9), 265-268 (1922). <sup>,</sup><sup>6</sup><sup>6</sup>6 M.G. Trocheries, โElectrodynamics in a rotating frame of reference,โ Phil. Mag. 40(310), 1143-1154 (1949). <sup>,</sup><sup>7</sup><sup>7</sup>7 Hyoitiro Takeno, โOn relativistic theory of rotating diskโ, Prog. Theor. Phys. 7(4), 367-376 (1952). in the first half of the twentieth century, the relativistically rotating reference frame continues to be a research topic of interest, and, in fact, has often generated significant discussion and debate.<sup>8</sup><sup>8</sup>8 Oyvind Gron., โRotating frames in special relativity analyzed in light of a recent article by M. Straussโ, Int. J. Theor. Phys. 16(8), 603-614 (1977). <sup>,</sup><sup>9</sup><sup>9</sup>9 Nicholas Sama, โOn the Ehrenfest paradox,โ Am. J. of Phys. 40, 415-418 (1972). <sup>,</sup><sup>10</sup><sup>10</sup>10 Gerald N. Pellegrini and Arthur R. Swift, โMaxwellโs equations in a rotating medium: Is there a problem?,โ Am. J. Phys. 63(8), 694-705 (1995). <sup>,</sup><sup>11</sup><sup>11</sup>11 Thomas A. Weber, โMeasurements on a rotating frame in relativity, and the Wilson and Wilson experiment,โ Am. J. Phys. 65 , 946-953 (1997). <sup>,</sup><sup>12</sup><sup>12</sup>12 Charles T. Ridgely., โApplying relativistic electrodynamics to a rotating material mediumโ, Am. J. Phys. 66 (2) 114-121 (1998). <sup>,</sup><sup>13</sup><sup>13</sup>13 Robert D. Klauber, โComments regarding recent articles on relativistically rotating framesโ, Am. J. Phys. 67(2), 158-159, (1999). <sup>,</sup><sup>14</sup><sup>14</sup>14 Thomas A. Weber, โResponse to โComments regarding recent articles on relativistically rotating frameโ \[Am. J. Phys. 67(2), 158 (1999)โ, Am. J. Phys. 67(2), 159-160 (1999). <sup>,</sup><sup>15</sup><sup>15</sup>15 Robert V. Krotkov, Gerald N. Pellegrini, Norman C. Ford, and Arthur R. Swift, โRelativity and the electric dipole moment of a moving, conducting, magnetized sphere,โ Am. J. Phys. 67(6), 493-498 (1999). <sup>,</sup><sup>16</sup><sup>16</sup>16 Hrvoje Nikolic, โEhrenfest paradox, non-time-orthogonal frames, and local observers,โ Los Alamos Nat Lab, xxx.lanl.gov, paper num gr-qc/9904078 (1999). (References cited in this section are not exhaustive.) Some<sup>17</sup><sup>17</sup>17 Franco Selleri, โNoninvariant one-way speed of light and locally equivalent reference frames,โ Found. Phys. Lett., 10, 73-83 (1997). <sup>,</sup><sup>18</sup><sup>18</sup>18 J. Paul Wesley, Classical Quantum Theory (Benjamin Wesley, Blumberg, Germany, 1996), pp 202-203. <sup>,</sup><sup>19</sup><sup>19</sup>19 R. Anderson, I. Vetharaniam, and G. E. Stedman, โConventionality of Synchronisation, Gauge Dependency, and Test Theories of Relativity,โ Phys. Rep., 295 (3&4), 93-180 (1998). , who have considered the thought experiment described below and/or the Sagnac experiment<sup>20</sup><sup>20</sup>20 G. Sagnac, Comptes Rendus, 157, 708 (1913). <sup>,</sup><sup>21</sup><sup>21</sup>21 E.J. Post, โSagnac effect,โ Mod. Phys. 39, 475-493 (1967). , have proposed the possible existence of a preferred reference frame, which is somehow disguised within our present understanding of special relativity.
The present article does not support this view, and in the final analysis, is consonant with the special and general theories of relativity. It does, however, illustrate one way in which the extant theory does appear to be self-contradictory. This seeming internal incongruity is resolved by analysis of the non-time-orthogonal nature (i.e., time is not orthogonal to at least one spatial dimension) of the rotating frame. In the process, however, non-time-orthogonal (NTO) frames are found to exhibit certain unique and somewhat surprising characteristics. Though some of these characteristics are not what might presently be considered typically relativistic, they are not only in agreement with empirical evidence, but explain what has heretofore been considered an anomalous experimental result.
## II. Rotating Frames in Thought and Practice
### A. Simple Gedanken Experiment
Consider the rotating disk of Figure 1 with a rim mounted light source capable of emitting light pulses in both directions along the disk circumference. A cylindrical mirror, polished side facing inward, is mounted on the rim as well.
Figure 1. Thought Experiment
At time T = 0, an observer attached to the light source triggers the emission of two very short light pulses, one in the clockwise direction, the other in the counter clockwise direction. Each light wave packet is 1<sup>o</sup> long at the rim radius, i.e., it is 1/360 of the rim circumference.
In the lab frame each light pulse has speed c, and each is reflected by the cylindrical mirror, such that it travels a circular route with radius equal to the disk radius r. As the two light pulses are travelling, the disk is rotating c.c.w. at an $`\omega `$high enough to produce relativistic rim velocities. Since the observer mounted on the light source is moving c.c.w. as well, the c.w. moving light pulse reaches him before the c.c.w. pulse does. From the point of view of the lab, this conclusion is inescapable, and since we are talking about separate detectable physical events, it must be true from the point of view of the observer on the disk as well.
The disk observer knows that each light signal traveled the same distance in his frame (he set up the experiment). He also knows that one of them took less time to travel that distance than the other. Hence, the only conclusion he can make is that the speed of light on the disk was greater for the c.w. travelling pulse than for the c.c.w. travelling pulse. Due to symmetry, it also seems inevitable that this conclusion holds true locally as well as globally.
Note this โexperimentโ did not entail any wave interference measurements. It dealt simply with arrival times of short photon wave packets and had nothing to do with the wave nature of those packets.
The conundrum, of course, is that, according to relativity theory, the speed of light is invariant and always equal to c in all directions, no matter what frame one is in. (Those who believe that the general theory of relativity may say otherwise are referred to the Appendix.) The problem is compounded when one considers that most analyses of relativistic rotation utilize an infinite series of local Lorentz frames instantaneously co-moving with the rotating frame at a given radius from the center of rotation. Calculations of things like spatial distance around the circumference are then made by summing (i.e., integrating) infinitesimal quantities (e.g., dx) from all of the local co-moving Lorentz frames.
But in Lorentz frames the speed of light is always c. So in effect, when one uses such frames one is assuming the local speed of light is c, isotropic and invariant. But, as we have illustrated in our thought experiment, this does not appear to be the case for rotating frames. And since the invariance of the speed of light is one postulate upon which the theory of relativity rests, one must immediately re-evaluate not only the suitability of such approaches, but also the very theory itself.
### B. The Sagnac Experiment
The analysis of the previous section literally reeks of a preferred frame (โabsoluteโ space), and proponents of such a thing commonly cite the Sagnac experiment (see Figure 2) as empirical proof<sup>22</sup><sup>22</sup>22 See, for example, reference 18. .
In the Sagnac experiment, a light beam is emitted radially from the center of a rotating disk and is split by a half-silvered mirror M at radius r. From there one part of the beam is reflected by mirrors appropriately placed on the disk such that it travels in one direction effectively around the circumference. The other half of the beam travels the same route over the same distance, but in the opposite direction. The beams then meet up again and are reflected back to the center where interference of the two beams results in a fringing, i.e., a displacement of one light wave with respect to the other.
If the speed of light on the disk were invariant, then as the rotational velocity of the disk increased, the fringe pattern would remain unchanged, similar to what one finds in the Michelson-Morley experiment. However, when this was done by Sagnac<sup>20</sup> and others<sup>21</sup> who have repeated his experiment, the fringe pattern did in fact change, indicating a dependence of the rotating frame speed of light on both direction and rotational speed.
The test results have experimental accuracy only to first order in v/c = $`\omega `$r/c, but they indicate that the speed of a light ray tangent to the circumference measured locally on the disk is equal to<sup>23</sup><sup>23</sup>23 Post (ref. 21) presents the Sagnac results in terms of the fringe shift (his equation (1)). Using (1) of the present paper, one can derive (1) in Post (use A = $`\pi `$r<sup>2</sup>), and vice versa.
$$|v_{light}|c\pm \omega r,$$
(1)
where the approximately equal sign implies accuracy to first order, and the sign in front of the last term depends on the relative direction of the rim tangent and light ray velocities. This same relationship can be easily derived using the logic of the previous section. It is obviously disconcerting, as (1) looks far more Galilean/Newtonian in nature than relativistic.
Some have attributed the Sagnac phenomenon to wave effects. Mashoon et al<sup>24</sup><sup>24</sup>24 Bahram Mashhoon, Richard Neutze, Mark Hannam, Geoffrey E. Stedman, โObservable frequency shifts via spin-rotation couplingโ, Phys. Lett. A, 249, 161-166 (1998). , for example, consider it to be a โmanifestation of the coupling of orbital angular momentum of a particle .. to rotationโ. For a wave this perturbation in the Hamiltonian induces a phase shift such as that measured in the Sagnac experiment. In somewhat similar fashion, Anandan<sup>25</sup><sup>25</sup>25 J. Anandan, โSagnac effect in relativistic and nonrelativistic physics,โ Phys. Rev. D 24(2), 338-346 (1981). asserts โ.. this effect depends only on the frequency of the beams โฆโ.
However, such analyses fail to answer the question raised by our thought experiment, which was not based on wave interference, but solely on arrival times of very short wave packets. As (1) is readily deduced from that thought experiment, it appears a more fundamental reconciliation with the theory of relativity is needed.
## III. Transformation to the Rotating Frame
We will adopt what is presently the most widely<sup>10,11,</sup><sup>26</sup><sup>26</sup>26 C. Moller, The Theory of Relativity (Clarendon Oxford, 1969). <sup>,</sup><sup>27</sup><sup>27</sup>27 L. D. Landau, and E. M. Lifshitz, The Classical Theory of Fields (Addison-Wesley, Reading, 1962), pp. 271-298. <sup>,</sup><sup>28</sup><sup>28</sup>28 Oyvind Gron,., โRelativistic description of a rotating diskโ, Am. J. of Phys. 43(10), 869-876 (1975). <sup>,</sup><sup>29</sup><sup>29</sup>29 Robert D. Klauber, โNew perspectives on the relatively rotating disk and non-time-orthogonal reference framesโ, Found. Phys. Lett. 11(5), 405-443 (1998). On page 421 Klauber lists assumptions upon which the transformation is based. <sup>,</sup><sup>30</sup><sup>30</sup>30 Adler, R., Bazin, M., and Schiffer, M., Introduction to General Relativity (McGraw-Hill, New York, 1975), 2<sup>nd</sup> ed., p. 121-122. , though not universally<sup>17,</sup><sup>31</sup><sup>31</sup>31 Post (reference 21) noted the transformation is determinable experimentally only to first order and suggested the presence of a factor he termed $`\gamma `$on the right hand side of (2.a) and in front of the second term on the right hand side of (2.c). He considered that this factor could be unity, the Lorentz contraction factor, or perhaps something else. Ehrenfest, Franklin, and Trocheries (references 4,5,6) considered transformations other than (2), which never proved completely satisfactory. Although Stedman (reference 19 and in private communication) contends there is latitude in the choice of transformation due to gauge freedom, he does consider the same transformation to the rotating frame used by Selleri (ref. 17). , accepted transformation (see (2.a-d) below) between the lab and rotating frames. This coordinate transformation, where upper case coordinates represent the inertial frame K, lower case denote the rotating frame k, and the axis of rotation is coincident with both the Z and z axes, is
$$cT=ct(2.a),$$
(2)
$$R=r(2.b),$$
$$\mathrm{\Phi }=\varphi +\omega t(2.c),$$
$$Z=z(2.d).$$
$`\omega `$ is the angular velocity of the disk, and t, the coordinate time for the rotating system, is the proper time of a standard clock located at the origin of the rotating coordinate frame, i.e., it is equivalent to any standard clock at rest in K. Note that t is only a coordinate. It is merely a label and cannot be expected to equal proper time at any given point on the disk (except, of course, at r=0).
The metric for the rotating system can be found from the line element for the standard cylindrical coordinate system of the Minkowski space K
$$ds^2=c^2dT^2+dR^2+R^2d\mathrm{\Phi }^2+dZ^2$$
(3)
Assuming ds is invariant, one can find dT, dR, d$`\mathrm{\Phi }`$, and dZ from (2), and insert into (3) to obtain the line element, and hence the metric, of the coordinate grid in k.
$$ds^2=c^2\left(1r^2\omega ^2/c^2\right)dt^2+dr^2+r^2d\varphi ^2+2r^2\omega d\varphi dt+dz^2$$
(4)
$$=g_{\alpha \beta }dx^\alpha dx^\beta $$
For a fixed point in the rotating frame (i.e., $`dr=d\varphi =dz=0`$) with $`ds^2=c^2d\tau ^2`$inserted into $`(4)`$, we find the local proper time on the disk to be
$$d\tau =\left(1r^2\omega ^2/c^2\right)^{1/2}dt=\left(1r^2\omega ^2/c^2\right)^{1/2}dT.$$
(5)
That this is the familiar Lorentz time dilation factor (with $`v=\omega r`$), in full accord with numerous cyclotron experiments, supports the contention that (2.a-d) is indeed the correct form of the transformation.
## IV. Speed of Light in NTO Frames
Note that the metric $`g_{\alpha \beta }`$of $`(4)`$ is non-diagonal, and hence the rotating frame is non-time-orthogonal (NTO). The presence of the non-zero line element term in $`d\varphi dt`$indicates that in 4D spacetime, the time axis is not orthogonal to the spatial axis for the circumferential direction. As we shall see, this has profound implications for measurement of the speed of light.
### A. Analytical Determination of the Speed of Light
Consider the path of a photon travelling in the circumferential direction, e.g., along the disk rim. For light, ds = 0 (see (3)), and for the path considered, dr=dz=0. Inserting these values in $`(4)`$, solving the resultant quadratic equation for $`d\varphi `$, dividing by dt, and multiplying by r, one obtains a coordinate velocity for the photon as seen from the rotating frame
$$v_{light,circum,coord}=\frac{rd\varphi }{dt}=r\omega \pm c.$$
(6)
To find the physical velocity one would measure with standard rods and clocks mounted on the rotating frame, we must use (5) to convert coordinate time dt in (6) to physical time on standard rotating clocks at radius r. This yields
$$v_{light,circum,phys}=\frac{rd\varphi }{d\tau }=\frac{1}{\sqrt{1r^2\omega ^2/c^2}}\frac{rd\varphi }{dt}=\frac{r\omega \pm c}{\sqrt{1v^2/c^2}},$$
(7)
which is the exact form of the approximate (first order) relationship (1) deduced from the Sagnac experiment and our thought experiment.
Note that had we started with a time orthogonal frame, there would have been no $`d\varphi dt`$ type off diagonal term in $`(4)`$, and hence no $`r\omega `$ term in (7). (For example, take the lab frame with $`\omega =0`$ in (7). Alternatively, set the off diagonal term to zero in $`(4)`$ and follow the steps used to derive (7).)
It is important to note that the non-relativistic looking, Newtonian-like, velocity addition relationship of (7) is a direct result of non-time-orthogonality. In NTO frames light speed is neither invariant nor isotropic. In the time orthogonal frames more typically dealt with in special and general relativity, it is invariant and isotropic. Further the degree of anisotropy in the speed of light is directly correlated with the magnitude of the NTO off diagonal term in the frameโs line element.
### B. Picturing NTO Frames and Light Speed
The effect of non-time-orthogonality can be visualized with the aid of Figure 3, which shows the time and spatial circumferential axes of both the lab frame k and the rotating frame K at radius r (=R). We set up our coordinates such that $`\varphi =z=\mathrm{\Phi }=Z=0`$, and note that only very small changes in $`\varphi ,\mathrm{\Phi },t,`$and T are considered. The thinner (orthogonal) lines represent the lab frame axes (upper case); the thicker (NTO) lines, the rotating frame (lower case).
Note that the effect of the time dilation factor appearing in $`(4)`$, (5), and (7) is masked, because we are plotting the time axis of the rotating frame as our coordinate time t (time on a clock at r=0). This is not the physical time on a standard clock fixed to that frame at radius $`r0`$. If one wishes, one can simply multiply the coordinate time t by the (second order) time dilation factor at any point in the discussion of this section to obtain exact relationships for physical quantities. For the present, however, the primary emphasis is on first order effects, as reflected in relationship (1).
From the line element $`(4)`$ and basic trigonometry (generalized Pythagorean theorem) one can determine the slope of the rotating frame time axis. Alternatively, set the RHS of (3) equal to RHS of the first line in $`(4)`$ with$`d\varphi =dr=dz=dR=dZ=0`$, use (2.a) and (2.b), and solve for $`Rd\mathrm{\Phi }/cdT`$.
MN is the path of a light ray and has null path length, i.e., ds = 0. Observe that for a given amount of coordinate time (which is the same in both k and K, i.e., c$`\mathrm{\Delta }`$T = c$`\mathrm{\Delta }`$t), the light ray travels a certain spatial distance l in k, but a greater spatial distance L in K. Hence the speed of light measured in k is less than that in K, and this corresponds with the plus sign before the c in (6). For a light ray in the opposite direction (minus sign in (6)) one can show graphically (with a light ray in the second quadrant of Figure 3 at right angle to MN) that the corresponding l distance is greater than L, and hence the velocity for that ray would be greater in k than in K.
Given that the slope of MN is unity, L = c$`\mathrm{\Delta }`$T. Dividing this by $`\mathrm{\Delta }`$T, one gets the speed of light in K as c. The k time axis has slope c/v = c/$`\omega `$r, so
$$l=Lc\mathrm{\Delta }T\left(\omega r/c\right).$$
(8)
Dividing this by $`\mathrm{\Delta }`$T, one arrives at (6) for the coordinate speed of light in k (with the plus sign for c since light ray MN is traveling in the direction of disk rotation). Note that larger values of $`\omega `$ or r mean a greater degree of non-time-orthogonality and light speed discrepancy from c.
Figure 4, presented for completeness, depicts a co-moving inertial (Lorentz) frame K<sub>1</sub>, having instantaneous velocity equal to the circumferential velocity of the rotating frame at r. Note that in both inertial frames K and K<sub>1</sub>, time is (Minkowski spacetime) orthogonal to 3D space and the speed of light ray MN equals c.
In general one can conclude that non-invariance and degree of anisotropy for the local, physically measurable speed of light are directly dependent on the slope of the time axis relative to the 3D space axes. All frames for which time is orthogonal to space have isotropic light speed equal to c. All NTO frames have anisotropic light speed not equal to c.
Note that both the rotating frame and the instantaneously co-moving Lorentz frame exhibit some of the same properties. They both have the same time dilation as seen from the lab, as well as the same 4D interval ds between events. However, because of differences in time orthogonality, each has different measured speeds for light. Hence, a co-moving Lorentz frame can not simply be assumed to be an appropriate local surrogate for the rotating frame.
## V. Comparison with Experiment
### A. Michelson-Morley Revisited
Although the analysis herein is supported by the Sagnac experiment, one might ask why then did Michelson and Morley not find the speed of light in the directions of galactic and solar orbital rotation different from that in other directions? The answer, the author submits, is that bodies in gravitational orbits follow geodesics, i.e., they are in โfree fallโ. That is, they are in locally inertial, time orthogonal frames and therefore obey Lorentzian mechanics.
Consider a planet in orbit around a star that doesnโt rotate relative to distant stars (i.e., one solar day equals one year). An observer inside a windowless box (similar to Einsteinโs enclosed gedanken elevator) on that planet could do tests (pails of water, Foucault pendulum, Coriolis effects, etc.) to determine that she is not rotating. Hence her frame, the frame of the planet, would be time orthogonal, and her experiments would also find the speed of light to be isotropic, equal to c, and independent of the planetโs orbital velocity.
If, however, her planet were rotating relative to distant stars at $`\omega `$, her measurements would detect a rotation rate of $`\omega `$, independent of orbital angular velocity around her own star. Her frame would then be NTO, with all of the concomitant phenomena described in Section IV. These phenomena would depend only on $`\omega `$r, the surface velocity of the planet relative to the Lorentz frame in which its axis of rotation is fixed. No variation in the speed of light would be found from the solar or galactic orbital velocities.
It is noteworthy that Michelson and Gale<sup>32</sup><sup>32</sup>32 A.A. Michelson, and H.G. Gale, โThe effect of the earthโs rotation on the velocity of light, Part II,โ Astrophys. J. 61, 140-145 (1925). See also A.A. Michelson, โThe effect of the earthโs rotation on the velocity of light, Part I,โ Astrophys. J. 61, 137-139 (1925). measured the Sagnac effect for the earthโs surface velocity in the 1920โs. And in order to maintain accuracy, the Global Positioning System must apply a Sagnac velocity correction to its electromagnetic signals<sup>33</sup><sup>33</sup>33 D. W. Allan, and M. A. Weiss, โAround-the-World Relativistic Sagnac Experiment,โ Science, 228, 69-70 (1985). .
### B. Modern Michelson-Morley Experiment
Although the original Michelson-Morley experiment and almost all subsequent tests of similar nature were not precise enough to detect any non-null effect due to the earth surface velocity, one such test was. In 1978, Brillet and Hall<sup>34</sup><sup>34</sup>34 A. Brillet and J. L. Hall, โImproved laser test of the isotropy of space,โ Phys. Rev. Lett., 42(9), 549-552 (1979). found a โnullโ effect at the $`\mathrm{\Delta }`$t/t = 3X10<sup>-15</sup> level, ostensibly verifying standard relativity theory to high order. However, to obtain this result they subtracted out a persistent โspuriousโ non-null signal of amplitude 2X10<sup>-13</sup> at twice the apparatus rotation frequency<sup>35</sup><sup>35</sup>35 Mark P. Haugan and Clifford M. Will, โModern tests of special relativity,โ Phys. Today, May 1987, 69-76. It is ironic that Haugan and Will cited the Brillet and Hall experiment as proof of the invariance of the speed of light. .
In 1981 Aspen<sup>36</sup><sup>36</sup>36 H. Aspen, โLaser interferometry experiments on light speed anisotropy,โ Phys. Lett., 85A(8,9), 411-414 (1981). pointed out that this โspuriousโ signal would correspond to a test apparatus velocity of 363 m/sec. The earth surface velocity due to its rotation at the test site is 355 m/sec.
## VI. Related Issues
In this section we briefly discuss several other rotating frame issues that are either prevalent in the literature or otherwise worth addressing.
### A. Relativistic Mass-Energy
Klauber<sup>37</sup><sup>37</sup>37 Ref. 29, pp. 427-429. uses transformation (2) to show that the mass-energy of an object fixed in the (NTO) rotating frame has typical relativistic dependence on speed, in this case $`\omega r`$. Given the well know cyclotron experiments, this lends further support for the correctness of transformation (2).
### B. Generalized Coordinates and Light Speed
Some may argue, in the spirit of generalized coordinates common to general relativity, that the velocity of light deduced in Section IV is merely a coordinate value. That is, it is not what one would measure with standard rods and clocks, but an expression in terms of arbitrary coordinates, which may in fact be anything one chooses. By choosing the appropriate generalized coordinate system we could then โdeduceโ any value we like.
This argument is in fact erroneous. One can calculate physical values for distance, time, velocity or any other measurable quantity from the metric of the particular coordinate grid employed. (See the Appendix for an example.) โPhysicalโ values are those one would measure using physical instruments such as standard rods and standard clocks. โCoordinateโ values are those one calculates using the arbitrary values for length and time associated with any arbitrarily chosen grid. In effect, physical values are those values one calculates when the associated basis vector of the generalized coordinate system has unit length. Given any coordinate value associated with a non-unit basis vector, one simply calculates the equivalent (physical) value associated with a unit basic vector pointed in the same direction. Coordinate values can be any number, while physical values are unique. For details we refer the reader to Misner, Thorne, and Wheeler<sup>38</sup><sup>38</sup>38 Charles W. Misner, Kip S. Thorne, and John A. Wheeler, Gravitation (Freeman, New York, 1973), pp. 37, 821-822, and many other places throughout the text. , Malvern<sup>39</sup><sup>39</sup>39 Lawrence E. Malvern, Introduction to the Mechanics of a Continuous Medium (Prentice-Hall, Englewood Cliffs, New Jersey, 1969), Appendix I, Sec. 5, pp. 606-613. , and Klauber<sup>40</sup><sup>40</sup>40 Ref. 29, pp. 426-427. .
Avoiding excessive complexity, we simply note that we have taken care that quantities in relationships such as (1) and (7) are indeed physical values. We note further that in the thought experiment of section II.A no coordinate grid whatsoever is used. The observer simply uses standard rods and clocks, which yield the same values regardless of the coordinate grid employed.
A related argument posits that the degree of orthogonality of any axis in generalized coordinates relative to any other axis is arbitrary, and hence we can choose our time axis in any direction we like. In fact, after using transformation (2), some<sup>41</sup><sup>41</sup>41 See for example, ref. 30, pp. 124. assert that we must then transform to a locally time orthogonal frame. But this is effectively the same as using local Lorentz co-moving frames. This in turn necessitates invariant, isotropic light speed and gives rise to problems already discussed, as well as the discontinuity in time described in the following subsection. While in generalized coordinates we can arbitrarily define our t coordinate in any of an infinite number of ways, if it is to represent the physical time nature chooses, then calculations done using it must match up with phenomena observed in the physical world.
We also note that any transformation to a rotating frame must incorporate the general form of (2.c). That is, the transformation of the azimuthal angle must include a term like $`\omega `$t, regardless of whether one believes other multiplicative factors (such as the second order Lorentz factor) should also be involved. When one squares the d$`\mathrm{\Phi }`$ of (2.c), or any other relation with a $`\omega `$dt term, and inserts the result in (3), one then ends up with off diagonal metric terms such as those of $`(4)`$. Hence no matter what physically reasonable transformation<sup>42</sup><sup>42</sup>42 For justification of (2.a), see ref. 29, pp. 416-418. one chooses, one must find that the rotating frame is NTO. And the primary effects of an NTO frame on observable phenomena are first order, i.e., they are independent of the presence or absence of a Lorentz factor in the transformation.
### C. Simultaneity in the Rotating Frame
Although this and the following subsection may seem counterintuitive to one cultured by a relativistic age, we suggest they deserve serious consideration as a possible way in which nature might actually work on rotating frames.
Using the co-moving local Lorentz frame methodology, one finds, due to the standard lack of agreement in simultaneity between Lorentz frames in relative motion, a quite bizarre result. (See Klauber<sup>43</sup><sup>43</sup>43 Ref. 29, pp. Pp. 413-415. .) If we consider a spatial path 360<sup>o</sup> around a given circumference, we find the clock at 360<sup>o</sup> has a different time on it than the clock at 0<sup>o</sup>, even though time remained constant all along the spatial path. This means the clock can not be synchronized with itself. It also implies that a continuous standard tape measure laid out around the circumference would not meet back up with itself at the same point in time. In other words, there would be a discontinuity in time.
Peres<sup>44</sup><sup>44</sup>44 Asher Peres, โSynchronization of clocks in a rotating frame,โ Phys. Rev. D, 18(6), pp. 2173-2174 (1978). noted this result as well, in addition to demonstrating that this methodology led to a โradial velocity of light \[that is\] not the same inward and outward.โ He concluded, โAll this is the heavy price which we are paying to make the azimuthal velocity of light โฆequal to c.โ
Reconciliation of analysis with reality occurs if simultaneity on a rotating disk is the same as that in the lab. Note that this is true for the time transformation of (2.a) wherein dt = 0 between two events, if dT = 0 between those same events.. Hence (1-$`\omega ^2`$r<sup>2</sup>/c<sup>2</sup>)<sup>1/2</sup>dt, the time passed on standard (physical) clocks in the rotating frame, is also zero. (Note that clocks running at different rates can still agree that no time passed on either one between events, and so can share a common simultaneity.) For this definition of simultaneity, there is no discontinuity in time. A line painted around a closed path on the rotating disk then does meet back up with itself at the same point in time.
### D. Length Contraction: To Be or Not to Be
The co-moving local Lorentz frames approach implies that standard rods on a rotating disk contract circumferentially. Many researchers<sup>45</sup><sup>45</sup>45 See ref. 3. have concluded from this that the disk surface is therefore curved, thereby ostensibly resolving the famous Ehrenfest paradox<sup>46</sup><sup>46</sup>46 P. Ehrenfest, โGleichf mrige Rotation starrer K rper und Relativit theorie,โ Phys. Z. 10, 918-918 (1909). <sup>,</sup><sup>47</sup><sup>47</sup>47 See for example, ref. 9. .
However, Tartagliaโs<sup>48</sup><sup>48</sup>48 A. Tartaglia, โLengths on rotating platformsโ, Found. Phys. Lett., 12(1), 17-28 (1999). interpretation of Ehrenfestโs paradox is more insightful. He notes that in Lorentz frames each observer sees rods in the otherโs frame as contracted, and โan observer on board the \[rotating\] disk would not perceive any curvature since in his reference frame \[there is no contraction\]โ.
Still further, if local Lorentz frames are valid surrogates for a rotating frame, then an observer on the rim of a rotating disk would likewise see the lab rods as contracted. And he would therefore conclude that the lab frame must be curved, which of course, it is not. Klauber<sup>49</sup><sup>49</sup>49 Ref. 29, pp 418-419, 431-433, 437, 439. and Tartaglia both conclude that internal contradictions in the theory disappear only if there is no Lorentz contraction effect between the disk and lab frames.
Transformation (2) actually implies this. Consider a small circumferentially aligned rod of proper length R$`\mathrm{\Delta }\mathrm{\Phi }`$ ($`\mathrm{\Delta }\mathrm{\Phi }`$ is small) in the lab, which to a lab observer must have simultaneous endpoints, both at time T= 0. Using (2.a), (2.b), and (2.c) one then finds the length of that rod as seen from the rotating disk to be r$`\mathrm{\Delta }\varphi `$ = R$`\mathrm{\Delta }\mathrm{\Phi }`$, i.e., no observed Lorentz contraction. The same logic works in reverse for a rod on the disk. Again, we emphasize that quantities discussed herein are physical, not merely coordinate, in nature.
The same conclusion may be drawn from the 4D line element ds, which is invariant between frames, NTO or not, inertial or not. In general between any two frames (notation should be obvious)
$$ds^2=cdT^2+dX^2+g_{XT}dXdT=cdt^2+dx^2+g_{xt}dxdt$$
(9)
where the off diagonal terms vanish for TO frames, and physical quantities are assumed (i.e., a coordinate grid with unit basis vectors is chosen.) Note that if the two frames share the same simultaneity, then when dT = 0, dt = 0, and therefore dx = dX. The length dx of a rod in one frame equals the length dX of the same rod as seen from the other frame, i.e., there is no Lorentz contraction. (Note that if dT = 0, but dt $``$ 0, then dx $``$ dX, and there is Lorentz contraction.<sup>50</sup><sup>50</sup>50 There is a little subtlety here. Calculation of the Lorentz contraction actually involves rod endpoints that appear to be different events to different observers. The present example assumes the same two events are measured by both observers. However, the conclusion remains valid. If the endpoint events look simultaneous to two different observers, then the rod length measured must be the same for each. )
The Lorentz contraction results from the Lorentz transformation, and in a sense is little more than an optical illusion. No Lorentz contracted object ever โfeelsโ contracted<sup>51</sup><sup>51</sup>51 See ref. 29, pp. 415, 418-419, 423. . The contraction appears because of the disagreement in simultaneity between frames, which is inherent within the Lorentz transformation.
We conclude that if two frames share the same simultaneity, then there is no Lorentz contraction effect between those frames. Hence, equivalence of simultaneity for the lab and rotating frames means no Ehrenfest paradox, as well as no discontinuity in time<sup>52</sup><sup>52</sup>52 Neither Lorentz contraction nor simultaneity differences can presently be measured directly by experiment. .
### E. Absolute Nature of Rotation
Translating systems display robust relativity. For such systems, absolute velocity does not exist, and there is no way to determine a preferred frame. As a result, light speed is invariant and transformations must be Lorentzian.
Rotating systems differ from translating systems in that one can determine oneโs angular velocity absolutely (in the sense of Mach). The preferred frame is then the non-rotating frame, which any observer can readily identify. Further, an observer within her own frame can, in fact, do tests that unambiguously determine her circumferential velocity relative to the Lorentz frame in which her axis of rotation is fixed.
If two types of frames have such fundamental dissimilarity at their cores, is it not presumptuous to assume, as has been the general practice, that the phenomena associated with each are identical? Is it correct to simply presume that invariant light speed, Lorentz contraction, and disagreement in simultaneity can be directly extrapolated to rotating frames? We suggest that it is not, and that the proper course of action consists of building a theory of relativistic rotation from empirical data on rotating frames alone, independently of preconceived conceptions.
## VII. Summary and Conclusions
The analysis of the previous section may, of course, run counter to a few long cherished ideas. But it does cleanly resolve certain major issues, and in the process leaves the essence of relativity theory intact.
Invariants like ds remain invariant. Every aspect of the theory for time orthogonal frames (the vast majority of applications) remains unchanged. Lorentz frames are still related by Lorentz transformations (with concomitant effects such as Lorentz contraction, etc.), and differential geometry continues its reign as descriptor of non-inertial systems. Neither special nor general relativity need be altered in any regard, provided of course, that NTO frames are appropriately interpreted.
All of the results obtained herein are derived from two postulates: i) transformation (2) relates rotating and non-rotating frames, and ii) the 4D line element ds is invariant. If these postulates are valid, it appears one must conclude the following.
NTO frames display non-invariant and non-isotropic local, physical speed of light, to a degree dependent on the degree of non-time-orthogonality. Lorentz frames are appropriate local surrogates for (curved or flat) TO frames, but not for (curved or flat) NTO frames. Rotating frames are truly (not superficially) NTO, and if one makes a straightforward interpretation, unfettered by preconceptions, of the most widely accepted transformation, one concludes that rotating and non-rotating frames share the same simultaneity. This in turn implies an absence of the Lorentz contraction effect between those frames.
NTO analysis for rotating frames predicts time dilation and mass energy increase with tangential speed in consonance with cyclotron experiments. It is also in full agreement with the Sagnac experiment and related thought experiments. Importantly, it also explains the persistent signal found in the Michelson-Morley type Brillet and Hall experiment, which has heretofore been considered inexplicable.
## Appendix A Appendix
The following discussion should be read only in the context of time orthogonal reference frames. Such frames make up the bulk of all applications, and virtually 100% of textbook problems. The conclusions drawn in this appendix are subsequently modified for NTO frames.
The speed of light in non-Lorentzian systems can be a source of confusion as it is sometimes (misleadingly) said that the speed of light in general relativity can be different than c. This is true if, for example, one measures the speed of light near a massive star using a clock based on earth. (Time on such a clock is effectively the coordinate time in a Schwarzchild coordinate system.) As is well known, due to the intense gravitation field, the passage of time close to the star is dilated relative to earth time, and using the earth clock, one would indeed calculate a light speed other than c. However, use of standard rods and clocks adjacent the light ray itself would result in a speed of precisely c. In the language of Section VI.B, the speed of light calculated with the earth clock is a coordinate speed, whereas that measured with rods and clocks proximate to the light ray is the physical speed.
Other confusion exists for scenarios where spacetime itself expands or contracts. For example, just after the big bang, space itself was expanding much like the surface of a balloon being blown up. A photon in space (analogous to an ant on the surface of the balloon) at a different location than an observer could then move away from the observer faster than c (analogous to faster than the ant can crawl on the surface) because the space (balloon surface) between the photon and the observer is itself expanding. Yet a photon spatially coincident with an observer could never be seen by that observer to have speed greater than c, and local standard rods and clocks adjacent any photon would find its speed equaling c regardless of the dynamical state of spacetime itself.
More mathematically, for a given generalized coordinate system in a non-inertial frame, we have
$$ds^2=g_{tt}c^2dt^2+g_{xx}dx^2+g_{yy}dy^2+g_{zz}dz^2$$
where g<sub>tt</sub> is negative. Consider a ray of light passing in the x direction such that dy = dz = 0. The coordinate speed of light is found by setting ds = 0 and solving for the length in generalized coordinates (the coordinate length) that the light ray travels divided by the time in generalized coordinates (the coordinate time), i.e.,
$$\frac{dx}{dt}=\sqrt{\frac{g_{tt}}{g_{xx}}}c.$$
On the other hand, the physical speed of light is the physical length divided by the physical time. Physical length (measured by standard rods) is $`\sqrt{g_{xx}}dx`$<sup>53</sup><sup>53</sup>53 See refs. 38,39, and 40. , and physical time (measured by standard clocks) is $`\sqrt{g_{tt}}dt`$. So the speed of light as measured by physical instruments, regardless of the generalized coordinates chosen, is
$$\frac{\sqrt{g_{xx}}dx}{\sqrt{g_{tt}}dt}=c,$$
and this is always equal to c (for time orthogonal frames.)
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# Brane Gases in the Early Universe
## I Introduction
In this paper we consider an approach to string cosmology in close analogy to the usual starting point of standard big-bang cosmology. We assume that the Universe started out small, dense, hot, and with all fundamental degrees of freedom in thermal equilibrium. We also assume that the background space is toroidal in all spatial dimensions <sup>*</sup><sup>*</sup>*In fact, we do not need to be this specific. The crucial assumption is the existence of one-cycles in all spatial directions.. Given these assumptions, the initial state will consist of a gas of all fundamental branes which the theory admits. We will study the equation of state of the individual types of branes, neglecting for simplicity brane interactions, and will use the results to determine the source terms in the equations of motion for the background. In particular, we will study obstructions to spatial dimensions becoming large. We will find that fundamental string winding modes dominate the evolution at late times, and prevent more than three spatial dimensions from becoming large, in agreement with the scenario proposed in in the context of perturbative string theory. We also argue that because of t-duality the usual singularities of a the homogeneous and isotropic big bang and inflationary cosmologies are not present.
The main goal of this paper is to generalize the considerations of to the context of our present understanding of string theory. In the 1980โs, it was believed that the only fundamental degrees of freedom of string theory were the fundamental strings. In this context, it was shown that t-duality and string winding modes could have a very important effect on early Universe cosmology. Assuming that the background space is toroidal and thus admits string winding modes, it was shown that t-duality could explain the absence of the initial big-bang singularity. In addition, it was speculated that string winding modes would only allow three spatial dimensions to become large. The second point was put on a firmer basis by the work of Tseytlin and Vafa , who discussed the effects of gases of strings on the background equations of motion which were taken to be those of dilaton gravity.
The first main point of the BV scenario is that t-duality will lead to an equivalence of the physics if the radius of the background torus changes (in string units) from $`R`$ to $`1/R`$. This corresponds to an interchange of momentum and winding modes. Thus, $`R`$ becoming small is equivalent to $`R`$ tending to infinity. Neither limit corresponds to a singularity for string matter. For example, the temperature $`T`$ obeys
$$T(\frac{1}{R})=T(R).$$
(1)
Thus, in string cosmology the big bang singularity can be avoided. The second point suggested in was that string winding modes would prevent more than three spatial dimensions from becoming large. The point is that string winding modes cannot annihilate in more than three spatial dimensions (by a simple classical dimension counting argument). In the context of dilaton cosmology, a gas of string winding modes (which has an equation of state $`\stackrel{~}{p}=(1/d)\rho `$, where $`\stackrel{~}{p}`$ and $`\rho `$ denote pressure and energy density, respectively, and $`d`$ is the number of spatial dimensions) will lead to a confining potential in the equation of motion for $`\lambda =log(a)`$, where $`a(t)`$ is the scale factor of the Universe . Note that this is not the result which would be obtained in a pure metric background obeying the Einstein equations. The dynamics of classical strings in higher dimensional expanding backgrounds was studied numerically in , confirming the conclusions of .
However, it is now clear that string theory has a much richer set of fundamental degrees of freedom, consisting - in addition to fundamental strings - of D-branes of various dimensionalities. The five previously known consistent perturbative string theories are now known to be connected by a web of dualities , and are believed to represent different corners of moduli space of a yet unknown theory called M-theory. Which branes arise in the effective string theory description depends on the particular point in moduli space. We will be making a specific assumption below.
The question we would like to address is whether the inclusion of the new fundamental degrees of freedom will change the main cosmological implications of string theory suggested in , namely the avoidance of the initial cosmological singularity, and the singling out of 3 as the maximal number of large spatial dimensions, in the context of an initial state which is assumed to be hot, dense and small, and of a background geometry which admits string winding modes.
Our concrete starting point is 11-dimensional M-theory compactified on $`S^1`$ to yield 10-dimensional Type II-A string theory. The resulting low energy effective theory is supersymmetrized dilaton gravity. As fundamental states, M-theory admits the graviton, 2-branes and 5-branes. After compactification, this leads to 0-branes, 1-branes, 2-branes, 4-branes, 5-branes, 6-branes and 8-branes as the fundamental extended objects of the 10-dimensional theory. The dilaton represents the radius of the compactified $`S^1`$. We are in a region of moduli space in which the string coupling constant $`g_s`$ is smaller than 1.
We assume that all spatial dimensions are toroidal (radius $`R`$), and that the Universe starts out small, dense, hot, and in thermal equilibrium. Thus, the Universe will contain a gas of all branes appearing in the spectrum of the theory. Note that this starting point is in close analogy with the hot big bang picture in standard cosmology, but very different from brane-world scenarios in which the existence of a particular set of branes is postulated from the outset without much justification from the point of view of cosmology.
There have been several interesting previous studies of the cosmology of brane gases. Maggiore and Riotto (see also ) studied the phase diagram of brane gases motivated by M-theory as a function of the string coupling constant and of the Hubble expansion rate (as a measure of space-time curvature) and discovered regions of the phase diagram in which brane gases determine the dynamics, and regions in which the effective action is no longer well described by a ten-dimensional supergravity action. Given our assumptions, we are in a region in moduli space in which the ten-dimensional effective description of the physics remains true to curvature scales larger than that given by the string scale. In this paper, we consider the time evolution of the system through phase space starting from some well-defined initial conditions. We will argue that as a consequence of t-duality, curvature scales where the ten-dimensional description breaks down are never reached.
In another interesting paper, Park et al. take a starting point very close to our own, a hot dense gas of branes. However, they did not consider the winding and oscillatory modes of the branes.
In the following section we will study the equation of state of the brane gases for all values of their spatial dimension p. We will separately analyze the contributions of winding and non-winding modes (the latter treated perturbatively). The results will be used as source terms for the equations of motion of the background dilaton gravity fields, following the approach of . We find that the winding modes of any p-brane lead to a confining force which prevents the expansion of the spatial dimensions, and that the branes with the largest value of p give the largest contribution to the energy of the gas in the phase in which the scale factor is increasing.
In Section 3 we use the results of the previous section to argue that the main conclusions of the scenario proposed in are unchanged: t-duality eliminates the cosmological singularity, and winding modes only allow three dimensions of space to become large. We point out a potential problem (the brane problem) of cosmologies based on theories which admit branes in their spectrum of fundamental states. This problem is similar to the well-known domain wall problem in cosmological models based on quantum field theory. It is pointed out that a phase of loitering (see e.g. ) yields a natural solution of this problem, and it is shown that the background equations of motion may well yield a loitering stage during the early evolution of the Universe. Some limitations of our considerations and avenues for future research are discussed in the final section.
## II Equation of State of Brane Gases
As mentioned in the Introduction, our starting point is Type II-A string theory on a 9-dimensional toroidal background space (with the time direction being infinite), resulting from the compactification of M-theory on $`S^1`$. The fundamental degrees of freedom are the fields of the bulk background (resulting from the graviton in M-theory), 1-branes, 2-branes, 4-branes, 5-branes, 6-branes and 8-branes.
The low-energy bulk effective action is given by
$`S_{bulk}={\displaystyle \frac{1}{2\kappa ^2}}{\displaystyle }d^{10}x\sqrt{G}e^{2\varphi }[R`$ $`+`$ $`4G^{\mu \nu }_\mu \varphi _\nu \varphi `$ (2)
$``$ $`{\displaystyle \frac{1}{12}}H_{\mu \nu \alpha }H^{\mu \nu \alpha }],`$ (3)
where $`G`$ is the determinant of the background metric $`G_{\mu \nu }`$, $`\varphi `$ is the dilaton, $`H`$ denotes the field strength corresponding to the bulk antisymmetric tensor field $`B_{\mu \nu }`$, and $`\kappa `$ is determined by the 10-dimensional Newton constant in the usual way.
The total action is the sum of the above bulk action and the action of all branes present. The action of an individual brane with spatial dimension p has the Dirac-Born-Infeld form
$$S_p=T_pd^{p+1}\zeta e^\varphi \sqrt{det(g_{mn}+b_{mn}+2\pi \alpha ^{}F_{mn})}$$
(4)
where $`T_p`$ is the tension of the brane, $`g_{mn}`$ is the induced metric on the brane, $`b_{mn}`$ is the induced antisymmetric tensor field, and $`F_{mn}`$ the field strength tensor of gauge fields $`A_m`$ living on the brane. The total action is the sum of the bulk action (2) and the sum of all of the brane actions (4), each coupled as a delta function source (a delta function in the directions transverse to the brane) to the 10-dimensional action.
The induced metric on the brane $`g_{mn}`$, with indices $`m,n,\mathrm{}`$ denoting space-time dimensions parallel to the brane, is determined by the background metric $`G_{\mu \nu }`$ and by scalar fields $`\varphi _i`$ (not to be confused with the dilaton $`\varphi `$) living on the brane (with indices $`i,j,\mathrm{}`$ denoting dimensions transverse to the brane) which describe the fluctuations of the brane in the transverse directions:
$$g_{mn}=G_{mn}+G_{ij}_m\varphi _i_n\varphi _j+G_{in}_m\varphi _i.$$
(5)
The induced antisymmetric tensor field is
$$b_{mn}=B_{mn}+B_{ij}_m\varphi _i_n\varphi _j+B_{i[n}_{m]}\varphi _i.$$
(6)
In addition,
$$F_{mn}=_{[m}A_{n]}.$$
(7)
In the string frame, the fundamental string has tension
$$T_f=\mathrm{\hspace{0.17em}2}\pi \alpha ^{},$$
(8)
whereas the brane tensions for various values of p are given by
$$T_p=\frac{\pi }{g_s}(4\pi ^2\alpha ^{})^{(p+1)/2},$$
(9)
where $`\alpha ^{}l_{st}^2`$ is given by the string length scale $`l_{st}`$ and $`g_s`$ is the string coupling constant Note that the s-duality between the fundamental and D-string tensions is more evident in the Einstein frame.. Note that all of these branes have positive tension.
In the following, we wish to compute the equation of state of the brane gases for a general value of p. For our considerations, the most important modes are the winding modes. If the background space is $`T^9`$, a p-brane can wrap around any set of p toroidal directions. The modes corresponding to these winding modes by t-duality are the momentum modes corresponding to center of mass motion of the brane. The next most important modes for our considerations are the modes corresponding to fluctuations of the brane in transverse directions. These modes are in the low-energy limit described by the brane scalar fields $`\varphi _i`$. In addition, there are bulk matter fields and brane matter fields.
Since we are mainly interested in the effects of a gas of brane winding modes and transverse fluctuations on the evolution of a spatially homogeneous Universe, we will neglect the antisymmetric tensor field $`B_{\mu \nu }`$. We will use conformal time $`\eta `$ and take the background metric to be given by
$$G_{\mu \nu }=a(\eta )^2diag(1,1,\mathrm{},1),$$
(10)
where $`a(\eta )`$ is the cosmological scale factor.
If the transverse fluctuations of the brane are small (in the sense that the first term on the right hand side of (5) dominates) and the gauge fields on the brane are small, then the brane action can be expanded as follows:
$`S_{brane}`$ $`=`$ $`T_p{\displaystyle d^{p+1}\zeta a(\eta )^{p+1}e^\varphi }`$ (12)
$`e^{\frac{1}{2}trlog(1+_m\varphi _i_n\varphi _i+a(\eta )^22\pi \alpha ^{}F_{mn})}`$
$`=`$ $`T_p{\displaystyle d^{p+1}\zeta a(\eta )^{p+1}e^\varphi }`$ (14)
$`(1+{\displaystyle \frac{1}{2}}(_m\varphi _i)^2\pi ^2\alpha _{}^{}{}_{}{}^{2}a^4F_{mn}F^{mn}).`$
The first term in the parentheses in the last line corresponds to the brane winding modes, the second term to the transverse fluctuations, and the third term to brane matter. We see that, in the low energy limit, the transverse fluctuations of the brane are described by a free scalar field action, and the longitudinal fluctuations are given by a Yang-Mills theory. The induced equation of state has pressure $`p0`$.
The above result extends to the case of large brane field and brane position fluctuations. It can be shown that large gauge field fluctuations on the brane give rise to the same equation of state as momentum modes ($`E1/R`$) and are thus also described by pressure $`p0`$. In the high energy limit of closely packed branes, the system of transverse brane fluctuations is described by a strongly interacting scalar field theory which also corresponds to pressure $`p0`$.
We will now consider a gas of branes and determine the equations of state corresponding to the various modes. The procedure involves taking averages of the contributions of all of the branes to the energy-momentum tensor, analogous to what is usually done in homogeneous cosmology generated by a gas of particles.
Let us first focus on the winding modes. From (12) it immediately follows that the winding modes of a p-brane give rise to the following equation of state:
$$\stackrel{~}{p}=w_p\rho \mathrm{with}w_p=\frac{p}{d}$$
(15)
where $`d`$ is the number of spatial dimensions (9 in our case), and where $`\stackrel{~}{p}`$ and $`\rho `$ stand for the pressure and energy density, respectively.
Since both the fluctuations of the branes and brane matter are given by free scalar fields and gauge fields living on the brane (which can be viewed as particles in the transverse directions extended in brane directions), the corresponding equation of state is that of โordinaryโ matter with
$$\stackrel{~}{p}=w\rho \mathrm{with}\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}w1.$$
(16)
Thus, in the absence of a scalar field sector living on the brane, the energy will not increase as the spatial dimensions expand, in contrast to the energy in the winding modes which evolves according to (as can again be seen immediately from (12))
$$E_p(a)T_pa(\eta )^p,$$
(17)
where the proportionality constant depends on the number of branes. Note that the winding modes of a fundamental string have the same equation of state as that of the winding modes of a 1-brane, and the oscillatory and momentum modes of the string obey the equation of state (16).
In the context of a hot, dense initial state, the assumption that the brane fluctuations are small will eventually break down. Higher order terms in the expansion of the brane action will become important. One interesting effect of these terms is that they will lead to a decrease in the tension of the branes . This will occur when the typical energy scale of the system approaches the string scale. At that point, the state of the system will be dominated by a gas of branes.
The background equations of motion are
$`d\dot{\lambda }^2+\dot{\phi }^2`$ $`=`$ $`e^\phi E`$ (18)
$`\ddot{\lambda }\dot{\phi }\dot{\lambda }`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^\phi P`$ (19)
$`\ddot{\phi }d\dot{\lambda }^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^\phi E,`$ (20)
where $`E`$ and $`P`$ denote the total energy and pressure, respectively,
$$\lambda (t)=log(a(t)),$$
(21)
and $`\phi `$ is a shifted dilaton field which absorbs the space volume factor
$$\phi =\mathrm{\hspace{0.17em}2}\varphi d\lambda .$$
(22)
In our context, the matter sources $`E`$ and $`P`$ obtain contributions from all components of the brane gas:
$`E`$ $`=`$ $`{\displaystyle \underset{p}{}}E_p^w+E^{nw}`$ (23)
$`P`$ $`=`$ $`{\displaystyle \underset{p}{}}w_pE_p^w+wE^{nw},`$ (24)
where the superscripts $`w`$ and $`nw`$ stand for the winding modes and the non-winding modes, respectively. The contributions of the non-winding modes of all branes have been combined into one term. The constants $`w_p`$ and $`w`$ are given by (15) and (16). Each $`E_p^w`$ is the sum of the energies of all of the brane windings with fixed p.
## III Brane Gases in the Early Universe
The first important conclusion of was that in the approach to string cosmology based on considering string gases in the early Universe, the initial cosmological (Big Bang) singularity can be avoided. The question we will now address is whether this conclusion remains true in the presence of branes with $`p>1`$ in the spectrum of fundamental states.
The two crucial facts leading to the conclusions of were t-duality and the fact that in the micro-canonical ensemble the winding modes lead to positive specific heat, leading to limiting Hagedorn temperature The Hagedorn temperature is not reached at finite energy density .. Both of these facts extend to systems with branes. First, as is obvious, each brane sector by itself preserves t-duality. Secondly, it was shown in that if two or more Hagedorn systems thermally interact, and at least one of them (let us say System 1) has limiting Hagedorn temperature, then at temperatures close to the Hagedorn temperature of System 1, most energy flows into that system, and the joint system therefore also has limiting Hagedorn behavior. Hence, as the Universe contracts, the t-duality fixed point $`R=1`$ is reached at a temperature $`T`$ smaller than the Hagedorn temperature, and as the background space contracts further, the temperature starts to decrease according to (1). There is no physical singularity as $`R`$ approaches 0.
Let us now turn to the dynamical de-compactification mechanism of 3 spatial dimensions suggested in . We assume that the Universe starts out hot, small and in thermal equilibrium, with all spatial dimensions equal (near the self-dual point $`R=1`$). In this case, in addition to the momentum and oscillatory modes, winding modes of all p-branes will be excited. By symmetry, it is reasonable to assume that all of the net winding numbers cancel, i.e. that there are an equal number of winding and anti-winding modes.
Let us assume that the Universe starts expanding symmetrically in all directions. As $`\lambda `$ increases, the total energy in the winding modes increases according to (17), the contribution of the modes from the branes with the largest value of $`p`$ growing fastest. In exact thermal equilibrium energy would flow from the winding modes into non-winding modes. However, this only can occur if the rate of interactions of the winding modes is larger than the Hubble expansion rate.
Generalizing the argument of , from a classical brane point of view it follows (by considering the probability that the world-volumes of 2 p-branes in space-time intersect) that the winding modes of p-branes can interact in at most $`2p+1`$ large <sup>ยง</sup><sup>ยง</sup>ยงLarge compared to the string scale. spatial dimensions. Thus, in $`d=9`$ spatial dimensions, there are no obstacles to the disappearance of $`p=8`$, $`p=6`$, $`p=5`$ and $`p=4`$ winding modes, whereas the lower dimension brane winding modes will allow a hierarchy of dimensions to become large. Since for volumes large compared to the string volume the energy of the branes with the largest value of $`p`$ is greatest, the 2-branes will have an important effect first. They will only allow 5 spatial dimensions to become large. Within this distinguished $`T^5`$, the 1-brane winding modes will only allow a $`T^3`$ subspace to become large. Thus, it appears that the mechanism proposed in will also apply if the Hilbert space of states includes fundamental branes with $`p>1`$.
To what extent can these classical arguments be trusted? It has been pointed out in that the microscopic width of a string increases logarithmically as the energy with which one probes the string increases. However, in our cosmological context, we are restricted to energy densities lower than the typical string density, and thus the effective width of the strings is of string scale . Similar conclusions will presumably apply to branes of higher dimensionality. However, no definite results are known since a rigorous quantization scheme for higher dimensional branes is lacking.
The cosmological scenario we have in mind now looks as follows: The Universe starts out near the self-dual point as a hot, dense gas of branes, strings and particles. The Universe begins to expand in all spatial directions as described by the background equations of motion (18 \- 20). As space expands and cools (and the brane tension therefore increases), the branes will eventually fall out of thermal equilibrium. The branes with the largest value of p will do this first. Space can only expand further if the winding modes can annihilate. This can be seen immediately from the background equations of motion (19). If the equation of state is dominated by winding modes (which it would be if the Universe would keep on expanding), then (with the help of (15) and (17)) it follows that the right hand side of that equation acts as if it comes from a confining potential
$$V_{eff}(\lambda )=\beta _pe^\phi e^{p\lambda }$$
(25)
where $`\beta _p`$ is a positive constant which depends on the brane tension.
The unwinding of p-branes poses no problem for $`p=8`$, $`p=6`$, $`p=5`$ and $`p=4`$ For a discussion of the microphysics of brane winding mode annihilation see e.g. and references therein.. The corresponding brane winding modes will disappear first. However, the $`p=2`$ branes will then only allow 5 spatial dimensions to expand further (which 5 is determined by thermal fluctuations). In this distinguished $`T^5`$, the one-branes and fundamental strings will then only allow a $`T^3`$ subspace to expand. Thus, there will be a hierarchy of sizes of compact dimensions. In particular, there will be 2 extra spatial dimensions which are larger than the remaining ones. The connection with the large-extra dimension scenario proposed in is intriguing. However, since the only fundamental scale in the problem is the string scale, and since there does not seem to be a way to naturally separate the temperature scales at which the various branes fall out of equilibrium, it appears difficult to produce the mm scale proposed in .
Note that even when winding mode annihilation is possible by dimension counting, causality imposes an obstruction. There will be at least 1 winding mode per Hubble volume remaining (see e.g. ). In our four-dimensional space-time, all branes with $`p2`$ will look like domain walls. This leads to the well known domain wall problem for cosmology since one wall per Hubble volume today will overclose the Universe if the tension of the brane is larger than the electroweak scale. Due to their large tension at low temperatures, even one-branes will overclose the Universe.
There are two ways to overcome this domain wall problem. The first is to invoke cosmic inflation at some stage after the branes have fallen out of equilibrium but before they come to dominate the energy of the Universe. The scenario would be as follows: initially, the winding modes dominate the energy density and determine the dynamics of space-time. Once they fall out of equilibrium, most will annihilate and the energy in the brane winding modes will become subdominant. The spatial dimensions in which unwinding occurs will expand. It is at this stage that the ordinary field degrees of freedom in the theory must lead to inflation, before the remnant winding modes (one per Hubble volume) become dominant.
In our context, however, there is another and possibly more appealing alternative - loitering . If at some stage in the Universe the Hubble radius More precisely, the โcausal horizonโ, meaning the distance which light can travel during the time when $`H^1`$ is larger than the spatial extent. becomes larger than the spatial extent of the Universe, there is no causal obstruction for all winding modes to annihilate and disappear. In the context of โstandardโ cosmology (fields as matter sources coupled to classical general relativity) it is very hard to obtain loitering. However, in the context of dilaton cosmology a brief phase of loitering appears rather naturally as a consequence of the confining potential (25) due to the winding modes. To see this, let us return to the background equations (18 \- 20). The phase space of solutions was discussed for general values of $`p`$ in .
For a fixed value of $`p`$, the phase space of solutions is two-dimensional and is spanned by $`l=\dot{\lambda }`$ and $`f=\dot{\phi }`$. If we start in the energetically allowed part
$$|l|<\frac{1}{\sqrt{d}}|f|$$
(26)
of the upper left quadrant of phase space with $`l>0`$ and $`f<0`$ corresponding to expanding solutions with $`\dot{\varphi }<0`$, then the solutions are driven towards $`l=0`$ at a finite value of $`f`$ (see Figure 1). More precisely, there are three special lines in phase space with
$$\frac{\dot{l}}{\dot{f}}=\frac{l}{f},$$
(27)
which correspond to straight line trajectories through the origin. They are
$$\frac{l}{f}=\pm \frac{1}{\sqrt{d}},\frac{p}{d}.$$
(28)
Solutions in the energetically allowed part of the upper left quadrant are repelled by the special line $`l/f=1/\sqrt{d}`$ and approach $`l=0`$. For $`l0`$, both $`\dot{l}`$ and $`\dot{f}`$ remain finite
$$\dot{l}\frac{p}{2d}f^2,\dot{f}\frac{f^2}{2}.$$
(29)
Hence, the trajectories cross the $`l=0`$ axis at some time $`t_1`$. This means that the expansion of space stops and the Universe begins to contract. The crossing time $`t_1`$is the first candidate for a loitering point.
Note that the dynamics of the initial collapsing phase is very different from the time reverse of the initial expanding phase. In fact, as the special line $`l/f=p/d`$ is approached, $`\dot{l}`$ changes sign again, and the trajectory approaches the static solution $`(l,f)=(0,0)`$ \- also implying a fixed value for the dilaton. At late times, the time evolution along the abovementioned special line corresponds to contraction with a Hubble rate whose absolute value is decreasing,
$$H(t)=\frac{1}{|H^1(t_0)|+\beta (tt_0)},$$
(30)
(where $`t_0`$ is some starting time along the special line) with
$$\beta =\frac{p}{2}+\frac{d}{2p}.$$
(31)
Thus, the evolution slows down and loitering is reached, even if the time evolution at $`t_1`$ is too rapid for loitering to occur then. Note that these considerations assume that the winding modes are not decaying. If they decay too quickly, this could obviously prevent loitering.
Note that given our cosmological starting point there is no horizon problem since space was initially of string size. However, two of the other problems of standard cosmology which the inflationary scenario successfully addresses, namely the flatness and the formation of structure problem, persist in our scenario. Note, in particular, that it seems necessary to have something like inflation of the large spatial dimensions in order to produce a Universe which is larger than the known Hubble radius. It is of great interest to explore the possibility of finding a solution of these problems in the context of string theory, e.g. along the lines of the string-driven inflationary models of .
## IV Discussion and Conclusions
In this paper we have generalized the approach to superstring cosmology pioneered in to include the contribution of branes. Our starting point is the assumption that the Universe starts out small, hot and in thermal equilibrium, with all of the spatial dimensions being equivalent and compact (string scale). We also assume that the background admits one cycles. We work in the corner of moduli space resulting from compactification of 11-d M-theory on $`S^1`$, and in which the string coupling constant is small.
We argue that, as a consequence of t-duality, the usual big bang singularity is absent in the resulting cosmology. Furthermore, since winding modes of all of the branes are excited in thermal abundance in the initial state, and since the energy in winding modes increases as the background spatial scale expands, the thermodynamics of the winding modes coupled to the background equations of motion, which in the corner of moduli space of M-theory which we consider are the dilaton-gravity equations, dominates the initial evolution of the background. Our thermodynamic considerations suggest that the mechanism first pointed out in , by which the string winding modes will prevent all but 3 spatial dimensions from becoming large, persists. In addition, the presence of winding modes of 2-branes may lead to a hierarchy in the sizes of the extra dimensions, with exactly internal dimensions being larger than the others.
We pointed out some further interesting characteristics of the resulting cosmology. In particular, the horizon problem is absent. However, in order to produce a Universe in which the large spatial dimensions exceed the present Hubble radius, it seems necessary to have a background evolution of the large spatial dimensions resembling inflation.
We also pointed out the existence of a brane problem, a problem for cosmology in theories with stable branes which is analogous to the domain wall problem in cosmological scenarios based on quantum field theories which admit stable domain walls. A phase of loitering in the background cosmological evolution will naturally solve this problem. Based on the background equations of motion, it appears that as long as the winding modes do not disappear, the background solutions approach a point of loitering.
It would be interesting to extend our considerations to other regions in moduli space, in particular to regions of strong string coupling. In those regions it appears that the effective ten-dimensional background description breaks down before the Hagedorn temperature is reached.
It is also important to point out that on Calabi-Yau threefold backgrounds one cycles are absent, and thus our cosmological scenario does not apply. Calabi-Yau three-folds are required if the four-dimensional low energy effective theory is to have $`N=1`$ supersymmetry. In the context of early Universe cosmology, however, it is not reasonable to demand $`N=1`$ supersymmetry. In particular, one could have maximal supersymmetry which is consistent with the toroidal background we are using. It would be interesting to explore whether consistent four-dimensional low-energy effective theories can be constructed from compact backgrounds which admit one-cycles.
Given that the tension of the branes exceeds that set by the string scale, the applicability of the homogeneous background field equations to brane gases might be questionable, in particular at late times when the branes are fairly widely separated <sup>\**</sup><sup>\**</sup>\**We thank D. Lowe for stressing this concern to us.. This issue deserves further study.
Acknowledgements
The research was supported in part by the U.S. Department of Energy under Contract DE-FG02-91ER40688, TASK A. Two of us (S.A. and D.E.) are supported by fellowships by the U.S. Department of Education under the GAANN program. We are grateful to E. Akhmedov, R. Easther, B. Greene, A. Jevicki, D. Lowe and S. Ramgoolam for many discussions about this project.
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# Drift without flux: Brownian walker with a space dependent diffusion coefficient
Brownian motion of spherical colloidal particles in the vicinity of a wall has been extensively studied, both theoretically and experimentally . It has been shown that the diffusion coefficients parallel or perpendicular to the wall were greatly reduced when the particles were close enough to the obstacle, i.e. within distances comparable to or less than their radius. When the Brownian particles are trapped in a more confined geometry, such as a porous medium, the theory is far more complicated and few experimental studies have been reported in model geometries, where the particles are trapped between two parallel walls. In this article, we report some new experimental results concerning the Brownian motion of particles trapped between two nearly parallel walls, so that the confinement, and thus the diffusion coefficient, become space dependent. As a result, we not only measure a diffusion coefficient which varies with the confinement, but also a drift of the particulesโ individual positions in the direction of the diffusion coefficient gradient, in the absence of any external force or concentration gradient. This drift was not accompanied by any net particle flux, i.e. statistically the same number of particles crossed any imaginary surface in both directions. We first discuss the general problem of a Brownian walker with a spatially dependent diffusion coefficient to explain the origin of the expected drift, and then present the experimental set-up and results.
As in our experiment the diffusion coefficient varies in only one direction, say $`x`$, we briefly sketch a heuristic derivation of the 1D Brownian walker algorithm. The velocity of a 1D Brownian particle subjected to a random force and a viscous drag follows the Langevin equation,
$$\frac{dv(t)}{dt}=\gamma v(t)+\mathrm{\Gamma }(t),$$
(1)
where $`\gamma ^1`$ is the velocity relaxation time and $`\mathrm{\Gamma }(t)`$ the random force per unit mass defined by its mean value $`\mathrm{\Gamma }(t)=0`$ and correlation function $`\mathrm{\Gamma }(t)\mathrm{\Gamma }(t^{})=q\delta (tt^{})`$. Using the equipartition theorem it can be shown that $`q`$ is related to the temperature $`T`$ and the particleโs mass, $`m`$, by the standard relation $`q=2\gamma kT/m`$. Discretizing the random function $`\mathrm{\Gamma }(t)`$ over time intervals $`\mathrm{\Delta }t>>\gamma ^1`$ allows us to drop in Eq.(1) the inertial term, $`dv/dt`$, and to replace the velocity $`v`$ by $`\mathrm{\Delta }x/\mathrm{\Delta }t`$. Choosing for $`\mathrm{\Gamma }(t)`$ the simplest random function, $`\mathrm{\Gamma }(t)=\pm \sqrt{q/\mathrm{\Delta }t}`$, leads to the well known Brownian walker algorithm,
$$x(t+\mathrm{\Delta }t)=x(t)\pm \sqrt{2D\mathrm{\Delta }t}$$
(2)
with $`D=kT/m\gamma `$. When the diffusion coefficient $`D`$, i.e. when the temperature $`T`$ and/or the drag coefficient $`\gamma `$ become position dependent, the above algorithm needs to be clarified. During each time interval $`\mathrm{\Delta }t`$, the walker makes a step to the right or to the left, but should the length of this position dependent step, $`\sqrt{2D\mathrm{\Delta }t}`$, be computed at the departure point $`x(t)=x`$, the arrival point $`x(t+\mathrm{\Delta }t)=x+\mathrm{\Delta }x`$ or at any point in between? These mathematical choices, often referred to as the Ito/Stratonovitch conventions, model different physical situations and the choice of convention is dictated only by the physics. We denote by $`D(x+\alpha \mathrm{\Delta }x)`$ the diffusion coefficient appearing in Eq.(2) where $`\alpha =0,1/2`$ and $`1`$ correspond to the Ito, Stratonovitch and isothermal choices respectively. As we will show, this last case models a situation where the temperature, $`T`$, is uniform but the drag coefficient, $`\gamma `$, is space dependent. Using in Eq.(2) the limited expansion, $`D(x+\alpha \mathrm{\Delta }x)D(x)+\alpha (dD/dx)\mathrm{\Delta }x`$ with $`\mathrm{\Delta }x=\pm \sqrt{2D(x)\mathrm{\Delta }t}`$, yields the algorithm for a Brownian walker with a position dependent diffusion coefficient:
$$x(t+\mathrm{\Delta }t)=x(t)\pm \sqrt{2D(x(t))\mathrm{\Delta }t}+\alpha \frac{dD}{dx}\mathrm{\Delta }t.$$
(3)
Depending on the value of $`\alpha `$, this model has very different implications concerning the equilibrium distribution of the Brownian walkers, their individual drift($`x(t)x(0)`$) and their net flux.
Averaging Eq.(3) over a large number of walkers shows the average position of a Brownian walker is no longer zero since the diffusion gradient term acts as an external force leading to particle drift. If this gradient is assumed to be constant, this drift increases linearly with time as
$$x(t)x(0)=\alpha \frac{dD}{dx}\mathrm{\Delta }t.$$
(4)
A first intuitive, but misleading, idea would be to conclude that the particles will migrate in the direction of the diffusion gradient, leading, therefore, to a concentration gradient. This is actually incorrect as we now show. Starting with a uniform particle distribution $`\rho _0`$, we check if this corresponds to an equilibrium state by determining the particle flux through an imaginary surface, $`S`$, placed perpendicular to the diffusion coefficient gradient at coordinate $`x`$ (see Fig. 1). During a time interval $`\mathrm{\Delta }t`$, all particles crossing $`S`$ from the left (or right) are half of those included in the volume $`SL_{right}`$ ($`SL_{left}`$) where $`L_{right}`$ ($`L_{left}`$) is the right (left) step terminating at $`x`$ taken by a walker during $`\mathrm{\Delta }t`$. The net particle flux to the right will thus be
$$J=\frac{\rho _0}{2}\frac{SL_{right}SL_{left}}{S\mathrm{\Delta }t}.$$
(5)
Equation (2) allows computing the length of these two steps which both end at the same point x:
$$L_{_{left}^{right}}=\sqrt{2D(x)\mathrm{\Delta }t}\pm (\alpha 1)\frac{dD}{dx}\mathrm{\Delta }t,$$
(6)
leading to the particle flux
$$J=\rho _0(1\alpha )\frac{dD}{dx}.$$
(7)
As a result, in the situation of maximum drift where $`\alpha =1`$, this flux will vanish (see left part of Fig. 1), meaning that the uniform particle distribution corresponds to an equilibrium. According to Boltzmann, this should correspond to an isothermal situation, the diffusion coefficient gradient arising only from a pure hydrodynamic effect, the spatial dependence of the drag coefficient $`\gamma `$. For all the other values of $`\alpha `$, the flux will be negative, leading to a concentration gradient of the particles in the direction opposite to that of the diffusion coefficient gradient. The maximum flux (and zero drift) is obtained for $`\alpha =0`$, as shown on the right part of figure 1. Equation (7) may be generalized to the case where the particle distribution is position dependent:
$$J=(1\alpha )\rho (x)\frac{dD}{dx}D\frac{d\rho }{dx},$$
(8)
which leads to the well known generalization of Fickโs law for a space dependent diffusion coefficient.
To confirm these results, we performed simulations of Brownian walkers following algorithm (3). We found that only the $`\alpha =1`$ case leads to a uniform distribution of particles with no net flux through any given surface while, at the same time, the average individual positions exhibit a drift in the direction of the diffusion coefficient gradient according to Eq(4). This situation of โdrift without fluxโ may be compared to the equilibrium situation of Brownian particles subjected to an external force, such as their weight: If one follows the motion of individual particles, an average downwards drift is observed; however, there is no net flux because of the vertical concentration gradient. In our isothermal case, the drift of individual particles from lower to larger D(x) region does not lead to a net flux because particles in the larger D(x) region diffuse further than particles in the lower D(x) region. This physical situation imposes the choice of $`\alpha =1`$ in algorithm (3), so that a particle coming from a low $`D(x)`$ region makes a right step just equal to the left step of that particle coming from a high $`D(x)`$ region and arriving at the same point (see left part of Fig. 1).
We have set up an experiment to check this drift without flux prediction where particles, observed under a microscope, undergo Brownian motion in a confined geometry. When the confinement, $`e`$, is of the order of the particle size, $`a`$, the diffusion coefficient strongly depends on the value of $`e/a`$. As the confinement was position dependent, we observed Brownian walkers in an isothermal situation but with a spatially dependent $`D(x)`$ due to a purely hydrodynamic effect.
Polystyrene spheres, of radius $`a=1\mu `$m, were suspended in a mixture of $`H_2O+D_2O`$ so as to cancel any sedimentation effects. Addition of a surfactant ($`2.2g/l`$ of SDS) helped minimize particle aggregation or adhesion to the walls. A drop of this mixture was placed between a flat disk and a planar convex lens, Fig. (2), of curvature radius $`R=15.5`$mm separated by an elastic O-ring. The flat and convex surfaces were then brought into contact at the center of the cell by gently squeezing the elastic joint, the remaining air providing the necessary sample compressibility. The spacing, $`e`$, between the flat and curved wall depends on the distance $`r`$ from the center of the cell as $`e=r^2/2R`$. The contact between the two walls as well as the dependence of $`e`$ on $`r`$ were carefully measured by monitoring the Newton rings observed under the microscope. We used as a light source a new super-radiant diode whose coherence length is less than $`100\mu `$m. This was important as this coherence length was long enough to observe the desired Newton rings, but short enough to avoid any other interference patterns due to all the cell interfaces, which were visible with an ordinary diode laser and which completely masked the relevant signal. The horizontal Brownian motion of the polystyrene balls was observed through a microscope equipped with a long range objective of magnification $`50X`$, followed by a CCD camera coupled to the microscope via an eyepiece of magnification $`8X`$. The video signal was processed in real time by a computer, which recorded, every $`3`$s the horizontal position, size and shape of all objects in a rectangular frame $`65\mu \mathrm{m}X100\mu `$m. This time interval was long enough to allow the image analysis of all particles present in the frame and small enough so that the particleโs average displacement was only a fraction of their diameter.
As the particlesโ confinement, $`e`$, was related to their distance, $`r`$, from the center of the cell, we were able to explore different confinement regions by moving the observation frame in the horizontal plane. The explored $`e`$ varied from $`2.5\mu `$m to $`11\mu `$m. The vertical Brownian motion of the particles over this small vertical range could not be monitored. However, we took that motion into account when interpreting the data by averaging the particleโs vertical position over the confinement range. The volume fraction of polystyrene balls, of the order of $`1\%`$, was chosen so as to allow the monitoring of a fairly large number of particles at the same time (around $`30`$ for $`e=3\mu `$m) to improve the statistics in the data analysis. Following the particlesโ positions from one frame to another, the program analyzed a great number of trajectories (more than $`10^5`$ for each run). When two particles got closer than twice their diameter, the program treated them as โdimersโ, their trajectories as โmonomersโ were ended at that time and were no longer used to determine the diffusion coefficient or the drift. In that way, the role of particle interactions, which have a range smaller than a couple of particle diameters, could be safely ignored.
The confinement dependence of the diffusion coefficient was determined as follows. A given observation frame was divided into 3 zones (see inset, Fig. (2)), each corresponding to an approximately constant $`e`$. For each zone we averaged the particleโs displacement squared, either in the x, or in the y directions, as a function of time, and checked that indeed they followed the usual diffusion law,
$$x^2=y^2=2D_{}(e)t.$$
(9)
By moving the frame to different locations, we were able to explore a range extending from $`e/2a=1.2`$ to $`11`$. The bulk diffusion constant, $`D_0`$, was determined using the well-known relation $`D_0=kT/6\pi \eta a`$, where $`\eta =0.99X10^3`$SI is the viscosity of the water- heavy water mixture, yielding $`D_0=1.92X10^{13}`$m<sup>2</sup>/s. The experimental values for $`D_{}/D_0`$ are shown in Fig. (3) (white squares) and fit remarkably well the available theoretical predictions (black dots) using the collocation method averaged over all the possible vertical positions $`z`$ of the particle for a given $`e`$, i.e. with $`az(ea)`$. For comparison, we also plotted (solid line) the analytical solution obtained using the Faxen expression for the position dependent drag of a particle moving parallel to a single wall, then adding the effect of each of the two walls, and averaging over the vertical position $`z`$. This solution clearly overestimates the reduction of the diffusion coefficient of a particle trapped between two parallel walls, particularly as the relative confinement, $`ฯต=e/2a`$ reaches its lower limit 1. We also plotted (dashed lines) in Fig. (3) the analytical expression for the drag of a particle trapped just in the middle of two parallel planes. The impossibility of averaging over $`z`$ an expression only known for $`z=e/2`$ explains the observed discrepancy which goes to zero as the relative confinement approaches its limit, $`1`$, where the z-average becomes irrelevant.
To demonstrate the existence of an individual drift of the particles, we fixed the center of the observation frame at a position $`y=0`$ and $`x=300\mu `$m, corresponding to an average $`e/2a=1.5`$ so that all particles present in the frame were outside the excluded volume (i.e. $`e2a`$), and had a diffusion coefficient with the largest $`x`$ dependence, but no y dependence (to first order). For the determination of $`x(t)x(0)`$ and $`y(t)y(0)`$, each trajectory was divided into independent paths lasting a time $`t`$, each contributing to the evaluation of the average drift during time $`t`$. The results are shown in Fig. (4), and reveal a drift in the Brownian walker position along the $`x`$ direction, and none in the $`y`$ direction along which the diffusion coefficient may be considered as constant. The statistics of the results clearly deteriorates as time increases: After recording trajectories for typically a dozen hours, more than a hundred thousand independent segments contributed to the determination of the drift at short times, whereas only up to a few thousand independent segments were left for $`t=200`$s. This is due to the fairly high particle concentration which lowers the lifetime of a โmonomerโ (time during which a particle doesnโt approach another one to within 2 particle diameters), but which was chosen as a compromise to have good statistics at short times while allowing us to follow each particle during a reasonable time, fixed at $`200`$s. It should be pointed out that in order to avoid any bias in the statistics, for a trajectory segment to be valid and included in the statistics, the position of a walker at instant $`t=0`$ had to be inside a region $`15\mu `$m away from the edges of the observation frame. This condition ensured that after diffusing for $`200`$s, the walker had less than $`0.5\%`$ chance to have covered $`15\mu `$m, and was thus still present in the observation frame. Failure to impose this condition resulted in the observation of a spurious drift, in the opposite direction, due to an artificial selection of walkers because of the experimental boundary conditions (limits of the observation frame).
To compare our experimental results with the theoretical predictions, Eq. (4), we evaluated the diffusion gradient encountered by the walkers present in the observation frame. It is important to realize that as a walker moves a distance $`dx`$ (see round inset in Fig. 2), its diffusion coefficient $`D_{}`$ varies first because the confinement varies (path (1) parallel to the bottom wall, i.e. at constant $`z`$), and second because at constant confinement, the particleโs height $`z`$ changes (path (2)). Adding both contributions and averaging over the vertical position z of the walker yields:
$$\frac{dD_{}(e,z)}{dx}_z=\frac{x}{R}\frac{D_{}(e,z)}{e}_z+\frac{x}{2R}\frac{D_{}(e,z)}{z}_z.$$
(10)
Using our experimental data and results of the collocation method, we found $`dD_{}(e,z)/dx_z2.2X10^9`$m/s. This value of the slope is used to plot the straight dotted line on Fig. 4. The experimental data are thus in good agreement with the predicted drift corresponding to the expected $`\alpha =1`$ value.
Finally, to claim drift without flux, it is not sufficient only to demonstrate the drift: we must also demonstrate the absence of flux. If a flux was due to the observed drift $`dD/dx`$, we would expect a radially outwards flux of $`\rho dD/dx`$ particles, which would empty our observation screen in less than a day. Furthermore, if this flux were to be balanced by a concentration gradient, one can show that a concentration change by $`30\%`$ over a distance of $`60\mu `$m would be necessary. Experimentally, we observed no flux and no concentration gradient over a period of a week or more, which is consistent with the Boltzmann requirement of a uniform concentration in the absence of a temperature gradient.
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# Limits of Crystallographic Methods for Detecting Space Topology
## 1 Introduction
The search for the topology of the spatial sections of the universe has made tremendous progress in the past years (see LachiรจzeโRey and Luminet (1995) for an introduction to the subject and early references, Luminet and Roukema (1999) and Uzan et al. (1999b) for a review of the late developments). Methods using twoโdimensional data sets, such as the cosmic microwave background maps planned to be obtained by the MAP and Planck Surveyor satellite missions, and threeโdimensional (3D) data sets, such as galaxy, cluster and quasar surveys with redshifts, have been developped.
Following our previous works (Lehoucq et al. 1996, Lehoucq et al. 1999 and Uzan et al. 1999a), we focus on the topological information that can be extracted from a 3D catalogue of cosmic objects. The key point of all the topologyโdetecting methods is based on the โtopological lens effectโ, i.e. on the fact that if the spatial section of the universe has at least one characteristic size smaller than the spatial scale of the catalogue, then one should observe different images of the same object. The original crystallographic method (Lehoucq et al. 1996) used a pair separation histogram (PSH) depicting the number of pairs of catalogueโs objects having the same threeโdimensional separations in the universal covering space, with the idea that spikes should stand out dramatically at characteristic lengths related to the size of the fundamental domain and to the holonomies of space. However, following a remark by Weeks (1998), we proved (Uzan et al. 1999a and Lehoucq et al. 1999) that sharp spikes can emerge in the PSH only if the holonomies of space are Clifford translations โ a result independently derived by Gomero et al. (1998). As a consequence, the PSH method does not apply to the detection of topology in universes with hyperbolic spatial sections.
Since then, various generalisations of the crystallographic method were proposed. Fagundes and Gaussman (1998) suggested to map the differences between the PSH of a simulated catalogue in a compact hyperbolic universe and the PSH in the corresponding simplyโconnected universe having the same distribution of objects and cosmological parameters. They noticed sharp oscillations on the scale of the bin width, modulated by a broad oscillation on the scale of the curvature radius. However in Uzan et al. (1999b), we calculated on one hand the differential PSH between a simulated distribution in a compact hyperbolic Weeks space and a similar distribution in the simplyโconnected hyperbolic space $`๐^\mathrm{๐}`$, on the other hand the differential PSH between two different distributions (with the same total number of objects) in the simplyโconnected hyperbolic space $`๐^\mathrm{๐}`$. Both curves (fig. 10 of Uzan et al. 1999b) exhibit the same pattern of sharp and broad oscillations, which shows that the topological significance of such a pattern is highly doubtful.
Fagundes and Gaussman (1999) next proposed a modified crystallographic method in which the topological images in simulated catalogues are pulled back to the fundamental domain before the set of 3D distances is calculated. The distribution of pair distances is expected to be peaked around zero. The main drawback of this method is that the nature of the signal strongly depends on the topological type, on the orientation of the fundamental domain and on the position of the observer inside the latter. Thus the pullback method can be useful only if the exact topology is already known. On the other hand, Gomero et al. (1999) introduced a mean PSH, aimed to reduce the statistical noise that may mask the topological signature. They claimed that such a technique should detect the contribution of non-translational isometries to the topological signal, whatever the curvature of space, but the applicability of their method to real data has not been demonstrated.
In order to improve the signalโtoโnoise ratio inherent to PSHโs , instead of reducing the statistical noise we can enhance the topological signal by collecting all distance correlations into a single index. Such was our purpose in Uzan et al. 1999a, where we reformulated cosmic crystallography as a collecting-correlated-pairs method (CCP). The CCP technique rests on the basic fact that in a multiply connected universe, equal distances appear more often than just by chance, whatever the curvature of space and the nature of holonomies. We also showed that the extraction of a topological signal drastically depends on the rather accurate knowledge of the cosmological parameters, namely the density parameter, $`\mathrm{\Omega }_0`$, and the cosmological constant parameter, $`\mathrm{\Omega }_{\mathrm{\Lambda }0}`$. This is due to the fact that all cosmic crystallography methods require the determination of 3D separations between two any images: observations use redshifts for determining the coordinate distance in the universal covering space (in addition to the angular positions on the celestial sphere), and the redshiftโdistance relation involves the cosmological parameters. Conversely, the detection of a topological signal would help to determine accurately the curvature parameters (see also Roukema and Luminet 1999).
However it is to be recognized that, in the framework of cosmic crystallography, the application of numerical simulations to real data rest on two idealized assumptions, namely:
1. all objects are strictly comoving
2. the catalogues of observed objects (quasars, galaxy clusters) are complete.
Although quoted in LachiรจzeโRey and Luminet (1995), the quantitative effects of such simplifications have not been fully discussed in the literature. In Lehoucq et al. (1996), the authors qualitatively argued that the distortion due to peculiar velocities was negligible, and they performed numerical simulations to study the influence of the angular resolution of the surveys. In Roukema (1996), the influence of the astrophysical uncertainties (mainly of spectroscopic measurements and of peculiar velocities) on a method trying to find quintuplets of quasars with the same geometry (and thus which may be topological images) was evaluated. It was shown that, in such a case, the most serious uncertainty comes from the radial peculiar velocities. In Uzan et al. (1999a), we discussed the effect of the peculiar velocities and of the errors in the spatial position arising from the imprecisions on the values of the cosmological parameters. To finish, the errors on the determination of the position and of the peculiar velocities were discussed in Roukema (1996), and the way that the constraints on $`\mathrm{\Omega }_0`$ and $`\mathrm{\Omega }_{\mathrm{\Lambda }0}`$ depend on the redshifts of multiple topological images and on their radial and tangential separations was calculated in Roukema and Luminet (1999).
The goal of the present article is to have a critical attitude on the methods we have developped so far, by listing all the sources of observational uncertainties and by evaluating their effects on the theoretical efficiencies of the PSH method (Lehoucq et al. 1996) and of the CCP method (Uzan et al. 1999a) respectively. In ยง 2, we discuss the nature of uncertainties and in ยง 3 the numerical methods used to estimate their effects on the topological signal. We then quantify the magnitude of each effect in the Euclidean and hyperbolic cases (ยง 4) (we postpone the case of elliptic spaces to a further study). In conclusion (ยง 5) we compare our results to the performances of current and future observational programs aimed to detect the topology of the universe.
Notations and descriptions
We keep the notations of our previous articles (Lehoucq et al. 1999, Uzan et al. 1999a). The local geometry of the universe is described by a FriedmannโLemaรฎtre metric
$$\mathrm{d}s^2=\mathrm{d}t^2+a^2(t)\left[\mathrm{d}\chi ^2+f^2(\chi )\left(\mathrm{d}\vartheta ^2+\mathrm{sin}^2\vartheta \mathrm{d}\phi ^2\right)\right],$$
(1)
where $`a`$ is the scale factor, $`t`$ the cosmic time, $`\chi `$ the comoving radial distance, and $`f(\chi )=(\mathrm{sin}\chi ,\chi ,\mathrm{sinh}\chi )`$ according to the sign of the space curvature $`k=(+1,0,1)`$. The time evolution of $`a`$ is obtained by solving the Friedmann and the conservation equations
$`H^2`$ $`=`$ $`\kappa {\displaystyle \frac{\rho }{3}}{\displaystyle \frac{k}{a^2}}+{\displaystyle \frac{\mathrm{\Lambda }}{3}},`$ (2)
$`\dot{\rho }`$ $`=`$ $`3H(\rho +P)`$ (3)
where $`\kappa 8\pi G/c^4`$, $`\rho `$ and $`P`$ are respectively the matter density and the pressure, $`\mathrm{\Lambda }`$ is the cosmological constant and $`H\dot{a}/a`$ is the Hubble parameter (with a dot referring to a time derivative). We also use the standard parameters
$$\mathrm{\Omega }\frac{\kappa \rho }{3H^2}\text{and}\mathrm{\Omega }_\mathrm{\Lambda }\frac{\mathrm{\Lambda }}{3H^2}.$$
(4)
This completely specifies the properties and the dynamics of the universal covering space. In the following, we assume that we are in the matter dominated era, so that $`P=0`$ and $`\rho a^3`$.
The topology of the spatial sections is described by the fundamental domain, a polyhedron whose faces are pairwise identified by the elements $`g`$ of the holonomy group $`\mathrm{\Gamma }`$ (see LachiรจzeโRey and Luminet (1995) for the details).
## 2 Observational uncertainties
Early methods used to put a lower bound on the size of the universe and based on the direct recognition of multiple images of given objects โ such as our Galaxy (Sokolov and Schvartsman 1974), the Coma cluster (Gott 1980) โ had to face a major drawback due to the fact that the same object would be seen at different lookback times. Thus, evolution effects such as photometric or morphologic changes will, in most cases, render the identification of objects impossible (see however Roukema and Edge, 1997). Statistical techniques such as the PSH method (Lehoucq et al. 1999) and the CCP method (Uzan et al. 1999a) are free from such biases.
The two main sources of uncertainties for all the statistical methods trying to detect the topology of the universe in catalogues of cosmic objects are
1. the errors in the positions of observed objects, which can be separated into:
1. the uncertainty in the determination of the redshifts due to spectroscopic imprecision; such an effect is purely experimental and exists even if the objects are strictly comoving
2. the uncertainty in the position due to peculiar velocities of objects, which induce peculiar redshift corrections
3. the uncertainty in the cosmological parameters, which induces an error in the determination of the radial distance (via the redshift โ distance relation)
4. the angular displacement due to gravitational lensing by large scale structure.
2. the incompleteness of the catalogue, which has two main origins:
1. selection effects implying that some objects are missing from the catalogue
2. the partial coverage of the celestial sphere, due either to the presence of the galactic plane, or to the fact that surveys are performed within solid angles much less than $`4\pi `$.
Effects of peculiar velocities
A peculiar velocity has two effects on the determination of the position of a cosmic object:
1. An integrated effect, coming from the fact that if a galaxy has a proper velocity, then its true position differs from its comoving one.
2. An instantaneous effect, due to the fact that the radial component of the proper velocity will add to the cosmological redshift an extra term, $`\mathrm{\Delta }z_{Dop}`$.
Assuming that the vector velocity $`v`$ of a galaxy is constant, which is indeed a good enough approximation to estimate the effects of peculiar velocities, a galaxy at redshift $`z`$ has moved from its comoving position by a comoving distance
$$\delta \mathrm{}=v\tau [z]$$
(5)
where $`\tau [z]`$ is the lookโback time, obtained by integrating the photon geodesic equation
$$\tau [z]=\frac{1}{H_0}_{\frac{1}{1+z}}^1\frac{\mathrm{d}x}{\sqrt{\mathrm{\Omega }_{\mathrm{\Lambda }0}x^2+(1\mathrm{\Omega }_0\mathrm{\Omega }_{\mathrm{\Lambda }0})+\mathrm{\Omega }_0/x}}.$$
(6)
This reduces to the well known expression
$$\tau [z]=\frac{2}{3H_0}\left(1\frac{1}{(1+z)^{3/2}}\right)$$
(7)
when $`\mathrm{\Omega }_0=1`$ and $`\mathrm{\Omega }_{\mathrm{\Lambda }0}=0`$.
Now, $`\delta \mathrm{}`$ can be expressed in terms of a peculiar redshift, $`\mathrm{\Delta }z_{\mathrm{pos}}`$, and of an angular displacement, $`\mathrm{\Delta }\theta _{\mathrm{pos}}`$, which depend on the galaxy velocity as
$`\mathrm{\Delta }z_{\mathrm{pos}}`$ $`=`$ $`{\displaystyle \frac{v_{}}{c}}{\displaystyle \frac{c\tau [z]}{\chi ^{}[z]}}`$ (8)
$`\mathrm{\Delta }\theta _{\mathrm{pos}}`$ $`=`$ $`{\displaystyle \frac{|v_{}|}{c}}{\displaystyle \frac{c\tau [z]}{\chi [z]}}`$ (9)
where a prime refers to derivative with respect to $`z`$, $`v_{}`$ is the component of the velocity $`v`$ with respect to the lineโofโsight direction $`\gamma `$, and $`v_{}`$ is the nonโradial peculiar velocity, defined as
$$v_{}v.\gamma \text{and}v_{}vv_{}\gamma .$$
(10)
$`\chi [z]`$ is the observer area distance, given by
$$\chi [z]=\frac{c}{a_0H_0}_{\frac{1}{1+z}}^1\frac{\mathrm{d}x}{\sqrt{\mathrm{\Omega }_{\mathrm{\Lambda }0}x^4+(1\mathrm{\Omega }_0\mathrm{\Omega }_{\mathrm{\Lambda }0})x^2+\mathrm{\Omega }_0x}}.$$
(11)
In figures 1 and 2, we respectively depict the variations of $`\mathrm{\Delta }z_{\mathrm{pos}}`$ and $`\mathrm{\Delta }\theta _{\mathrm{pos}}`$ as a function of the redshift.
Concerning the instantaneous effect ii), the redshift uncertainty $`\mathrm{\Delta }z_{\mathrm{Dop}}`$ can be related to the galaxy proper velocity as follows.
If we consider the trajectory of a photon $`x^\mu (s)`$, $`s`$ being the affine parameter along the null geodesic, the relation between the emission (E) wavelength $`\lambda _\mathrm{E}`$ and the reception (R) wavelength $`\lambda _\mathrm{R}`$ can be expressed as
$$\frac{\lambda _\mathrm{R}}{\lambda _\mathrm{E}}=\frac{\left(k^\mu u_\mu ^{\mathrm{obs}}\right)_\mathrm{E}}{\left(k^\mu u_\mu ^{\mathrm{gal}}\right)_\mathrm{R}}$$
(12)
where $`k^\mu dx^\mu /ds`$ is the tangent vector to the photon geodesic; $`u_\mu ^{\mathrm{obs}}`$ and $`u_\mu ^{\mathrm{gal}}`$ are respectively the 4โvelocity of the observer and of the galaxy. Neglecting the perturbations of the metric and of the matter and focusing on the Doppler effect, one can easily show that
$$k^\mu u_\mu =\frac{1}{a(t)}\left[1+\gamma ^i\frac{v_i}{c}\right],$$
(13)
where $`\gamma ^i`$ is the direction in which the galaxy is observed and $`v^i`$ is its (Newtonian) velocity. The observed (spectroscopic) redshift, $`z_{\mathrm{obs}}`$, and the cosmological redshift, $`z_{\mathrm{cosm}}`$, respectively defined by
$$1+z_{\mathrm{obs}}\frac{\lambda _R}{\lambda _E},1+z_{\mathrm{cosm}}\frac{a(t_E)}{a(t_R)}$$
(14)
are thus related by
$$1+z_{\mathrm{obs}}=(1+z_{\mathrm{cosm}})\left[1+\gamma ^i\frac{(v_i^{\mathrm{gal}}v_i^{\mathrm{obs}})}{c}\right].$$
(15)
Assuming that we can substract the Earthโs velocity, an error $`\mathrm{\Delta }zz_{\mathrm{obs}}z_{\mathrm{cosm}}`$ in the determination of $`z_{\mathrm{obs}}`$ can be interpreted in terms of a radial peculiar velocity given by
$$\mathrm{\Delta }z_{\mathrm{Dop}}=\frac{v_{}^{\mathrm{gal}}}{c}(1+z_{\mathrm{cosm}}).$$
(16)
As expected, a galaxy receding from us, i.e. with $`\gamma ^iv_i^{\mathrm{gal}}>0`$ will have an additional redshift, whereas a galaxy drawing nearer to us will induce a Doppler blueshift correction.
The considerations above are also useful for discussing the effects of angular displacement $`\mathrm{\Delta }\theta _{\mathrm{gl}}`$ due to gravitational lensing by large scale structure (Mellier 1999, van Waerbeke et al. 2000). A typical value $`\mathrm{\Delta }\theta _{\mathrm{gl}}1`$ arcsec corresponds to a peculiar velocity a few km.s<sup>-1</sup> as long as the redshift is less than 5 (as can be seen from Fig. 2). Thus effect $`(A4)`$ is much less important than effect $`(A2)`$.
Effects of catalogue incompleteness
Objects can be missing from a catalogue due to strong evolution effects (this is particularly the case for quasars), and to absorption of light by intergalactic gas or dust in some sky directions. Another well known selection effect is the Malmquist bias: statistical samples of astronomical objects which are limited by apparent magnitude have mean absolute magnitudes which are different than those of distanceโlimited samples. An apparent magnitudeโlimited sample contains, if luminosity function has a finite width, some very luminous objects which, in spite of their large distances, can jump the apparent-magnitude limit of the catalogue. While one looks further and further out one finds more and more luminous objects. The consequence for an apparentโmagnitude catalog (which is the case with galaxy catalogues) is that dwarf galaxies fade out quickly with distance, and finally at the largest distances the extremely luminous and rare galaxies are the only ones which can enter the catalogue.
In addition to effects (B1) and (B2), some noise spikes may appear due to gravitational clustering of objects. The large scale distribution of galaxies shows a variety of landscapes containing voids, walls, filaments and clusters in a complex 3D sponge-like pattern. For instance, since there are many galaxies in clusters, the distances associated to cluster-cluster separations may appear as fake spikes in the PSH for galaxies. Typical cluster โ cluster separations are $`130h^1`$ Mpc (Guzzo et al. 1999). Such effects currently occur in Nโbody simulations that show clustering (Park and Gott 1991). The trouble can be relieved if galaxy clusters instead of galaxies are used as typical objects for probing the topology, although in that case supercluster-supercluster separations might also introduce noise spikes (at a lower level).
## 3 Numerical implementation
In the following, we perform numerical simulations to evaluate separately the magnitudes of the effects listed above. As usual, we start from a random distribution of cosmic objects in the fundamental domain. In a first step we generate a complete catalogue of comoving objects by unfolding the distribution in the universal covering space. We refer to this catalogue as the ideal catalogue since it would correspond to ideal observations. In a second step we introduce the various errors $`(A_i)(B_i)`$ in order to build more realistic catalogues which depart from the ideal one.
### 3.1 Uncertainties on positions
To study numerically the errors on the positions, we first assume that the cosmological parameters are known with good enough accuracy, so that we do not discuss (A3). Then we perform the following calculations.
1. To evaluate the effect of the observational imprecision, we give to each object of the ideal catalogue (built on the assumption that the objects are strictly comoving) a redshift error to be added to its ideal redshift. We assume that the distribution of the redshift error is Gaussian, with mean value $`\overline{z}=0`$ and dispersion $`\mathrm{\Delta }z`$.
Note that, since $`\mathrm{\Delta }z`$ is absolute, the relative error will be less important when we deal with catalogues of higher redshift objects.
2. To discuss (A2), we associate a peculiar velocity to each point of the catalogue before unfolding. Hence, there will be correlations between the velocities of two topological images. It follows that, assuming that the time evolution of the peculiar velocity is small, the observed velocities of two topological images of the same object can have different directions but similar magnitudes (see e.g. Roukema and Bajtlik 1999). As in the previous case, we assume that the velocity distribution is Gaussian. Moreover, velocities at different points of space are assumed to be uncorrelated both in magnitude and direction, which is indeed not strictly the case in real data since large scale streaming motions of galaxies have been observed (Strauss and Willick 1995).
More precisely, to generate the catalog of topological images taking into account the peculiar velocities, we proceed as follows:
1. We generate a random collection of points $`M_i`$ (named original sources) uniformly distributed inside the fundamental domain. Each point is assigned a velocity $`v_i`$ according to a Gaussian distribution with mean $`\overline{v}`$ and dispersion $`\mathrm{\Delta }v`$.
2. We unfold this catalogue of original sources to obtain the set of points $`M_{k,i}`$, images of $`M_i`$ by the holonomies $`g_k\mathrm{\Gamma }`$.
3. Each image $`M_{k,i}`$ has a redshift $`z_{k,i}`$ corresponding to a lookโback time $`\tau _{k,i}\tau (z_{k,i})`$ calculated with formula (6).
4. The final catalogue accounting for the peculiar velocities of the original sources is obtained by applying the $`g_k`$ to $`M_i+v_i\times \tau _{k,i}`$ for all couples $`(k,i)`$.
Since $`\mathrm{\Delta }z`$ can be interpreted either as an uncertainty on the redshift determination or as due to a radial peculiar velocity, both errors (A1) and (A2) can be investigated with the same calculations. The interpretation of $`\mathrm{\Delta }z`$ in terms of a radial peculiar velocity is useful to compare the orders of magnitude of (A1) and (A2). In real data, the error due to (A2) is expected to be greater than the error due to (A1).
### 3.2 Catalogue incompleteness
Starting from an ideal catalogue, we simulate two kinds of incomplete catalogues:
1. To account for effect (B1), we randomly throw out $`p\%`$ of the objects from the ideal catalogue. In our various runs we vary $`p`$ while keeping approximately constant the number of catalogue objects, which means that we must increase the number of โoriginalโ objects in the fundamental domain when $`p`$ is increasing.
2. To simulate effect (B2), we generate a catalogue limited in solid angle by selecting only the objects that lie in a beam of aperture $`\theta `$. The sky coverage is related to $`\theta `$ via
$$q\%=\frac{1}{2}\left(1\mathrm{cos}(\theta /2)\right).$$
(17)
Again, we keep constant the number of objects in the catalogue when we vary $`\theta `$. Table 1 below gives typical numbers.
As a matter of fact, more than the aperture angle, the depth of the survey will be critical for the multiplication factor, since within a beam of given angle, if the redshift cut-off is great enough to encompass a distance (in the universal covering space) N times greater than the size of the fundamental domain, at least N topological images will be expected in the beam.
## 4 General Results
### 4.1 Euclidean spaces
We first consider the Euclidean case and we apply the PSH crystallographic method as described in Lehoucq et al. (1996), restricting to a typical situation where $`\mathrm{\Omega }_0=0.3`$ and $`\mathrm{\Omega }_{\mathrm{\Lambda }0}=0.7`$. We choose the topology of the universe to be a cubic 3โtorus $`T_1`$ (see e.g. LachiรจzeโRey and Luminet (1995) for description), with identification length $`L=3,000\mathrm{Mpc}`$ for a Hubble parameter $`H_0=75\mathrm{km}.\mathrm{s}^1\mathrm{Mpc}^1`$. In the various runs, the number of objects in the catalogue is kept constant at $`8500`$ (this is the order of magnitude of the number of objects in current quasar catalogues). We examine separately the effects of errors in position due to redshift uncertainty $`\mathrm{\Delta }z`$ and peculiar velocities $`\mathrm{\Delta }v`$, and the effects of catalogue incompleteness due to selection effects and partial sky coverage. Each of these effects will contribute to spoil the sharpness of the topological signal. For a given depth of the catalogue, namely a redshift cut-off $`z_{cut}`$, we perform the runs to look for the critical value of the error at which the topological signal fades out.
Figure 3 gives the critical redshift error $`\mathrm{\Delta }z_l`$ above which the topological spikes disappear.
The effect of peculiar velocities is very weak, since we find that the peculiar velocity must exceed $`\mathrm{\Delta }v_l=10,000\mathrm{km}.\mathrm{s}^1`$ when $`z_{cut}=1`$ and $`\mathrm{\Delta }v_l=40,000\mathrm{km}.\mathrm{s}^1`$ when $`z_{cut}=5`$ in order to make the topological signal disappear. As already pointed out in LachiรจzeโRey and Luminet (1995), the observed peculiar velocities of galaxies have typical values much less that $`\mathrm{\Delta }v_l`$.
The effects of catalogue incompleteness are summarized in tables 2 and 3. The topological signal would be destroyed only for a very large rejection percentage or a small aperture angle.
### 4.2 Hyperbolic spaces
We now turn to universes with hyperbolic spatial sections and apply the CCP method as described in Uzan et al. (1999a). The obtention of a topological signal (the soโcalled CCP index) strongly depends on the correct determination of the cosmological parameters. In Uzan et al. (1999a) we discussed the problems arising from the spanning of the cosmological parameters space with a required โaccuracy binโ $`\epsilon `$. Latest independent constraints on these cosmological parameters from the cosmic microwave background (de Bernardis et al. 2000), the study of supernovae (Efstathiou et al. 1999), of large scale structure at $`z=2`$ (Roukema and Mamon, 2000) and of gravitational lensing (Mellier 1999), make us hope that we can restrict further the parameters space to apply efficiently the CCP method. In the following, we assume that the cosmological parameters are known with the required accuracy.
We choose the topology of the universe to be described by a Weeks manifold (Weeks, 1985) (see also Lehoucq et al. (1999) for the numerical implementation of this topology), assuming the cosmological parameters given by $`\mathrm{\Omega }_0=0.3`$ and $`\mathrm{\Omega }_{\mathrm{\Lambda }0}=0`$. In such a case the number of copies of the fundamental domain within the horizon is about 190. However our simulated catalogues are much smaller than the horizon volume. We fix the number of objects to $`1300`$, which is a compromise between a realistic catalogue population and a reasonable computing time.
Again we perform the tests by varying the errors $`\mathrm{\Delta }z`$ and $`\mathrm{\Delta }v`$ around their average values $`\overline{z}=0`$ and $`\overline{v}=0`$ until when the CCP index falls down to noise level.
Contrarily to the Euclidean case, both slight changes in redshift and in peculiar velocity induce errors in position which dramatically eradicate the topological signal: an error of only $`50\mathrm{kms}^1`$ in velocity and an error in redshift of the order of the bin accuracy $`\epsilon 10^6`$ drown the CCP index into noise.
The incompleteness effects are less dramatic, as shown in figures 4 and 5. Again, for a given redshift cutโoff, we performed the runs by varying $`p`$ and $`\theta `$ in order to find the critical values at which the topological signal disappears.
## 5 Conclusions and perspectives
Our numerical results have now to be compared with the precisions of present experimental 3D data, and to the performances of observational programs started or expected to be achieved in the next decade.
At present day, a typical precision practical for the spectroscopic uncertainty is $`\mathrm{\Delta }z0.001`$ for an object such as a quasar. In clusters, spectroscopic redshifts can be found very precisely for individual galaxies.
Concerning peculiar velocities, the typical dispersion velocity is $`1000\mathrm{km}.\mathrm{s}^1`$ in rich clusters. The X-ray velocity of the peak of the X-ray distribution would also provide a way to estimate the true cluster redshift, including the peculiar velocity of the cluster as a whole. For a quasar, a conservative upper limit to the peculiar velocity, assuming the quasar to be at the centre of a galaxy, can be taken as $`\mathrm{\Delta }z0.002`$. So, from an experimental point of view, the uncertainties on the redshifts will be dominated by effect (A2), namely the peculiar velocities, rather than by the spectroscopic imprecision.
The main limitation of present 3D samples is the small volume of existing redshift data. Future surveys will significantly improve both the redshift cutโoff and the sky coverage. For instance, the Sloan Digital Sky Survey (SDSS) (Loveday 1998) will map in detail one-quarter of the entire sky, determining the positions and absolute brightnesses of more than 100 million celestial objects. It will also measure distances to more than a million galaxies and quasars. More precisely, SDSS will map a contiguous $`\pi `$ steradians area in the north Galactic cap, up to a limiting magnitude 23 for two thirds of the observing time, together with three southern stripes centred at RA $`\alpha =5^{}`$ and with central declinations of $`\delta =+15^{},0^{}`$ and $`10^{}`$ for the remaining one third of the time. Concerning the distance determinations, $`10^6`$ galaxies and $`10^5`$ quasars will be observed spectroscopically with a resolution $`\mathrm{\Delta }z/z5.10^4`$. The main galaxy sample will consist of $`900,000`$ galaxies up to magnitude 18, with a median redshift $`z0.1`$. A second galaxy sample will consist in $`100,000`$ luminous red galaxies to magnitude 19.5 with a median redshift $`z0.5`$. Precision on redshifts can be estimated for the reddest galaxies to $`\mathrm{\Delta }z0.02`$. A sample of $`100,000`$ quasars will be observed, an order of magnitude larger than any existing quasar catalogue. The complete survey data will become public by 2005. We can also mention the ESOโVLT Virmos Deep Survey Project (Lefรจvre 2000), a comprehensive imaging and redshift survey of the deep universe based on more than 150 000 redshifts.
In the field of Xโray observations, the XMM satellite (Arnaud 1996) will provide deep insight on Xโray galaxy clusters and active galactic nuclei. Also, the XEUS project under study by ESA (Parmar et al. 1999) will be a longโterm Xโray observatory at 1 keV, with a limiting sensitivity around 250 times better than XMM, allowing XEUS to study the properties of galaxy groups at $`z=2`$ and active galactic nuclei at $`z3`$.
In the present paper we have investigated how the various observational uncertainties will spoil the topological signal expected to arise in the ideal situation when crystallographic methods are applied to complete catalogues of perfectly comoving objets with zero proper velocities and whose 3Dโpositions are known with infinite accuracy. By numerical simulations we have introduced random errors for each possible uncertainty, and we varied the parameters to determine the limits at which the topological signal vanishes.
Our numerical calculations of the spoiling effects due to the various uncertainties $`(A_i)(B_i)`$ clearly show that the crystallographic methods are stable (in the sense that the topological signal is robust when data depart from the ideal ones) in the Euclidean case, but highly unstable in the hyperbolic case. Indeed in a small multiโconnected flat space, realistic values of peculiar velocities of objects, errors in redshift determinations and partial sky coverage will not make the PSH method to fail. This can be understood by the fact that the topological images of a given object are related together by Clifford translations which enhance the topological signal. On the contrary, in a compact hyperbolic model, holonomies are not Clifford translations. The topological signal, built as a CCP index, is destroyed as soon as small errors are introduced in the position of objects, due either to peculiar velocities or to redshift measurement imprecision.
The same kind of critical analysis should be made with the 2D topologyโdetecting methods based on the analysis of CMB data. For instance, in the pairs of matched circles method (Cornish et al. 1998), it would be necessary to investigate how deviations to the ideal situation, such as a non zero thickness of the last scattering surface or peculiar motions of the emitting primordial plasma regions, would alter the pattern of perfectly matched circles.
To conclude, let us comment on the future of observational cosmic topology, at the light of observational constraints on the curvature parameters recently provided by BOOMERanG and MAXIMA balloon measurements (de Bernardis et al. 2000, Hanany et al. 2000). Under specific assumptions such as a cold dark matter model and a primordial density fluctuation power spectrum, the range of values allowed for the energy-density parameter $`\mathrm{\Omega }=\mathrm{\Omega }_0+\mathrm{\Omega }_{\mathrm{\Lambda }0}`$ is restricted to $`0.88<\mathrm{\Omega }<1.12`$ with $`95\%`$ confidence. This means that the curvature radius of space is as least as great as the radius of the observable universe (delimited by the last scattering surface). Such results, if accepted, leave open all three cases of space curvature as well as most of multiโconnected topologies.
A strictly flat space is quite improbable. Even inflationary scenarios predict a value of $`\mathrm{\Omega }`$ asymptocally close to 1, but not strictly equal. From a topological point of view, a strictly Euclidean space would be interesting since we have shown that the PSH method is robust enough to provide a topological signal even when realistic uncertainties on the data are taken into account.
For compact hyperbolic spaces, if we accept the recent observational constraints $`0.88<\mathrm{\Omega }_0+\mathrm{\Omega }_{\mathrm{\Lambda }0}<1`$, a number of spaceforms such as the Weeks or the Thurston manifolds still have topological lengths smaller than the horizon size (Weeks: SnapPea). However the topological lens effects would be weaker in spaceforms with $`\mathrm{\Omega }`$ close to 1 than in spaceforms with $`\mathrm{\Omega }0.3`$ (see figs. 2 and 3 of Lehoucq et al. 1999). Furthermore, only the CCP method can be applied for detecting the topology, and the present article shows that the topological signal will fall to noise level as soon as uncertainties smaller than the experimental ones are taken into account in the simulations.
Eventually, elliptical spaceforms appear to be the most interesting case. On a theoretical point of view, as far as we know, no inflationary model is able to drive the density parameter to a value greater than 1, so that if space happened to be really elliptical, new models should be built in order to explain, e.g., the primordial fluctuations spectrum or the horizon problem. From a topological point of view, since the volumes of (all closed) spaceforms are not bounded below (the order of the holonomy group can be arbitrarily large), one can always find an elliptical space which fits into the Hubble radius even if $`1<\mathrm{\Omega }_0+\mathrm{\Omega }_{\mathrm{\Lambda }0}<1.12`$. On the other hand, the holonomies of such spaces are Clifford translations, so that we can hope to apply the robust PSH method to detect a topological signal. This will be the purpose of our subsequent paper.
Acknowledgements: It is a pleasure to thank N. Aghanim, R. Juszkiewicz, Y. Mellier and B. Roukema for discussions.
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# IFUP-TH/2000-13 Linked-Cluster Expansion of the Ising Model
## 1 Introduction
The high-temperature (strong-coupling) series expansion is one of the most successful tools for the study of physical systems near a critical point.
High-temperature series are analytic; the radius of convergence is usually quite large, often reaching the boundary of the high-temperature phase. This property allows the application of powerful techniques of resummation and analytical continuation , which can yield very precise and reliable results, provided that long series are available. It is therefore worthwhile to push the computation of high-temperature series as far as our algorithms and computers allow.
The most successful technique for the computation of high-temperature series of $`3D`$ spin models is the linked-cluster expansion (LCE), which is well suited for the fully computerized approach required to reach very high orders of the expansion.
Several detailed discussion of the LCE appeared in the literature. Wortisโ review covers most of the basic topics, and provides many graphical rules fit for algorithmic implementation. Nickel performed a remarkable computation for the $`3D`$ Ising model on the bcc lattice , which had not been surpassed until the present work; in a more detailed paper, Nickel and Rehr also present several clever algorithms which we found very useful . Lรผscher and Weisz, describing their application of the LCE to lattice field theory, also provide several important implementation hints .
Unfortunately, the notation found in the literature is by no means uniform. Therefore we will review the relevant aspects of the LCE, which can be found in Ref. , not only to make our paper more self-contained, but also to explain notations carefully and to remark the correspondence with Refs. . We will follow the notations of Ref. whenever possible.
The paper is organized as follows:
Sect. 2 introduces the relevant graph theory concepts and definitions.
Sect. 3 presents the generalized Ising model we focus on.
Sects. 4, 5, and 6 review the LCE, with special focus on the two-point Greenโs functions.
Sect. 7 describes our algorithm for the computation of three-point and higher Greenโs functions.
Sect. 8 is devoted to programming details.
Sect. 9 displays a (small) selection of the series generated.
Forthcoming papers will be devoted to the analysis of the series, using the techniques presented in Ref. , and to the generation and analysis of the series for $`XY`$ systems.
We will not give proofs of our formulae. The only nontrivial step in the proofs of Sects. 57 is to show that the symmetry factors compensate exactly the different number of contributions that may appear on the two sides of the equations; it is typically a straightforward, if tedious, exercise in combinatorics. The proofs of Sects. 46 are given or sketched in Ref. . The proofs of Sect. 7 are especially easy, since the symmetry factor of a 1-irreducible tree graph is always 1.
## 2 Graphology
In this section we introduce a number of graph theory concepts relevant for the LCE. We refer the reader to Ref. for a comprehensive introduction to the subject.
A graph is a set of vertices and edges (also named links or bonds in the literature). Each edge $`l`$ is incident with two distinct vertices, its extrema (we do not allow the extrema to coincide); the set of extrema will be denoted by $`l`$; we write $`l=\{i(l),f(l)\}`$; the choice of an โinitialโ and a โfinalโ vertex is arbitrary. Two vertices are adjacent if they are the extrema of the same edge. We will denote the number of vertices and edges of a graph by $`v`$ and $`e`$ respectively. We also consider arcs or oriented edges, incident out of the initial vertex $`i(l)`$ into the final vertex $`f(l)`$.
The valence $`n(i)`$ of a vertex $`i`$ is the number of edges incident with $`i`$.
An $`r`$-rooted graph is a graph with $`v`$ vertices, $`vr`$: $`r`$ roots or external vertices and $`vr`$ internal vertices. We will assign the indices $`1`$, โฆ, $`r`$ to the roots and the indices $`r+1`$, โฆ, $`v`$ to the internal vertices. In the drawings, roots will be denoted by open dots and internal vertices by filled dots.
Two $`r`$-rooted graphs are isomorphic if there exists a one-to-one correspondence $`\pi `$ of their internal vertices and edges such that the incidence relations are preserved, i.e., $`(\pi (l))=\{\pi (i(l)),\pi (f(l))\}`$ ($`\pi (i)=i`$ for the roots). From now on, we will identify isomorphic graph, and silently assume that all sets of graphs we define contain only non-isomorphic graphs.
The symmetry factor $`S(๐ข)`$ of an $`r`$-rooted graph $`๐ข`$ is the number of isomorphisms of $`๐ข`$ into itself, i.e. the number of of permutations of internal vertices and edges preserving the incidence relations.
We will also consider $`p`$-ordered $`r`$-rooted graphs, $`pr`$, i.e. the classes of rooted graphs isomorphic up to permutations of $`rp`$ roots. We will assign indices $`1`$, โฆ, $`p`$ to the fixed roots and indices $`p+1`$, โฆ, $`r`$ to the roots which can be permuted. The symmetry factor $`S(๐ข)`$ is defined as the number of isomorphisms of $`๐ข`$ into itself, including permutation of the roots $`p+1`$, โฆ, $`r`$. The symmetry factor divided by $`(rp)!`$ is called the modified symmetry factor $`S_E(๐ข)`$ (cfr. Ref. ); it need not be an integer. The $`r`$-rooted graphs defined above are ordered ($`r`$-ordered); unless otherwise specified, we will assume that graphs are ordered. 0-ordered graphs are unordered.
We will also discuss ($`p`$-ordered) $`r`$-rooted graphs whose edges and/or vertices are assigned a label; let us consider e.g. the case of an edge label $`a(l)`$ and vertex label $`b(i)`$. Two labelled graphs $`(๐ข,a,b)`$ and $`(๐ข^{},a^{},b^{})`$ are equivalent (isomorphic) if there exists an isomorphism of $`๐ข`$ into $`๐ข^{}`$ such that $`a`$ is mapped into $`a^{}`$ and $`b`$ is mapped into $`b^{}`$. The symmetry factor $`S(๐ข,a,b)`$ is the number of isomorphisms of $`(๐ข,a,b)`$ into itself.
A pair of vertices $`i`$ and $`j`$ is connected if there exists a sequence of vertices $`k_1`$, โฆ, $`k_n`$, with $`k_1=i`$ and $`k_n=j`$, and a sequence of edges $`l_1`$, โฆ, $`l_{n1}`$ such that $`l_a=\{k_a,k_{a+1}\}`$, $`a=1`$, โฆ, $`n1`$. A graph is connected if every pair of its vertices is connected. In the following, unless otherwise noted, we will assume that every graph is connected.
A sequence of distinct vertices $`k_1`$, โฆ, $`k_n`$ and distinct edges $`l_1`$, โฆ, $`l_n`$ is called a loop of length $`n`$ if $`l_a=\{k_a,k_{a+1}\}`$, $`a=1`$, โฆ, $`n1`$, and $`l_n=\{k_n,k_1\}`$. The number of independent loops (also known as cyclomatic number) of a connected graph is $`ev+1`$.
A connected graph is called a tree graph if it contains no loop. A tree graph has $`v1`$ edges.
## 3 The model
We wish to compute the high-temperature (HT) expansion of the $`q`$-point functions of a generalized Ising model on a $`D`$-dimensional Bravais lattice $`\mathrm{\Lambda }`$; notice that Bravais lattices enjoy inversion symmetry at each lattice site. The model is defined by the generating functional
$$\mathrm{exp}(W[h])=\frac{1}{Z}\underset{i}{}\left[๐\varphi _if(\varphi _i)\mathrm{exp}(h_i\varphi _i)\right]\mathrm{exp}\left(K\underset{ij}{}\varphi _i\varphi _j\right),$$
(1)
where $`\varphi _i`$ is a scalar field, $`f`$ is an even non-negative function or distribution decreasing faster than $`\mathrm{exp}(\varphi ^2)`$ as $`\varphi \mathrm{}`$, normalized by the condition
$$๐\varphi f(\varphi )=1,$$
$`K=\beta J`$, the sum runs over all pairs of nearest neighbours, and the normalization $`Z`$ is fixed by requiring $`W[0]=0`$.
The connected $`q`$-point function (denoted by $``$ in Ref. ) at zero magnetic field is defined by
$$G_q(x_{i_1},\mathrm{},x_{i_q})\varphi _{i_1}\mathrm{}\varphi _{i_q}^{\mathrm{con}}|_{h=0}=\frac{^qW[h]}{h_{i_1}\mathrm{}h_{i_q}}|_{h=0},$$
(2)
where $`x_i(x_1(i),\mathrm{},x_D(i))`$ is the coordinate vector of the lattice site $`i`$. $`G_q`$ is invariant under the lattice symmetry group, including (discrete) translations, and under permutation of its arguments; it is customary to write $`G_2(x_{i_1},x_{i_2})`$ in the form $`G_2(x_{i_1}x_{i_2})`$. We will apply the LCE to the computation of $`G_q`$.
## 4 Unrenormalized expansion
Let us parametrize the potential $`f`$ in terms of the bare vertices $`\mu _0(2n)`$, defined by the generating function
$$\mathrm{exp}\left[\underset{n}{}\frac{\mu _0(2n)}{(2n)!}h^{2n}\right]=๐\varphi f(\varphi )\mathrm{exp}(h\varphi ).$$
(3)
These quantities are named bare semi-invariants and denoted by $`M_{n}^{}{}_{}{}^{0}`$ in Ref. ; they are named cumulant moments and denoted by $`\mu _{2n}`$ in Ref. ; ฮผ0(2n)=(2n1)!!
m
2nconsubscript๐02๐double-factorial2๐1subscriptsuperscript
m
con2๐\mu_{0}(2n)=(2n-1)!!\,\vbox{\hbox{\hbox{}\kern 2.5pt$\scriptstyle\circ$}\hbox{$m$}}^{\rm con}_{2n} in the notations of Ref. . Without loss of generality, we can rescale $`\varphi `$ and $`K`$ to fix $`\mu _0(2)=1`$.
For a generic $`r`$-rooted graph $`๐ข`$, we define the bare external vertex factor
$$V_0^{(e)}(n_1,\mathrm{},n_r;๐ข)=\underset{i=1}{\overset{r}{}}\mu _0(n(i)+n_i),$$
(4)
the bare internal vertex factor
$$V_0^{(i)}(๐ข)=\underset{i=r+1}{\overset{v}{}}\mu _0(n(i)),$$
(5)
and the bare edge factor
$$L_0(x_1,\mathrm{},x_v;๐ข)=\underset{l=1}{\overset{e}{}}[K\theta (x_{i(l)}x_{f(l)})],$$
(6)
where $`\theta (x)=\delta (x1)`$ and $`x`$ is the lattice distance between 0 and $`x`$.
It is convenient to focus on the contributions of $`r`$-rooted graphs to $`q`$-point functions. To this purpose we introduce the auxiliary $`r`$-point functions $`X`$, whose unrenormalized LCE is
$$X_r(x_1,\mathrm{},x_r;n_1,\mathrm{},n_r)=\underset{๐ข^{(r,0)}}{}\underset{x_{r+1},\mathrm{},x_v}{}\frac{V_0^{(e)}(n_1,\mathrm{},n_r;๐ข)V_0^{(i)}(๐ข)L_0(x_1,\mathrm{},x_v;๐ข)}{S(๐ข)},$$
(7)
where $`^{(r,0)}`$ is the set of all $`r`$-rooted connected graphs.
$`X`$ is invariant under simultaneous permutation of coordinates and valences:
$$X_r(x_1,\mathrm{},x_r;n_1,\mathrm{},n_r)=X_r(x_{\pi (1)},\mathrm{},x_{\pi (r)};n_{\pi (1)},\mathrm{},n_{\pi (r)}),$$
but not over independent permutation of coordinates and valences. Furthermore, $`X`$ is invariant over the lattice symmetry group, e.g. it is translation-invariant:
$$X_r(x_1,\mathrm{},x_r;n_1,\mathrm{},n_r)=X_r(x_1+x,\mathrm{},x_r+x;n_1,\mathrm{},n_r).$$
$`X_1(x;n)`$ is independent of $`x`$, and it will be denoted by $`X_1(n)`$; it will also be denoted by $`\mu (n)`$ in its role of renormalized vertex. $`X_2(x_1,x_2;n_1,n_2)`$ only depends on the difference $`x_2x_1`$, and it will be denoted by $`X_2(x_2x_1;n_1,n_2)`$. Since, by invariance under space inversion, $`X_2(x;n_1,n_2)=X_2(x;n_1,n_2)`$, we also have $`X_2(x;n_1,n_2)=X_2(x;n_2,n_1)`$.
The sum over the location of internal vertices of the $`\theta `$ functions is by definition the (free) lattice embedding number of $`๐ข`$ with fixed roots:
$$\underset{x_{r+1},\mathrm{},x_v}{}\underset{l=1}{\overset{e}{}}\theta (x_{i(l)}x_{f(l)})=E(x_1,\mathrm{},x_r;๐ข).$$
(8)
Therefore
$$X_r(x_1,\mathrm{},x_r;n_1,\mathrm{},n_r)=\underset{๐ข^{(r,0)}}{}\frac{V_0^{(e)}(n_1,\mathrm{},n_r;๐ข)V_0^{(i)}(๐ข)K^{e(๐ข)}E(x_1,\mathrm{},x_q;๐ข)}{S(๐ข)}.$$
(9)
Finally, the $`q`$-point functions are computed as:
$`G_q(x_1,\mathrm{},x_q)={\displaystyle \underset{\mathrm{partitions}}{}}X_r(x_{i_{11}},\mathrm{},x_{i_{r1}};u_1,\mathrm{},u_r){\displaystyle \underset{l=1}{\overset{r}{}}}\delta _{u_l}(x_{i_{l1}},\mathrm{},x_{i_{lu_l}}),`$ (10)
where the $`q`$-point delta function is
$$\delta _1(x_1)=1,\delta _2(x_1,x_2)=\delta (x_1x_2),\mathrm{},\delta _q(x_1,\mathrm{},x_q)=\underset{l=2}{\overset{q}{}}\delta (x_1x_l)$$
and $`\{\{i_{11},\mathrm{},i_{1u_1}\},\mathrm{},\{i_{r1},\mathrm{},i_{ru_r}\}\}`$ is a generic partition of $`\{1,\mathrm{},q\}`$ into $`r`$ sets of size $`u_1`$, โฆ, $`u_r`$. We will call a root bearing a factor of $`\delta _u`$ a $`u`$-th order root.
The two-point function is simply
$$G_2(x)=X_2(x;1,1)+\delta (x)X_1(2).$$
## 5 Vertex-renormalized expansion
The graph $`๐ข\backslash i`$ is obtained by deleting from $`๐ข`$ the vertex $`i`$, i.e. by removing $`i`$ and all edges incident with $`i`$. A vertex $`i`$ of a rooted graph $`๐ข`$ is called an articulation point if there exist vertices of $`๐ข\backslash i`$ not connected to a root. A rooted graph is called 1-irreducible if it does not contain any articulation point.
Any $`r`$-rooted ($`r>1`$) connected graph $`๐ข`$ can be decomposed in a unique way into a 1-irreducible $`r`$-rooted 1-skeleton $`๐ฎ`$ and a 1-rooted 1-decoration for each vertex; $`๐ข`$ is reconstructed by decorating each vertex, identifying the root of its decoration with the vertex; an example is presented in Fig. 1.
Since the only 1-irreducible 1-rooted graph is the single-vertex graph, we use a different definition for 1-rooted graphs. A 1-rooted graph is called a 1-skeleton if it has no articulation points except the root. Any 1-rooted connected graph $`๐ข`$ can be decomposed in a unique way into a 1-rooted 1-skeleton $`๐ฎ`$ and a 1-rooted 1-decoration for each vertex except the root, which is left undecorated.
A 1-rooted connected graph is called a 1-insertion if $`๐ข\backslash 1`$ is connected.
The LCE can be reorganized by summing together all contributions from graphs having the same 1-skeleton, incorporating 1-decorations into renormalized vertices $`\mu (n)=X_1(n)`$ (named semi-invariants and denoted by $`M_n`$ in Ref. ; $`\mu (n)=(n1)!!m_n`$ in the notations of Ref. ). The unrenormalized LCE of $`\mu (n)`$ is given by Eq. (7).
The $`r`$-point function can be computed restricting the sum in Eq. (7) or (9) to the (much smaller) set $`^{(r,1)}`$ of 1-irreducible $`r`$-rooted graphs:
$`X_r(x_1,\mathrm{},x_r;n_1,\mathrm{},n_r)`$ $`=`$ $`{\displaystyle \underset{๐ข^{(r,1)}}{}}{\displaystyle \underset{x_{r+1},\mathrm{},x_v}{}}{\displaystyle \frac{V_1^{(e)}(n_1,\mathrm{},n_r;๐ข)V_1^{(i)}(๐ข)L_0(x_1,\mathrm{},x_v;๐ข)}{S(๐ข)}}`$ (11)
$`=`$ $`{\displaystyle \underset{๐ข^{(r,1)}}{}}{\displaystyle \frac{V_1^{(e)}(n_1,\mathrm{},n_r;๐ข)V_1^{(i)}(๐ข)K^{e(๐ข)}E(x_1,\mathrm{},x_r;๐ข)}{S(๐ข)}},`$
where the internal and external renormalized vertex factors are
$$V_1^{(e)}(n_1,\mathrm{},n_r;๐ข)=\underset{i=1}{\overset{r}{}}\mu (n(i)+n_i)$$
(12)
and
$$V_1^{(i)}(๐ข)=\underset{i=r+1}{\overset{v}{}}\mu (n(i)).$$
(13)
Eq. (7) requires a sum over all connected graphs, and therefore it is impractical for the computation of $`\mu (n)`$ at large orders of the LCE. We introduce the renormalized moments $`q(n)`$ (named self-fields and denoted by $`G_n`$ in Ref. ), defined by
$$q(n)=\underset{๐ข_{n,\mathrm{in}}^{(1,0)}}{}\underset{x_2,\mathrm{},x_v}{}\frac{V_0^{(i)}(๐ข)L_0(x_1,\mathrm{},x_v;๐ข)}{S(๐ข)},$$
(14)
where $`_{n,\mathrm{in}}^{(1,0)}`$ is the set of 1-insertions with root of valence $`n`$. The following equations hold:
$$q(n)=\underset{๐ข_{n,\mathrm{in}}^{(1,1)}}{}\underset{x_2,\mathrm{},x_v}{}\frac{V_1^{(i)}(๐ข)L_0(x_1,\mathrm{},x_v;๐ข)}{S(๐ข)},$$
(15)
where $`_{n,\mathrm{in}}^{(1,1)}=_{n,\mathrm{in}}^{(1,0)}^{(1,1)}`$, and $`^{(1,1)}`$ is the set of 1-rooted 1-skeletons;
$$\mu (n)=\mu _0(n)+\underset{s=1}{\overset{\mathrm{}}{}}\frac{1}{s!}\underset{l_1=1}{\overset{\mathrm{}}{}}\mathrm{}\underset{l_s=1}{\overset{\mathrm{}}{}}q(l_1)\mathrm{}q(l_s)\mu _0(n+l_1+\mathrm{}+l_s).$$
(16)
Since $`q(2n1)=0`$, and odd values of $`n`$ and $`l_i`$ do not contribute to Eq. (16). $`q`$ and $`\mu `$ can now be computed recursively in parallel order by order in $`K`$, since, once Eq. (15) is expanded in powers of $`K`$ and truncated, the coefficient of the highest power of $`K`$ of $`q`$ in the l.h.s. depends only on coefficients of lower powers of $`K`$ of $`\mu `$ in the r.h.s.
## 6 Edge-renormalized expansion
A pair of distinct vertices $`i`$ and $`j`$ of a rooted graph $`๐ข`$ is called an articulation pair if there exist vertices of $`๐ข\backslash i\backslash j`$ not connected to a root, or if $`i`$ and $`j`$ are joined by more than one edge. A rooted graph is called 2-irreducible if it does not contain any articulation pair.
Any 1-irreducible $`r`$-rooted ($`r>2`$) graph $`๐ข`$ can be decomposed in a unique way into a 2-irreducible $`r`$-rooted 2-skeleton $`๐ฎ`$ and a 1-irreducible 2-rooted 2-decoration for each edge (oriented in a canonical way, e.g. by choosing $`i(l)<f(l)`$); $`๐ข`$ is reconstructed by replacing each edge with its decoration, identifying the first and second decoration root with the initial and final vertex of the edge respectively. An example is shown in Fig. 2.
We use a different definition for 2-rooted graphs, since the only 2-irreducible 2-rooted graph is the bond graph (no internal vertices and only one edge); the roots of all other 1-irreducible 2-rooted graphs are an articulation pair. We call a 1-irreducible 2-rooted graph a 2-skeleton if it does not contain any articulation pair except the pair consisting of the two roots. Any 1-irreducible 2-rooted graph can be decomposed in a unique way into a 2-rooted 2-skeleton $`๐ฎ`$ and a 1-irreducible 2-rooted 2-decoration for each edge except edges connecting the roots, which are left undecorated.
The LCE can be reorganized by summing together all contributions from graphs having the same 2-skeleton, incorporating all 2-decorations into renormalized edges.
We start by decomposing Eq. (11) for $`r>1`$ into
$`W_r(x_1,\mathrm{},x_r;n_1,\mathrm{},n_r)`$ $`=`$ $`{\displaystyle \underset{๐ข_{n_1,\mathrm{},n_r}^{(r,1)}}{}}{\displaystyle \underset{x_{r+1},\mathrm{},x_v}{}}{\displaystyle \frac{V_1^{(i)}(๐ข)L_0(x_1,\mathrm{},x_v;๐ข)}{S(๐ข)}},`$ (17)
$`X_r(x_1,\mathrm{},x_r;n_1,\mathrm{},n_r)`$ $`=`$ $`{\displaystyle \underset{s_1,\mathrm{},s_r}{}}\left[{\displaystyle \underset{i=1}{\overset{r}{}}}\mu (n_i+s_i)\right]W_r(x_1,\mathrm{},x_r;s_1,\mathrm{},s_r),`$ (18)
where $`_{n_1,\mathrm{},n_r}^{(r,1)}`$ is the set of 1-irreducible $`r`$-rooted graphs with roots of valence $`n_1`$, โฆ, $`n_r`$. $`W_r`$ enjoys the same symmetry properties of $`X_r`$. $`W_1(x,n)`$ is undefined.
$`W_r`$ can be computed by assigning an initial and a final valence $`i_l`$, $`f_l`$ to each oriented edge of a 2-rooted graph; the valence $`i_l`$ is incident with $`i(l)`$ and $`f_l`$ is incident with $`f(l)`$:
$`W_r(x_1,\mathrm{},x_r;n_1,\mathrm{},n_r)={\displaystyle \underset{๐ข^{(r,2)}}{}}{\displaystyle \underset{x_{r+1},\mathrm{},x_v;i_1,\mathrm{},i_e,f_1,\mathrm{},f_e}{}}{\displaystyle \underset{i=1}{\overset{r}{}}}\delta (n_i\nu _i(i_1,\mathrm{},i_e,f_1,\mathrm{},f_e;๐ข))`$
$`\times {\displaystyle \frac{V_2^{(i)}(i_1,\mathrm{},i_e,f_1,\mathrm{},f_e;๐ข)L_2(x_1,\mathrm{},x_v;i_1,\mathrm{},i_e,f_1,\mathrm{},f_e;๐ข)}{S(๐ข)}},`$ (19)
where $`^{(r,2)}`$ is the set of 2-irreducible $`r`$-rooted graphs, $`\nu _i`$ is the sum of all the valences incident with the vertex $`i`$,
$$V_2^{(i)}(i_1,\mathrm{},i_e,f_1,\mathrm{},f_e;๐ข)=\underset{i=r+1}{\overset{v}{}}\mu (\nu _i(i_1,\mathrm{},i_e,f_1,\mathrm{},f_e;๐ข)),$$
(20)
and
$$L_2(x_1,\mathrm{},x_v;i_1,\mathrm{},i_e,f_1,\mathrm{},f_e;๐ข)=\underset{l=1}{\overset{e}{}}W_2(x_{i(l)}x_{f(l)};i_l,f_l).$$
(21)
For $`r=2`$ Eqs. (19) and (21) require a slight modification: we define $`^{(2,2)}`$ as the set of 2-rooted 2-skeletons; for edges incident with at least one root, we replace $`W_2(x_1x_2;i,f)`$ with $`K\theta (x_1x_2)\delta _{i1}\delta _{f1}`$.
The sum over graphs can be restricted to a subset of $`^{(r,2)}`$; we will discuss here the case $`r=2`$; the next Section will be devoted to the case $`r3`$. Let us start by classifying 1-irreducible 2-rooted graphs into several classes.
An internal vertex $`i`$ of a 1-irreducible 2-rooted graph $`๐ข`$ is called a nodal point if the roots of $`๐ข\backslash i`$ are not connected. A 1-irreducible 2-rooted graph is nodal (also named articulated or separable in the literature) if it contains one or more nodal points; otherwise it is non-nodal.
A 1-irreducible 2-rooted graph is simple if $`๐ข\backslash 1\backslash 2`$ is connected, and 1 is not adjacent to 2. By definition, all nodal graphs are simple. A 1-irreducible 2-rooted graph is a ladder graph if it is not simple and it is not the bond graph.
A 1-irreducible 2-rooted graph is elementary if it is both simple and non-nodal.
We have divided 1-irreducible 2-rooted graphs into four disjoint classes: bond, nodal, ladder, and elementary graphs. Let us separate the contributions to $`W_2`$ according to the four classes:
$$W_2(x;n_1,n_2)=W_2^{\mathrm{bo}}(x;n_1,n_2)+W_2^{\mathrm{no}}(x;n_1,n_2)+W_2^{\mathrm{la}}(x;n_1,n_2)+W_2^{\mathrm{el}}(x;n_1,n_2),$$
i.e. bond, nodal, ladder, and elementary contributions respectively.
The bond contribution is trivial. Nodal contributions can be factorized into a product of non-nodal contributions:
$`W_2^{\mathrm{no}}(x;n_1,n_2)={\displaystyle \underset{x_3;i_1,i_2}{}}W_2^{\mathrm{nn}}(x_3;n_1,i_1)\mu (i_1+i_2)W_2^{\mathrm{nn}}(xx_3;i_2,n_2)`$
$`+{\displaystyle \underset{x_3,x_4;i_1,i_2,i_3,i_4}{}}W_2^{\mathrm{nn}}(x_3;n_1,i_1)\mu (i_1+i_2)`$
$`\times W_2^{\mathrm{nn}}(x_4x_3;i_2,i_3)\mu (i_3+i_4)W_2^{\mathrm{nn}}(xx_4;i_4,n_2)+\mathrm{},`$ (22)
which can be written recursively as
$$W_2^{\mathrm{no}}(x;n_1,n_2)=\underset{x_3;i_1,i_2}{}W_2^{\mathrm{nn}}(x_3;n_1,i_1)\mu (i_1+i_2)W_2(xx_3;i_2,n_2).$$
(23)
Likewise, ladder contributions can be factorized into a product of non-ladder contributions:
$$W_2^{\mathrm{la}}(x;n_1,n_2)=\underset{s=2}{\overset{\mathrm{}}{}}\frac{1}{s!}\underset{i_1,\mathrm{},i_s,f_1,\mathrm{},f_s}{}\delta \left(n_1_{t=1}^si_t\right)\delta \left(n_2_{t=1}^sf_t\right)\underset{t=1}{\overset{s}{}}W_2^{\mathrm{nl}}(x,i_t,f_t).$$
(24)
Elementary contributions can be computed by setting $`r=2`$ into Eq. (19), and restricting the sum to $`_{\mathrm{el}}^{(2,2)}`$, the set of elementary 2-rooted 2-skeletons. The sum can be further restricted to $`_{0,\mathrm{el}}^{(2,2)}`$, the set of unordered elementary 2-rooted 2-skeletons, provided that we replace $`S(๐ข)`$ with $`S_E(๐ข)`$ and we symmetrize the result:
$`W_2^{\mathrm{el}}(x;n_1,n_2)={\displaystyle \frac{1}{2}}{\displaystyle \underset{๐ข_{0,\mathrm{el}}^{(2,2)}}{}}{\displaystyle \underset{x_3,\mathrm{},x_v;i_1,\mathrm{},i_e,f_1,\mathrm{},f_e}{}}{\displaystyle \underset{l=1}{\overset{2}{}}}\delta (n_i\nu _i(i_1,\mathrm{},i_e,f_1,\mathrm{},f_e;๐ข))`$
$`\times {\displaystyle \frac{V_2^{(i)}(i_1,\mathrm{},i_e,f_1,\mathrm{},f_e;๐ข)L_2(x_1,\mathrm{},x_v;i_1,\mathrm{},i_e,f_1,\mathrm{},f_e;๐ข)}{S_E(๐ข)}}+n_1n_2.`$ (25)
The last ingredient we need for a fully edge-renormalized expansion are the renormalized vertices $`\mu (n)`$; they can be computed by combining Eq. (16) with
$$q(n_1+n_2)=\frac{n_1!n_2!}{(n_1+n_2)!}(W_2^{\mathrm{no}}(0;n_1,n_2)+W_2^{\mathrm{el}}(0;n_1,n_2)),$$
(26)
reflecting the fact that the contributions to $`q(n_1+n_2)`$ in Eq. (15) can be obtained from the contributions to $`W_2(0;n_1,n_2)`$ in Eq. (17) by identifying the roots of the 2-rooted graph, i.e. by suppressing the second root and reattaching all the edges incident with it to the first root, provided that the roots are not adjacent.
Expanding in powers of $`K`$ Eqs. (23), (24), (25), (26), and (16), we can compute $`W_2^{\mathrm{el}}(x;n_1,n_2)`$, $`W_2^{\mathrm{no}}(x;n_1,n_2)`$, $`W_2^{\mathrm{la}}(x;n_1,n_2)`$, $`q(n)`$, and $`\mu (n)`$ in parallel order by order in $`K`$. The only step which involves a summation over graphs is Eq. (25), where the sum only runs over $`_{0,\mathrm{el}}^{(2,2)}`$, a relatively small set.
## 7 Three-point and higher functions
A vertex (internal or external) $`i`$ of a 1-irreducible $`r`$-rooted graph $`๐ข`$ is called a nodal point if $`๐ข\backslash i`$ is not connected. By definition of 1-irreducibility, each connected component of $`๐ข\backslash i`$ must contain at least one root. A 1-irreducible $`r`$-rooted graph $`๐ข`$ is nodal if it contains one or more nodal points; otherwise it is non-nodal.
In the rest of this section, we will assume that every graph is 2-irreducible.
A nodal point $`j`$ of a 2-irreducible $`r`$-rooted graph $`๐ข`$ is called a tree-insertion point if at least one of the connected components of $`๐ข\backslash j`$ is a tree graph. The order $`t`$ of $`j`$ is the number of roots of $`๐ข`$ contained in all tree graph components of $`๐ข\backslash j`$, plus 1 if $`j`$ itself is a root. A 2-irreducible $`r`$-rooted graph is called compact if it contains no tree-insertion points.
Let us consider a 2-irreducible $`r`$-rooted graph $`๐ข`$. We generate a compact 2-irreducible graph $`๐ข^{}`$, called the compact kernel of $`๐ข`$, by removing all the tree graphs attached to every tree-insertion point, and promoting all internal tree-insertion points to root. $`๐ข`$ is obtained by attaching a 2-irreducible tree graph to each root of $`๐ข^{}`$. If $`๐ข^{}`$ is not a tree graph (and therefore it has at least 3 roots), the decomposition is unique; otherwise, $`๐ข`$ itself is a tree graph. An example is shown in Fig. 3.
By summing all contribution of 2-irreducible graphs with the same compact kernel, we can write $`W_q`$ as
$`W_q(x_1,\mathrm{},x_q;n_1,\mathrm{},n_q)=W_q^{\mathrm{tr}}(x_1,\mathrm{},x_q;n_1,\mathrm{},n_q)`$
$`+{\displaystyle \underset{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\mathrm{partitions}}{r>2}}}{}}{\displaystyle \underset{y_1,\mathrm{},y_r,i_1,\mathrm{},i_r}{}}W_r^{\mathrm{co}}(y_1,\mathrm{},y_r;i_1,\mathrm{},i_r){\displaystyle \underset{l=1}{\overset{r}{}}}Y_{u_l+1}(y_l,x_{i_{l1}},\mathrm{},x_{i_{lu_l}};i_l,n_{i_{l1}},\mathrm{},n_{i_{lu_l}}),`$
(27)
where $`W_q^{\mathrm{tr}}(x_1,\mathrm{},x_q;n_1,\mathrm{},n_q)`$ is the tree graph contribution to $`W_q`$, $`W_r^{\mathrm{co}}(y_1,\mathrm{},y_r;i_1,\mathrm{},i_r)`$ is the compact graph contribution to $`W_r`$, and
$`Y_{t+1}(y,x_1,\mathrm{},x_t;i,n_1,\mathrm{},n_t)`$ $`=`$ $`{\displaystyle \underset{j}{}}\mu (i+j)W_{t+1}^{\mathrm{tr}}(y,x_1,\mathrm{},x_t;j,n_1,\mathrm{},n_t)`$
$`+`$ $`{\displaystyle \underset{j}{}}\delta (yx_1)\delta (n_1ij)W_t^{\mathrm{tr}}(x_1,\mathrm{},x_t;j,n_2,\mathrm{},n_t),`$
$`Y_2(y,x;i,n)`$ $`=`$ $`{\displaystyle \underset{j}{}}\mu (i+j)W_2^{\mathrm{tr}}(y,x;j,n)+\delta (yx)\delta (ni)`$ (28)
(the first and second term correspond to an internal and external $`t`$-tree-insertion point respectively); notice that $`W_2^{\mathrm{tr}}=W_2`$. Eq. (27) can be combined with Eq. (10) to give
$`G_q(x_1,\mathrm{},x_q)=G_q^{\mathrm{tr}}(x_1,\mathrm{},x_q)`$
$`+{\displaystyle \underset{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\mathrm{partitions}}{r>2}}}{}}{\displaystyle \underset{y_1,\mathrm{},y_r,i_1,\mathrm{},i_r}{}}W_r^{\mathrm{co}}(y_1,\mathrm{},y_r;i_1,\mathrm{},i_r){\displaystyle \underset{l=1}{\overset{r}{}}}Z_{u_l+1}(y_l,x_{i_{l1}},\mathrm{},x_{i_{lu_l}};i_l),`$ (29)
where
$`Z_{t+1}(y,x_1,\mathrm{},x_t;i)`$ $`=`$ $`{\displaystyle \underset{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\mathrm{partitions}}{\mathrm{of}\{1,\mathrm{},t\}}}}{}}{\displaystyle \underset{s_1,\mathrm{},s_r}{}}Y_{r+1}(y,x_{i_{11}},\mathrm{},x_{i_{r1}};i,s_1,\mathrm{},s_r)`$ (30)
$`\times {\displaystyle \underset{l=1}{\overset{r}{}}}\mu (s_l+u_l)\delta _{u_l}(x_{i_{l1}},\mathrm{},x_{i_{lu_l}}).`$
$`Z_{t+1}(y,x_1,\mathrm{},x_t;i)`$ is symmetric under permutations of $`x_1,\mathrm{},x_t`$ and lattice symmetries, e.g. simultaneous translation of $`y`$ and $`x_1,\mathrm{},x_t`$.
These formulae can be written graphically, according to the rules presented in Table 1. A sum over all dummy coordinates $`y`$ and all dummy valences $`i,f`$ and a sum over all inequivalent permutation of external coordinates $`x`$ or coordinate-valence pairs $`x,n`$ are understood. Notice that, despite the graphical notation, all pairs of roots of $`W_q^{\mathrm{nn}}`$ and $`W_q^{\mathrm{co}}`$ are equivalent.
Eq. (29) can be expressed by writing all polygons with a letter โcโ having 3 to $`q`$ vertices, placing a crossed dot with a positive integer label at each vertex in all inequivalent ways, the sum of the labels being $`q`$, and adding the tree contribution. The case $`q=5`$ is shown in Fig. 4.
$`G_q`$ and $`Z_q`$ can be computed by adding the contributions of all 2-irreducible $`r`$-rooted tree graphs with roots labelled by positive integers with sum $`q`$; for $`Z_q`$, the first root must be drawn as a square. The case $`Z_3`$ is shown in Fig. 5.
The next step is to write an expression of $`W_q^{\mathrm{co}}`$ in terms of $`W_r^{\mathrm{nn}}`$. For $`q=3`$ we have simply $`W_3^{\mathrm{co}}=W_3^{\mathrm{nn}}`$. The case $`q=4`$ is shown in Fig. 6. For larger values of $`q`$, the number of non-nodal contributions to $`W_q^{\mathrm{co}}`$ grows rapidly, and a systematic approach is needed.
Let us define for a (connected or non-connected) graph $`๐ข`$ and a vertex $`i`$ the graph $`๐ข/i`$: let $`๐ข^{}`$ be the connected component of $`๐ข`$ containing $`i`$; if $`๐ข^{}\backslash i`$ is connected, set $`๐ข/i=๐ข`$; otherwise, for each connected component $`๐ข_k^{}`$ of $`๐ข^{}\backslash i`$ generate the graph $`\overline{๐ข}_k`$ by adding a new vertex, internal or external like $`i`$, and joining it to all the vertices adjacent to $`i`$ in $`๐ข^{}`$, by the same number of edges; replace $`๐ข^{}`$ with the connected components $`\overline{๐ข}_k`$. The edges and vertices of $`๐ข/i`$ are in one-to-one correspondence with the edges and vertices of $`i`$, except for the new vertices which all correspond to $`i`$ (a nodal point of $`๐ข`$). Notice that $`๐ข/i/j=๐ข/j/i`$.
Let $`๐ข`$ be a compact 2-irreducible $`r`$-rooted graph ($`r3`$) with $`t`$ nodal points $`i_1,\mathrm{},i_t`$. Observe that all the connected components $`๐ข_l`$ of $`\overline{๐ข}๐ข/i_1/\mathrm{}/i_t`$ are non-nodal. Generate a new graph $`๐ฏ`$, the nodal skeleton of $`๐ข`$, in the following way: for each $`๐ข_l`$ with $`v3`$, containing $`n`$ roots corresponding to non-nodal roots of $`๐ข`$, write a root $`l`$ of $`๐ฏ`$ with a label $`n0`$; for each node $`i_k`$ of $`๐ข`$ write an unlabelled vertex $`i_k`$ of $`๐ฏ`$, internal or external like $`i_k`$; join $`i_k`$ to all the labelled roots $`l`$ such that $`๐ข_l`$ contain a vertex corresponding to $`i_k`$, and with all unlabelled vertices which are adjacent to $`i_k`$ in $`๐ข`$.
A nodal skeleton is a connected 1-irreducible tree graph, but it is not in general 2-irreducible (it may contain 2-valent internal vertices). A nodal skeleton enjoys the following properties: each 2-valent internal vertex is adjacent to a labelled root; labelled roots are never adjacent; unlabelled roots are at least 2-valent; $`m`$-valent roots with label $`n`$ satisfy $`n+m3`$. Every nodal skeleton can be generated by adding labels to some of the roots and by splicing 2-valent internal vertices into a 2-irreducible tree graph.
Every 1-irreducible tree graph, with some roots carrying a non-negative integer label, satisfying the above properties, is the nodal skeleton of a set of $`q`$-rooted 2-irreducible compact graphs, with $`q`$ equal to the sum of the labels plus the number of unlabelled roots. The contribution of this set to $`W_q^{\mathrm{co}}`$ can be computed by replacing each $`m`$-valent root of $`๐ฏ`$ carrying a label $`n`$ with an $`(n+m)`$-sided polygon whose vertices are the $`m`$ vertices adjacent to the root and $`n`$ new (unlabelled) roots, and applying the rules of Table 1.
An example of the construction of the nodal skeleton and its evaluation is presented in Fig. 7. The set of all nodal skeletons contributing to $`W_5^{\mathrm{co}}`$ is shown in Fig. 8; see also Fig. 6.
We could carry further the reduction of the set of graphs to be summed over, e.g. by identifying ladder graphs along the lines of Sect. 6. This is rather complicated for arbitrary $`r`$, and goes beyond the scope of the present work. Moreover, the zero-momentum projection described below is not applicable to the ladder graph reduction. Therefore we compute $`W_r^{\mathrm{nn}}`$ by restricting the sum of Eq. (19) to $`_{\mathrm{nn}}^{(r,2)}`$, the set of non-nodal 2-irreducible $`r`$-rooted graphs.
The above considerations can be simplified considerably if we are only interested in moments of the $`q`$-point functions, e.g.
$`\chi _q`$ $``$ $`{\displaystyle \underset{x_2,\mathrm{},x_q}{}}G_q(x_1,\mathrm{},x_q),`$
$`M_2^{(q)}`$ $``$ $`{\displaystyle \underset{x_2,\mathrm{},x_q}{}}(x_1x_2)^2G_q(x_1,\mathrm{},x_q).`$ (31)
Let us introduce the moments of $`Z_q`$, $`W_q^{\mathrm{co}}`$ and $`W_q^{\mathrm{nn}}`$:
$`\zeta _q(n)`$ $``$ $`{\displaystyle \underset{x_2,\mathrm{},x_q}{}}Z_q(x_1,\mathrm{},x_q;n),`$
$`\omega _q^{\mathrm{co}}(n_1,\mathrm{},n_q)`$ $``$ $`{\displaystyle \underset{x_2,\mathrm{},x_q}{}}W_q^{\mathrm{co}}(x_1,\mathrm{},x_q;n_1,\mathrm{},n_q),`$ (32)
the corresponding definition for $`\chi _q^{\mathrm{tr}}`$ and $`\omega _q^{\mathrm{nn}}(n_1,\mathrm{},n_q)`$, all independent second moments, etc.
Eq. (29) can be projected over zero momentum to give
$$\chi _q=\chi _q^{\mathrm{tr}}+\underset{r=3}{\overset{q}{}}\underset{\genfrac{}{}{0pt}{}{u_1\mathrm{}u_r}{u_1+\mathrm{}+u_r=q}}{}\frac{q!}{_{l=1}^ru_l!}\underset{i_1,\mathrm{},i_r}{}\omega _r^{\mathrm{co}}(i_1,\mathrm{},i_r)\underset{l=1}{\overset{r}{}}\zeta _{u_l+1}(i_l).$$
(33)
The computation of $`\zeta _q`$, $`\chi _q^{\mathrm{tr}}`$, and $`\omega _q^{\mathrm{co}}`$ is also easy; the graphical rules can be immediately projected over zero momentum, suppressing all coordinates $`x`$ and $`y`$ and removing all space delta functions; the only nontrivial part is the counting of the number of inequivalent permutations of roots.
$`\zeta _q(n)`$ can be computed by summing over all inequivalent 2-irreducible 1-ordered $`r`$-rooted tree graphs, with the roots labelled by positive integers $`u_1,\mathrm{},u_r`$ with sum $`q`$. The number of inequivalent permutations of the roots 2, โฆ, $`r`$ is
$$\frac{(q1)!}{S(u_11)!_{l=2}^ru_l!},$$
where $`S`$ is the symmetry factor of the labelled graph.
The computation of $`\chi _q^{\mathrm{tr}}(n)`$ is very similar, but we sum over unordered graphs, and the number of inequivalent permutations of the roots is
$$\frac{q!}{S_lu_l!}.$$
$`\omega _q^{\mathrm{co}}(n_1,\mathrm{},n_q)`$ can be computed by summing over all inequivalent unordered $`r`$-rooted nodal skeletons, with $`p`$ roots labelled by non-negative integers $`u_1,\mathrm{},u_p`$ with $`rp+u_1+\mathrm{}+u_p=q`$, with a weight $`1/S`$, where $`S`$ is the symmetry factor of the labelled graph, with unlabelled roots assigned an arbitrary distinct label (e.g. $`1`$).
Finally, we can compute $`\omega _q^{\mathrm{nn}}(n_1,\mathrm{},n_q)`$ by summing over unordered non-nodal 2-irreducible $`q`$-rooted graphs, provided that we use the modified symmetry factor and we symmetrize the result under permutation of the valences.
The computation of the second moment of the above quantities proceeds along the same lines. The factor $`(x_ix_j)^2`$ in Eq. (31) is dealt with in the following way: the two roots are connected by a chain of terms with a space structure of the form
$$\underset{x_j,y_1,\mathrm{},y_n}{}(x_ix_j)^2f_1(x_iy_1)f_2(y_1y_2)\mathrm{}f_{n+1}(y_nx_j)$$
(we have dropped the dependency on coordinates lying outside the branch connecting $`i`$ with $`j`$). Let us write
$$(x_ix_j)^2=(x_iy_1)^2+(y_1y_2)^2+\mathrm{}+(y_nx_j)^2+\text{cross terms}.$$
The cross terms do not contribute to the sum, and the result is
$$f_1^{(2)}f_2^{(0)}\mathrm{}f_{n+1}^{(0)}+f_1^{(0)}f_2^{(2)}\mathrm{}f_{n+1}^{(0)}+\mathrm{}+f_1^{(0)}f_2^{(0)}\mathrm{}f_{n+1}^{(2)},$$
where
$$f_i^{(0)}=\underset{y}{}f_i(y),f_i^{(2)}=\underset{y}{}y^2f_i(y).$$
Therefore we can compute the second moment by taking each contribution to the zero-momentum quantity, promoting one of the zero-momentum factors along the branch connecting the roots $`i`$ and $`j`$ to second moment, and summing over all possible choices.
By dealing with moments, we avoid the need of storing all the values of $`Z`$, $`W^{\mathrm{co}}`$, and $`W^{\mathrm{nn}}`$, which can rapidly exhaust all available memory. The extension to higher moments is straightforward but cumbersome.
## 8 Programming details
We wrote a set of computer programs to implement the automatic evaluation of the edge-renormalized LCE on the simple cubic lattice (sc) and on the body-centered cubic lattice (bcc) ($`3D`$); the same programs evaluate the LCE on two different representations of the square lattice ($`2D`$) and on the $`1D`$ lattice.
The computation of $`q`$-point functions is performed for a generic potential, keeping $`\mu _0(2n)`$ symbolic; each term of the series is a polynomial in $`\mu _0(2n)`$ with rational coefficients. We also implemented the same computation for a specific potential; this requires much less memory and is somewhat faster (up to 30%), but not enough to give up the flexibility of a generic potential.
To speed up search and insertion into ordered sets of data, graph sets and polynomials in $`\mu _0(2n)`$ are implemented as AVL trees (height-balanced binary trees) (cfr. e.g. Ref. , Chapt. 6.2.3), using the ubiqx library. Rational numbers and (potentially) large integers are handled by the GNU multiprecision (gmp) library.
Given the complexity of the procedure, it is crucial to perform a number of checks in order to flush out all algorithm and program errors. In $`1D`$ our series are compared with exact results for the spin-1/2 and the spin-1 Ising model; this is already a very stringent check, especially of the graph sets (cfr. Ref. ). In $`2D`$, our results are compared with the series for $`\chi `$ and $`M_2`$ for spin-1/2 published in Ref. . In $`3D`$, our results are compared with the lower-order series already available, for $`\chi `$ and $`M_2`$ for specific potentials in Refs. , and for $`\chi `$, $`M_2`$, and $`\chi _q`$ for a generic potential in Ref. . $`q(n)`$ can be computed from different combinations of $`n_1`$ and $`n_2`$ in Eq. (26); their agreement is non-trivial. Finally, for the spin-1/2 Ising model on any lattice, the series for $`G_q`$, rewritten in terms of $`v=\mathrm{tanh}K`$, must have integer coefficients.
### 8.1 Graph generation
A program generates the table of all unordered elementary 2-rooted 2-skeletons contributing to the desired order; the algorithm follows Ref. . Starting from the graph drawn in Fig. 9, we apply recursively the following modified Heap rules :
(a) join any two distinct vertices by a new edge, provided the two vertices are not already adjacent;
(b) insert a new internal vertex on any edge and join it to any vertex, excluding the edge extrema, by a new edge;
(c) insert two new internal vertices on any two distinct edges, and join them by a new edge;
(d) do not join the roots by a new edge.
The modifications to the original Heap rules (a), (b), and (c), marked in italics, prevent the generation of 2-reducible graphs. Rule (d) prevents the generation of ladder graphs.
The reduction of graphs to canonical form is performed using a generalization of the algorithm of Ref. . Graphs are stored in a compact form similar to the one of Ref. .
To reduce the proliferation of graphs at higher orders, it is extremely important to know the order (โstrict boundโ) $`o_s`$ at which a given graph will enter in the expansion (it is not trivially $`e`$, since we require even valence of all internal vertices, and, being interested in bipartite lattices, even length of all loops).
We must also keep in mind that some graphs donโt contribute at the desired order, but graphs generated from them might contribute. We define the โHeap boundโ $`o_H(๐ข)`$ as the minimum of $`o_s`$ on the set of graphs including $`๐ข`$ and all graphs generated from it. We also define the two bounds for $`(๐ข,\eta )`$, i.e. the minimal order when the vertex $`i`$ is forced to be embedded in a lattice site of parity $`\eta (i)`$.
We apply the modified Heap rules to the graph $`๐ข`$ in the following way: assign a parity label $`\eta (i)`$ to each vertex, in all the ways compatible with the Heap bound; apply the modified Heap rules assigning all possible parity labels to the new vertices; discard immediately the generated graph-parity pairs not satisfying the Heap bound; discard vertex parity information and store the generated graphs not isomorph to previously generated graphs. Finally, save into a file only the graphs satisfying the strict bound.
The generation of the elementary 2-rooted 2-skeletons contributing to the 25th order required 41 hours of computation on one CPU of a Compaq ES-40, and ca. 300 Mbytes of RAM. The number of inequivalent elementary 2-skeletons for each order of $`o_s`$ is reported in Table 2.
A similar program generates all unordered non-nodal 2-irreducible $`r`$-rooted graphs for $`r3`$. The table is initialized by applying rules (a) and (b) recursively, starting from each 2-irreducible $`r`$-rooted tree graph, until the result is non-nodal. Rules (a), (b), and (c) are then applied recursively. The number of inequivalent non-nodal 2-irreducible $`r`$-rooted graphs for each order of $`o_s`$ is reported in Table 2.
We remark that this graph tables can be used for the LCE of any spin model with $`\varphi \varphi `$ symmetry on a bipartite lattice in any dimension.
The generation of all the required families of tree graphs is straightforward.
### 8.2 Computation of the $`q`$-point functions
A separate program reads the table of elementary 2-rooted 2-skeletons and computes all components of $`W_2`$.
The evaluation of $`\mu `$ and of bond, nodal, and ladder contribution to $`W_2`$ is a straightforward application of the formulae of Sect. 6.
The evaluation of elementary contribution dominates the computation time, and must be optimized as much as possible. Assume that all lower-order contributions to $`W_2`$ have been computed. For each unordered elementary 2-skeleton $`๐ข`$ with $`o_s`$ not larger than the desired order, all inequivalent assignations $`๐ข,(n,l)`$ of edge valence parity $`n`$ and length parity $`l`$ compatible with the desired order, with even length of all loops, and with even valence of all internal vertices, are generated. We have implemented two different algorithms for the computation of the contribution of $`๐ข,(n,l)`$ to the two-point function.
The first algorithm is essentially the one used by Nickel and Rehr in Ref. : all inequivalent 1-irreducible 2-rooted graphs with a 2-skeleton compatible with $`๐ข,(n,l)`$ are generated, and their contributions are computed according to Eq. (11). In the second algorithm, the contribution is computed according to Eq. (25), and 1-irreducible graphs are not needed. The first algorithm is more efficient for 2-skeletons with large $`o_s`$, while the second algorithm is more efficient for small $`o_s`$; for each value of $`o_s`$ we select the algorithm which is (presumably) more efficient. On the sc lattice, the speed-up obtained over the use of either algorithm for all skeletons grows with the order, and is about a factor of 4 at order 23. On the bcc lattice, the first algorithm is very efficient, since the embedding number factorizes into a product of 1-dimensional embedding numbers ; we still use the second algorithm for the simplest 2-skeletons ($`o_s10`$), since the computation of the corresponding 1-irreducible 2-rooted graphs contributing to orders higher than 21 is extremely time- and memory-consuming.
Keeping in RAM all components of $`W_2`$ for a generic potential would be problematic. Most of these components are needed only to compute nodal and ladder contributions, and can be kept on disk; keeping in RAM just the components needed to compute elementary contribution is manageable.
The computation of the 25th-order LCE for the two-point function on the bcc lattice required ca. 400 hours of computation, and ca. 700 Mbytes of RAM.
A similar program reads the table of non-nodal 2-irreducible $`r`$-rooted graphs and the components of $`W_2`$, and computes $`\omega _q^{\mathrm{nn}}`$. The computation of $`\chi _q`$ is then straightforward. We computed $`\chi _4`$, $`\chi _6`$, and $`\chi _8`$ to 21st, 19th, and 17th order respectively on the bcc lattice.
The computation of the same quantities on the sc lattice is much slower (but does not requires more RAM); so far, we obtained $`W_2`$ to 23th order, with an effort not much smaller than the 25th order on the bcc lattice. The computation of $`W_2`$ and $`\chi _q`$ to the same orders as on the bcc lattice is in progress, but it will require a non-trivial amount of time.
## 9 Selected results
All high-temperature series computed in the present work are available for the most general potential, in the form of polynomials in the bare vertices $`\mu _0(2n)`$. The general results are extremely lengthy, and are only useful for further computer processing.
We present here a selection of high-temperature series for the spin-1/2 Ising model, i.e. for
$$f(\varphi )=\frac{1}{2}\left(\delta (\varphi +1)+\delta (\varphi 1)\right).$$
For sake of compactness, all the series are written in terms of $`v=\mathrm{tanh}K`$. Series for other specific potentials are available upon request from the author.
Although we computed all components of $`G_2(x,y)`$, we report here only $`\chi \chi _2`$ and $`M_2M_2^{(2)}`$ (cfr. Eq. (31)). For the $`q`$-point functions, we only computed $`\chi _q`$.
On the bcc lattice, we obtained
$`\chi `$ $`=`$ $`1+8v+56v^2+392v^3+2648v^4+17864v^5+118760v^6+789032v^7`$ (34)
$`+\mathrm{\hspace{0.17em}5201048}v^8+34268104v^9+224679864v^{10}+1472595144v^{11}+9619740648v^{12}`$
$`+\mathrm{\hspace{0.17em}62823141192}v^{13}+409297617672v^{14}+2665987056200v^{15}+17333875251192v^{16}`$
$`+\mathrm{\hspace{0.17em}112680746646856}v^{17}+731466943653464v^{18}+4747546469665832v^{19}`$
$`+\mathrm{\hspace{0.17em}30779106675700312}v^{20}+199518218638233896v^{21}+1292141318087690824v^{22}`$
$`+\mathrm{\hspace{0.17em}8367300424426139624}v^{23}+54141252229349325768v^{24}`$
$`+\mathrm{\hspace{0.17em}350288350314921653160}v^{25}+\mathrm{O}(v^{26});`$
$`M_2`$ $`=`$ $`8v+128v^2+1416v^3+13568v^4+119240v^5+992768v^6+7948840v^7`$ (35)
$`+\mathrm{\hspace{0.17em}61865216}v^8+470875848v^9+3521954816v^{10}+25965652936v^{11}`$
$`+\mathrm{\hspace{0.17em}189180221184}v^{12}+1364489291848v^{13}+9757802417152v^{14}`$
$`+\mathrm{\hspace{0.17em}69262083278152}v^{15}+488463065172736v^{16}+3425131086090312v^{17}`$
$`+\mathrm{\hspace{0.17em}23896020585393152}v^{18}+165958239005454632v^{19}+1147904794262960384v^{20}`$
$`+\mathrm{\hspace{0.17em}7910579661767454248}v^{21}+54332551216709931904v^{22}`$
$`+\mathrm{\hspace{0.17em}372033905161237212392}v^{23}+2540342425838560175616v^{24}`$
$`+\mathrm{\hspace{0.17em}17301457207110720278440}v^{25}+\mathrm{O}(v^{26});`$
$`\chi _4`$ $`=`$ $`2+64v+1168v^2+16576v^3+201232v^4+2204608v^5+22411504v^6`$ (36)
$`+\mathrm{\hspace{0.17em}215447872}v^7+1981980688v^8+17602809920v^9+151865668752v^{10}`$
$`+\mathrm{\hspace{0.17em}1278888344256}v^{11}+10550227820400v^{12}+85510907958720v^{13}`$
$`+\mathrm{\hspace{0.17em}682500568307184}v^{14}+5374496030148928v^{15}+41821018545214608v^{16}`$
$`+\mathrm{\hspace{0.17em}321992795063663936}v^{17}+2455641803116052752v^{18}+18567879503614668736v^{19}`$
$`+\mathrm{\hspace{0.17em}139310655514229882000}v^{20}+1037854026688655887552v^{21}+\mathrm{O}(v^{22});`$
$`\chi _6`$ $`=`$ $`16+1088v+36416v^2+853952v^3+15974528v^4+255491264v^5+3638767040v^6`$ (37)
$`+\mathrm{\hspace{0.17em}47395195712}v^7+574950589568v^8+6581949043264v^9+71803170318144v^{10}`$
$`+\mathrm{\hspace{0.17em}752047497945024}v^{11}+7606707093034368v^{12}+74649010982738112v^{13}`$
$`+\mathrm{\hspace{0.17em}713458387977120192}v^{14}+6661638582474716480v^{15}+60923519621981242752v^{16}`$
$`+\mathrm{\hspace{0.17em}546923327751320201536}v^{17}+4828463182433394315584v^{18}`$
$`+\mathrm{\hspace{0.17em}41987611565592990702272}v^{19}+\mathrm{O}(v^{20});`$
$`\chi _8`$ $`=`$ $`272+31744v+1673728v^2+58110976v^3+1538207872v^4+33584739328v^5`$
$`+\mathrm{\hspace{0.17em}634387677184}v^6+10699575811072v^7+164723097021568v^8`$
$`+\mathrm{\hspace{0.17em}2352360935459840}v^9+31540880634427392v^{10}+400802365468148736v^{11}`$
$`+\mathrm{\hspace{0.17em}4862781935250449280}v^{12}+56665753776838026240v^{13}`$
$`+\mathrm{\hspace{0.17em}637305912177206767104}v^{14}+6945658883867865975808v^{15}`$
$`+\mathrm{\hspace{0.17em}73600395257678784586368}v^{16}+760476823195422275111936v^{17}+\mathrm{O}(v^{18}).`$
On the sc lattice, we obtained
$`\chi `$ $`=`$ $`1+6v+30v^2+150v^3+726v^4+3510v^5+16710v^6+79494v^7+375174v^8`$ (39)
$`+\mathrm{\hspace{0.17em}1769686}v^9+8306862v^{10}+38975286v^{11}+182265822v^{12}+852063558v^{13}`$
$`+\mathrm{\hspace{0.17em}3973784886}v^{14}+18527532310v^{15}+86228667894v^{16}+401225368086v^{17}`$
$`+\mathrm{\hspace{0.17em}1864308847838}v^{18}+8660961643254v^{19}+40190947325670v^{20}`$
$`+\mathrm{\hspace{0.17em}186475398518726}v^{21}+864404776466406v^{22}+4006394107568934v^{23}`$
$`+\mathrm{O}(v^{24});`$
$`M_2`$ $`=`$ $`6v+72v^2+582v^3+4032v^4+25542v^5+153000v^6+880422v^7+4920576v^8`$ (40)
$`+\mathrm{\hspace{0.17em}26879670}v^9+144230088v^{10}+762587910v^{11}+3983525952v^{12}+20595680694v^{13}`$
$`+\mathrm{\hspace{0.17em}105558845736}v^{14}+536926539990v^{15}+2713148048256v^{16}+13630071574614v^{17}`$
$`+\mathrm{\hspace{0.17em}68121779384520}v^{18}+338895833104998v^{19}+1678998083744448v^{20}`$
$`+\mathrm{\hspace{0.17em}8287136476787862}v^{21}+40764741656730408v^{22}+199901334823355526v^{23}`$
$`+\mathrm{O}(v^{24}).`$
### Acknowledgments
We would like to thank Michele Caselle and Ettore Vicari for many useful discussions. Special thanks are due to Paolo Rossi, who worked out exact results for the $`1D`$ spin-1 Ising model, and to Andrea Pelissetto, for his critical reading of the manuscript.
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# Cosmological Reionization Around the First Stars: Monte Carlo Radiative Transfer
## 1 Introduction
An increasing number of works has been recently dedicated to the study of the reionization of the universe (Gnedin & Ostriker 1997; Haiman & Loeb 1998; Valageas & Silk 1999; Ciardi et al. 2000, CFGJ; Miralda-Escudรฉ, Haehnelt & Rees 2000; Chiu & Ostriker 2000; Bruscoli et al. 2000; Gnedin 2000; Benson et al. 2000) which use both analytical and numerical approaches. An important refinement has been introduced by the proper treatment of a number of feedback effects ranging from the mechanical energy injection to the H<sub>2</sub> photodissociating radiation produced by massive stars (CFGJ). However, probably the major ingredient still lacking for a physically complete description of the reionization process is the correct treatment of the transfer of ionizing photons from their production site into the intergalactic medium (IGM). Attempts based on different approximated techniques have been proposed (Razoumov & Scott 1999; Abel, Norman & Madau 1999; Norman, Paschos & Abel 1998; Gnedin 2000; Umemura, Nakamoto & Susa 1999), which sometimes are not readily implemented in cosmological simulations. Therefore, it is crucial to develop exact and fast methods that can eventually yield a better treatment of the propagation of ionization fronts in the early universe. The first three papers are based on the ray-tracing method, whereas Gnedin (2000) uses the so called local optical depth approximation. Finally Umemura, Nakamoto & Susa (1999) implemented a time-independent ray-tracing method. Monte Carlo (MC) methods have been widely used in several physical/astrophysical areas to tackle radiative transfer problems (for a reference book see Cashwell & Everett 1959) and they have been shown to result in fast and accurate schemes. Here we build up on previous experience of our group (Bianchi, Ferrara & Giovanardi 1996; Ferrara et al. 1996; Ferrara et al. 1999; Bianchi et al. 2000) in dealing with MC problems to present a case study of cosmological H<sub>II</sub> regions produced by the first stellar sources. The study is intended as a test of the applicability of the adopted techniques in conjunction with cosmological simulations in terms of convergency, accuracy, and speed of the scheme.
As an example of interesting cosmological application, we study the reionization produced by a stellar source of total mass $`M=2\times 10^8M_{}`$ turning on at $`z12`$, located at a node of the cosmic web. The study includes a Spectral Energy Distribution of a zero-metallicity stellar population, and two Initial Mass Functions (Salpeter/Larson); the IGM spatial density distribution is deduced from high resolution cosmological simulations described below.
In a forthcoming paper we will then improve the results of CFGJ by exploiting the strength of the MC method to release the approximations relative to the radiative transfer made in that paper.
## 2 Numerical Simulations
We have studied the evolution of small scale cosmic structure in a standard $`\mathrm{\Lambda }`$+CDM model with $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ and $`\mathrm{\Omega }_m=0.3`$. The baryon contribution to $`\mathrm{\Omega }_m`$ is $`\mathrm{\Omega }_b=0.04`$, the adopted value of the Hubble constant is $`H_0=70`$ km sec<sup>-1</sup> Mpc<sup>-1</sup>.
The initial conditions are produced with the COSMICS <sup>1</sup><sup>1</sup>1Bertschinger, http://arcturus.mit.edu/cosmics/ code adopting a normalization $`\sigma _8=1.14`$ for the power spectrum. The initial particle distribution (at redshift $`z=103.2`$) is made out of a $`256^3`$ grid in a box of $`4`$ Mpc (comoving), thus giving a mass per particle of about $`1.5\times 10^5M_{}`$. We adopt vacuum boundary conditions; this requires the initial box to be reshaped. Assuming that tidal forces from large scale fluctuations are negligible for this problem, we cut a sphere of radius $`2`$ Mpc from the initial cube and follow the evolution of the system in isolation. In practice, we study the evolution of a central sphere of radius about $`1.2`$ Mpc in full mass resolution with a coarse resampling of the most external particles. In the hi-resolution (central) region the number of particles is 2313847, whereas the rebinning left us with 10008 particles of variable mass (increasing with distance from the center) in the external $`1.22.0`$ Mpc shell.
We evolve this initial matter distribution with GADGET (Springel, Yoshida & White 2000), a parallel N-body tree-code with an SPH scheme to describe hydrodynamical processes. SPH particles are placed on top of hi-resolution dark matter particles with the same velocity and a low initial temperature (about $`150`$ K, the temperature of intergalactic medium as determined by adiabatic expansion of the universe after decoupling); the mass of each fluid particle is rescaled according to the baryon and matter fraction quoted above. The (spline) softening length is $`1`$ kpc for the gas particles, and $`3`$ kpc for hi-resolution dark matter particles, respectively. For the problem at study here, the simulation is stopped at $`z5`$, thus avoiding the very strong nonlinear phase in which our assumption of zero tidal effects from large scale structure would break down.
The gas-dynamics evolution is purely adiabatic in order to reasonably limit the integration time. Thus, we are not able to follow the details of cooling processes that ultimately lead to the formation of very dense and cold star forming regions. However, the density field of the mildy overdense IGM should be reasonably well described by this assumption. Our mass resolution does not also allow a full resolution of the Jeans mass, and therefore the presence of sub-grid structure that we cannot track is expected. This is a problem common to all presently available simulations, implying that the treatment of IGM clumping is only approximate, as recently pointed out by Haiman, Abel & Madau (2000).
Simulations were performed on the Cray-T3E at Rechenzentrum-Garching, the Joint Computing Center of the Max-Planck-Gesellschaft and Max-Planck-Institut fuer Plasmaphysik.
Finally, from the SPH particle temperature and density distribution we reconstruct a central cubic region of $`128^3`$ cells, the grid values of any field being evaluated from particle values with an SPH-like interpolation. For the reionization study presented below, we have concentrated, as an example, on a source located at redshift $`z12`$ immersed in the density field calculated from the simulations described above. The location and mass of the dark matter halos in the simulation have been determined by means of a friend-of-friend algorithm. We have then selected a halo suitably located in the node of a cosmic filamentary structure and associated a luminous source to it. The total halo mass is $`M=2\times 10^8M_{}`$; the corresponding stellar mass has been determined following the same prescriptions of CFGJ, i.e.
$$M_{}=(\mathrm{\Omega }_b/\mathrm{\Omega }_m)Mf_bf_{},$$
(1)
where $`f_b=0.08`$ is the fraction of virialized baryons able to cool and become available to form stars, and $`f_{}=0.15`$ is the star formation efficiency. Although with our cell size of 15 kpc we are marginally able to resolve the source (comoving) virial radius (about 20 kpc), we do not attempt to treat the the escape of ionizing photons (Dove, Shull & Ferrara 2000; Wood & Loeb 2000) from the galaxy interstellar medium. Instead we take a fixed value for this quantity equal to $`f_{esc}=0.2`$, the upper limit obtained by the above studies, by which we multiply the source ionizing photon flux.
## 3 PopIII Star Spectrum and IMF
The SEDs of metal free star systems used in this paper were calculated implementing the synthetic evolutionary code developed for a Simple Stellar Population (Brocato et al. 1999, Brocato et al. 2000). For a detailed description of the computational procedure we refer to those papers; here we just outline that the results of this code have been largely tested on young LMC clusters and old galactic globular clusters showing a very good agreement. The present synthesis models rely on the set of homogeneous evolutionary computations for stars with a cosmological amount of heavy elements and helium (Z=10<sup>-10</sup> and Y=0.23) presented by Cassisi & Castellani (1993) and Cassisi, Castellani & Tornambรจ (1996) and properly extended to higher masses in order to cover a wide range of ages. All the major evolutionary phases are taken into account, both the H and the He burning phase for stars with original masses in the range 0.7$``$40 M, for which the He burning phase is followed up to the Carbon ignition or, alternatively, to the onset of the thermal pulses. Models atmospheres are by Kurucz (1993) for metallicity Z=2$`\times 10^7`$. The available temperatures range from 3700 K to 50000 K, so the spectra at higher temperatures and gravity are simply extrapolated in order to keep the homogeneity of the grid. The wavelength range is from 91 ร
to 160 $`\mu `$m and the spectral resolution is typically 10-20 ร
in the UV to optical bands. We consider an instantaneous burst of star formation assuming that masses range from 1 $`M_{}`$, as recently suggested for Pop III stars by Nakamura & Umemura (1999a, 1999b), to 40 $`M_{}`$ and two cases for the IMF: $`i)`$ a standard Salpeter law, and $`ii)`$ a Salpeter function at the upper mass end which falls off exponentially below a characteristic stellar mass $`M_c`$ (Larson 1998). The last case is taken into account since both theory of thermo-dynamical condition of primordial gas and the absence of observations of solar mass, metal free stars seem to suggest at birth a paucity of low mass stars. Moreover, the mass interval used in this paper properly describes population models of age down to 5 Myr, the lifetime of a 40 M star; therefore, the burst models of this age are actually populated by 1-40 M stars. The resulting grid of SEDs were calculated at intervals of 1, 10 and 50 Myr between 1 and 20 Myr, 20 and 100 Myr, and 100 and 150 Myr, respectively. Other calculations of zero-metallicity star SEDs are present in the literature (see e.g. Cojazzi et al. 2000; Tumlinson & Shull 2000), but they do not show the detailed SED time evolution.
In Fig. 1 we show the adopted SEDs for the Salpeter IMF and the Larson IMF with $`M_c=5M_{}`$. The two sets of curves do not differ sensibly in their spectral shape evolution but the Larson IMF has about 4 times larger specific photon flux at the Lyman continuum. In Fig. 2 we show the time evolution of the total number of ionizing photons (in units of $`10^{60}`$) emitted by a solar mass of metal free stars with the two IMFs. After about $`10^7`$ yr both curves flatten as a consequence of the decreasing number of surviving massive stars.
It is important to note that, during the first 5 Myr, a zero-metallicity stellar cluster with a Salpeter IMF has an ionizing photon rate about 25 times higher than the analogous one for a 1% solar metallicity used by CFGJ. As we will see later, this has important effects on the reionization process.
## 4 Monte Carlo Radiative Transfer
The application of MC schemes to radiative transfer problems requires that the radiation field is discretized in a representative number of photon packets, $`๐ฉ_p`$. The processes involved (e.g. packet emission and absorption) are then treated statistically by randomly sampling the appropriate distribution function. Let us consider a source of bolometric luminosity $`L`$ and lifetime $`t_s`$ ($`L`$ can be a function of time as our method easily allows to treat short-lived and variable sources). Packets have the same energy, $`_p`$, but contain a different number of monochromatic photons, $`N_\gamma `$, of frequency $`\nu `$ so that for the $`j`$-th packet it is $`_p=N_\gamma (j)h\nu (j)`$; their rate is $`\dot{๐ฉ}_p=L/_p`$ and the simulation time when the $`j`$-th photon packet is emitted is then $`t=jdt`$ ($`j=1,..,๐ฉ_p`$), where $`dt=\dot{๐ฉ}_p^1`$. The number of photons in each packet emitted by the source during $`t_s`$ is:
$$N_\gamma (j)=\frac{E_s}{๐ฉ_ph\nu (j)}=\frac{Lt_s}{๐ฉ_ph\nu (j)}.$$
(2)
The packet frequency, $`\nu `$, is obtained by sampling the source SED, $`S(\nu )`$: if $`R_\nu `$ is a random number, $`\nu `$ is given implicitly by:
$$R_\nu =_{\nu _H}^\nu ๐\nu ^{}S(\nu ^{})\left[_{\nu _H}^{\nu _{max}}๐\nu ^{}S(\nu ^{})\right]^1,$$
(3)
where $`\nu _H`$ is the threshold for hydrogen ionization.
We adopt spherical coordinates ($`r,\theta ,\varphi `$) with origin at the source location ($`x_s,y_s,z_s`$) in the box, such that the cartesian coordinates are:
$`x`$ $`=`$ $`r\mathrm{sin}\theta \mathrm{cos}\varphi 0\varphi 2\pi ,`$
$`y`$ $`=`$ $`r\mathrm{sin}\theta \mathrm{sin}\varphi 0\theta \pi ,`$ (4)
$`z`$ $`=`$ $`r\mathrm{cos}\theta .`$
Assuming that the emission is isotropic and spherically symmetric, photon packets will propagate along the direction:
$$\widehat{๐ฌ}[\widehat{r},\mathrm{arccos}(12R_\theta )\widehat{\theta },2\pi R_\varphi \widehat{\varphi }],$$
(5)
where $`R_\theta `$ and $`R_\varphi `$ are random numbers. Note that the above choice automatically ensures a constant surface density of photon packets, thus preventing flux anisotropies at the poles due to angle sampling; this can be seen as follows. As we are sampling $`R_\theta `$ from a uniform distribution, the fraction of photon packets emitted in the angular range $`\theta _1<\theta <\theta _2`$ is equal to
$$d๐ฉ_p(\theta _1,\theta _2)=(R_{\theta _2}R_{\theta _1}),$$
(6)
where $`R_{\theta _i}=(1/2)(1\mathrm{cos}\theta _i)`$ from eq. 5. As the area fraction of the sphere belt delimited by $`\theta _1`$ and $`\theta _2`$ is
$$dA(\theta _1,\theta _2)=\frac{1}{2}(\mathrm{cos}\theta _1\mathrm{cos}\theta _2)$$
(7)
one sees that the surface density $`d๐ฉ_p/dA`$ of photon packets is constant and independent of angles.
Once emitted, a packet of frequency $`\nu `$ will travel for a finite path length of optical depth:
$$\tau _i(\nu )=\sigma (\nu )\underset{k=1}{\overset{i}{}}(N_{\mathrm{HI}})_k=\sigma (\nu )\underset{k=1}{\overset{i}{}}(n_{\mathrm{HI}})_kf\mathrm{\Delta }x,$$
(8)
before being absorbed in the $`i`$-th grid cell. Here $`\sigma (\nu )`$ is the photoionization cross section, $`(N_{\mathrm{HI}})_k=(n_{\mathrm{HI}})_kf\mathrm{\Delta }x`$ is the IGM neutral hydrogen column density through the $`k`$-th cell of size $`\mathrm{\Delta }x`$; the cell with $`k=0`$ contains the source location and it is not included in the opacity count, as its contribution is already included in $`f_{esc}`$. For simplicity, the results presented here only include H opacity; inclusion of He, molecules (as H<sub>2</sub> and HD) and heavy elements, if present, is straightforward, although computationally expensive for molecules as one has to properly compute both the line widths and their Doppler shifts. The factor $`f`$ accounts for the fact that the path through a cell is in the range $`0<\mathrm{}\sqrt{3}\mathrm{\Delta }x`$ depending on its inclination.
To limit the computational cost, we do not attempt to calculate the actual path and we treat this effect statistically as follows. Given a cubic box of unit size $`\mathrm{\Delta }x=1`$, via an independent Monte Carlo procedure, we randomly pick up on a face the coordinates of an entering ray and the direction of its propagation. We then derive the differential probability distribution function, $`P(\mathrm{})`$, of the path lengths (Fig. 3) and define $`f=0.56`$ as the median value of such distribution.
The value of $`i`$ in eq. 8 is defined as the minimum one for which the following condition is satisfied:
$$\tau _i(\nu )\tau =\mathrm{ln}(1R_\tau ),$$
(9)
with $`R_\tau `$ a random number. The previous relation is obtained by sampling the photon absorption probability distribution.
In each cell $`k=1,..,i`$ along the path from the source to the absorption site of the $`j`$-th packet, we update the value of the hydrogen ionization fraction; $`x=n_{\mathrm{H}^+}/n_\mathrm{H}`$, $`n_{\mathrm{H}^+}`$ and $`n_\mathrm{H}`$ are the ionized and total hydrogen number density, respectively. This is done by advancing the time-dependent ionization equation in discretized form. At time $`t=t^{n+1}`$ it is:
$$x_k^{n+1}=x_k^n+\mathrm{\Delta }t[\gamma _c(x_k^n)\gamma _r(x_k^n)]+\frac{N_\gamma (j)}{(n_\mathrm{H})_k(\mathrm{\Delta }x)^3}\delta _{ki};$$
(10)
$`\gamma _c`$ and $`\gamma _r`$ are the collisional ionization and recombination rates, respectively (Cen 1992); they are evaluated at the temperature of the gas in the cell. The initial temperature is increased to $`T=10^4`$ K after absorption of a photon packet and set equal to $`T=2.734(1+z)`$, i.e. the CMB temperature, after the cell recombines (see below). This is a reasonable approximation, although $`T`$ could be smaller at higher redshift (Ciardi, Ferrara & Abel 2000). The last term, describing photoionization due to photon packet absorption, is present only for the cell in which absorption takes place ($`k=i,\delta _{ki}=1`$); in the other cells along the path ($`k<i,\delta _{ki}=0`$) we simply follow the gas recombination. In case the derived value of $`x_k^{n+1}>1`$, we set $`x_k^{n+1}=1`$, and we use the extra photons to ionize the cell $`x_{k+1}`$. In practice, we write $`N_\gamma (j)=N_\gamma ^{(1)}(j)+N_\gamma ^{}(j)`$, where the number of photons required to obtain $`x_k^{n+1}=1`$ is:
$$N_\gamma ^{(1)}(j)=\{(1x_k^n)\mathrm{\Delta }t[\gamma _c(x_k^n)\gamma _r(x_k^n)]\}(n_\mathrm{H})_k(\mathrm{\Delta }x)^3.$$
(11)
The most delicate point of the ionization fraction evaluation is to find an appropriate expression for $`\mathrm{\Delta }t`$. In principle, one could update the ionization state on the entire grid after each photon packet emission. However, this is computationally too expensive. For this reason we choose to update the ionization state of cells along the packet path and we deduce $`\mathrm{\Delta }t`$ as follows.
The cell in which the source is located is enclosed by a cubic surface (shell). If we define the order of this shell to be $`g=1`$, the number of its member cells is given by:
$$N_s(g)=(2g+1)^3(2g1)^3=2(12g^2+1)g1,$$
(12)
or $`N_s(1)=26`$. The second order shell ($`g=2`$) is instead made of 98 cells and so on. As the photon packet directions are isotropically distributed, statistically a cell in a shell of order $`g`$ will be along a path every $`[N_s(g)]^1`$ packets. Therefore the typical time scale between two subsequent ionization fraction updates is $`\mathrm{\Delta }t=N_s(g)dt`$; this is the value we use in eq. 10 above. The order $`g`$ is calculated as:
$$g=\mathrm{๐๐๐ฅ}[k_x,k_y,k_z],$$
(13)
where the $`k_i`$โs indicate the discrete cartesian coordinates of the cell in the reference frame eq. 4. The above approximation sets a lower limit on $`๐ฉ_p`$, as in order to be valid it is necessary that the update time $`\mathrm{\Delta }t`$ is shorter than the average recombination time, $`t_r`$, in the box. Hence to advance the simulation up to a time $`t_s`$ the number of packets required is:
$$๐ฉ_p\stackrel{>}{}24N^2\frac{t_s}{t_r},$$
(14)
where $`N^3`$ is the total numer of cells in the box. This detailed treatment of recombination is important to accurately evaluate the optical depth through eq. 8.
We then introduce an ionization vector (size $`N^3`$) whose elements are initially set to zero. The vector element corresponding to a cell for which the ionization fraction is increased above $`x_{th}=0.99`$ as a result of an absorption event is updated to contain the value:
$$j_r=\mathrm{๐ผ๐๐ก}\left[\frac{t_r(x_k,n_k)+t}{dt}\right],$$
(15)
where $`t_r(x_k,n_k)`$ is the local recombination time and $`t`$ is the current simulation time. If further absorption does not take place before the time $`j_rdt`$, we allow for recombination emission from the given cell; if instead another absorption event takes place, we re-update the vector element to a new value of $`j_r`$. As already mentioned, after recombination, the cell temperature is set at the CMB temperature at that redshift.
Following a cell recombination a packet is emitted along a random direction from the same location. The probability that the emitted packet has $`h\nu (j)>13.6`$ eV is given by the ratio between the number of recombinations to the ground level and the total one, i.e. 50% for $`T=10^4`$ K. The frequency of the re-emitted packet is determined through eq. 3, where the source SED is substituted by the volume emissivity of the diffuse radiation, given by the Milne relation:
$`S_d(\nu )`$ $`=`$ $`{\displaystyle \frac{2h\nu ^3}{c^2}}{\displaystyle \frac{g_i}{g_{i+1}}}\left({\displaystyle \frac{h^2}{2\pi m_pkT_e}}\right)^{3/2}\times `$ (16)
$`\sigma (\nu )\mathrm{e}^{h(\nu \nu _H)/(kT_e)}n_{H^+}n_e,`$
where $`g_i`$ and $`g_{i+1}`$ are the hydrogen ground level statistical weights; $`T_e`$ and $`n_e`$ are the electron temperature and number density. The propagation and absorption of re-emitted packets is then followed in the same manner as for primary ones. The only difference is that the opacity contribution of the re-emitting cell is now included.
The method is easily generalized to include an arbitrary number of point sources once the packet emission rates, and hence their emission sequence for the various sources are given.
### 4.1 Testing the Method
The method described in the previous Section has been tested against the analytical solution for the temporal evolution of the radius of an H<sub>II</sub> region produced by a source with constant ionizing rate and expanding in an homogeneous medium (Donhaue & Shull 1987; Shapiro & Giroux 1987). The H<sub>II</sub> region equivalent radius, $`R_n`$, is numerically derived by first assuming that a cell $`k`$ is in the ionized state when $`x_k>x_{th}=0.99`$. If $`N_{ion}`$ is the number of ionized cells in the box, the volume occupied by the ionized region is $`V=\mathrm{\Delta }x^3N_{ion}`$ and the equivalent radius is then $`R_n=(3V/4\pi )^{1/3}`$. In all the tests presented here the gas density is set equal to average IGM one in an Einstein-de Sitter universe with $`\mathrm{\Omega }_bh^2=0.019`$ at $`z=11.6`$, i.e. $`n_\mathrm{H}=4.5\times 10^4`$ cm<sup>-3</sup>; the adopted spectrum is monochromatic, with a photon energy equal to 13.6 eV.
First, we check for convergency of our scheme. We run three simulations with different values of $`๐ฉ_p=(5,50,500)\times 10^5`$. The source ionizing photon rate is $`\dot{๐ฉ}_\gamma =2.5\times 10^{50}`$ s<sup>-1</sup>. For each run we calculate the residuals with respect to the highest resolution run $`๐ฉ_p=500\times 10^5`$; they are defined as:
$$D^y=\frac{R_n^yR_n^{500}}{R_n^{500}},$$
(17)
where the superscript $`y=5,50`$ refers to the low- and medium-resolution runs, respectively. The residuals are shown in Fig. 4 at different evolutionary times together with the associated mean quadratic errors; the error on $`R_n`$ is assumed to be equal to $`\mathrm{\Delta }x`$. Satisfactory convergency is achieved by the medium-res run; subsequent tests have hence been done using $`๐ฉ_p=50\times 10^5`$. For reference, this run takes only $`300`$ s (CPU time) on a SUN ULTRA10/333MHz workstation.
Next, we assess the accuracy of the solution. In Fig. 5a we show the time evolution of $`R_n`$, as produced by sources of different ionizing photon rates $`\dot{๐ฉ}_\gamma =(2.5,12.5,25,125)\times 10^{50}`$ s<sup>-1</sup>, compared to the corresponding analytical radius, $`R_a`$. Fig. 5b shows the residuals (defined analogously to eq. 17 with $`R_n^{500}`$ substituted by $`R_a`$) with respect to the analytical solution. In the worst case (very faint source, early evolutionary times) the residual is about 20%; This relatively large discrepancy occurs because the size of the ionized region is only a few cells. At later times and/or for more luminous sources the agreement is remarkably good, i.e. within a few percent.
Finally, we check the effects of varying the spatial resolution. To this aim we have considered runs with different box sizes, keeping a fixed linear number of cells $`N=128`$. In Fig. 6 we show the residuals for runs with a box size $`(1,0.2,0.1)\times N\mathrm{\Delta }x`$ with respect to the highest spatial resolution run for a source ionizing photon rate $`\dot{๐ฉ}_\gamma =2.5\times 10^{50}`$ s<sup>-1</sup>.
## 5 Results and Discussion
The combined cosmological and radiative transfer simulations described above allow us to determine the evolution of the ionized region around the selected zero-metallicity stellar object. Illustrative slices extracted from our simulation box through the source location are shown in Fig. 7 at different times (10, 20, 40, 60, 100 Myr) after the source has been turned on and a for Larson IMF, together with the initial density field. The above time interval corresponds to less than two redshift units at $`z12`$, and it is therefore not unreasonable for the purposes of this paper to neglect the IGM density evolution As seen from Fig. 7, the ionization front (I-front) breaks out from the galaxy very rapidly, leaving behind a very clumpy ionization structure in the immediate surroundings where the IGM is overdense and, consequently, recombination times are shorter. As an aside, it is interesting to note that, as the circular velocity of the parent galaxy is a few km s<sup>-1</sup>, the ionized gas (whose sound speed is of order of 10 km s<sup>-1</sup>) will very likely be able to leave the galaxy, thus quenching further star formation. If this is a widespread phenomenon, it might have strong implications for the evolution of dwarf galaxies as already outlined by some authors (Barkana & Loeb 1999; Ferrara & Tolstoy 2000). Once the I-front expands in the IGM, it is channeled โ for similar reasons โ by the large filaments into the underdense volumes (voids) and its structure becomes very complex and jagged, assuming, in a 2D representation, a characteristic butterfly shape. The typical overdensity of the clumps encountered by the expanding I-front is $`30`$. Qualitatively, the typical final extent of the ionized region is of about 1 comoving Mpc; this size is reached already after 60 Myr. After that time, the rapid decrease of the source ionization power (Fig. 2) slows down the expansion. The final stage of the evolution is constituted by a relic H<sub>II</sub> region which slowly starts to recombine on a time scale which in the voids can be as long as 0.8 Gyr. The MC technique shows here all his power in following the details of the I-front evolution. For example, it tracks remarkably well the channeling induced by the large scale structure: the ionization front cannot propagate inside the densest filaments due to their large optical depth and short recombination time. In addition, also the effects of shadowing produced by isolated clumps are clearly recognized from the fingers protruding from the H<sub>II</sub> region to the left of the source in the upper panel of Fig. 8. The ionization cone visible below the source, caused by a large overdensity located close to the source in that direction, is also a result of shadowing. Thus it seems that our scheme is highly suitable at least for this type of cosmological radiative transfer calculations.
We now turn to the differences induced by our two assumptions concerning the IMF. Fig. 8 shows a comparison between the final stages (100 Myr after the source turn on) of the I-front evolution for a Larson (upper panel) and a Salpeter (bottom) IMF. The source stellar mass and other properties of the simulations are the same as those discussed above. For a Salpeter IMF, the volume of the ionized region is smaller by a factor 8, although the shape is very similar to the one for a Larson IMF case at an earlier stage, roughly corresponding to 10 Myr (see Fig. 7). This was expected from the two adopted SEDs. In fact, the total number of ionizing photons (see Fig. 2) per stellar mass formed integrated over the entire source lifetime and spectral extent is $`5\times 10^{61}`$ ($`10^{61}`$) for the Larson (Salpeter) IMF. Thus, the IMF might play an important role for the reionization of the universe; in addition, zero-metallicity stars have larger ionizing power as already stressed previously and recently addressed by other authors (Cojazzi et al. 2000; Tumlinson & Shull 2000).
To assess the differences between the evolution of an I-front propagating in an inhomogeneous medium, as in the simulations discussed here, and in an homogeneous gas at the mean IGM density, we have calculated the ratio $`๐ฑ=V/V_h`$ between the ionized volumes in our simulation and in the homogeneous case. This is found to be roughly independent of time and equal to $`๐ฑ=0.45`$ for the Larson IMF. In the Salpeter IMF case $`๐ฑ`$ varies slightly with time from 0.125 to 0.2. Hence, this result confirms that the ionized volume tends to be smaller in the inhomogeneous case, but suggests that the effect is not dramatic.
## Acknowledgments
We would like to thank S. Cassisi for providing us with the stellar tracks; N. Yoshida and V. Springel for advises in various numerical problems; the referee T. Abel for useful comments. This research was supported in part by the National Science Foundation under Grant No. PHY94-07194 (SM) and by the Italian Ministery of University, Scientific Research and Technology (MURST) (GR).
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# The Enrichment History of the Intergalactic Medium - Measuring the C IV/H I Ratio in the Ly๐ผ ForestThe data presented herein were obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among the California Institute of Technology, the University of California and the National Aeronautics and Space Administration. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. .
## 1. Introduction
Our understanding of the intergalactic medium (IGM) has undergone a paradigm shift in recent years, largely as a result of powerful hydrodynamical simulations. When the Ly$`\alpha `$ forest was first observed in the late 1960s, the rich field of H I absorption blueward of the QSOโs Ly$`\alpha `$ emission was interpreted as the detection of discrete intergalactic clouds (e.g. Lynds & Stockton 1966; Lynds 1971; Sargent et al. 1980). In order to understand the existence of such isolated absorbers, various theories of cloud confinement, including self-gravity (Melott 1980), cold dark matter minihaloes (Rees 1986) and the presence of an inter-cloud medium (e.g. Sargent et al. 1980; Ostriker & Ikeuchi 1983) were proposed. However, the advent of hydrodynamical simulations, which model the growth of structure in the high redshift universe, provided a significant revision to our picture of the IGM (see the recent review by Efstathiou, Schaye and Theuns 2000). It has been found that in the presence of a UV ionizing background, the โbottom-upโ hierarchy of structure formation knitted a complex, but smoothly fluctuating โcosmic webโ in the IGM (e.g. Cen et al. 1994; Hernquist et al. 1996; Bi & Davidsen 1997). The absorption in the Ly$`\alpha `$ forest is caused not by individual, confined clouds, but by a gradually varying density field characterized by overdense sheets and filaments and extensive, underdense voids. The advance in theoretical simulations has been matched by increasingly high quality data, as comprehensively reviewed by Rauch (1998). One of the major discoveries concerning the IGM has been the identification of metal absorption lines associated with many of the Ly$`\alpha `$ forest clouds (Cowie et al. 1995; Tytler et al. 1995). Thus, whilst the Ly$`\alpha `$ forest was once thought to be chemically pristine, it has now been well-established that a large fraction of the high column density Ly$`\alpha `$ clouds ($`N`$(H I) $`>`$ 14.5), associated with collapsing, over-dense structures, contain metals (most notably C IV) โ the signature of enrichment by the products of stellar nucleosynthesis (Songaila & Cowie 1996).
The presence of metals in the Ly$`\alpha `$ forest may be reasonably explained either by in-situ enrichment (local star formation in the H I cloud itself or in a nearby galaxy) or by early pre-enrichment by a high redshift episode of Population III stars. Whilst the effects of supernova feedback are still not fully understood and therefore the spatial extent of wind-driven ejecta is poorly constrained, enrichment by galactic winds and superbubbles is unlikely to be an efficient way to distribute metals over distances large in comparison with the mean separation between galaxies (MacLow & Ferrara 1999). Instead, models have appealed to larger scale processes such as merging and turbulent diffusion as the dominant mixing mechanisms (Gnedin & Ostriker 1997; Gnedin 1998; Ferrara, Pettini & Shchekinov 2000). Such processes would take time to smooth out the metallicity of the IGM so that at $`z3`$ the metal enrichment is still expected to be very patchy. Whilst the deep potential wells of galaxies inhibit efficient, widespread distribution of metals far from their sites of formations, small regions of star formation at high redshift may be able to eject their nucleosynthetic products for more homogeneous mixing (e.g. Nath & Trentham 1997 and references therein). An episode of Population III star formation may therefore have spread metals far from their sites of formation, seeding โsterileโ regions of the IGM with metals (Ostriker & Gnedin 1996). Clearly, distinguishing between in-situ and Pop III scenarios has important implications for understanding not only the first generation of stars, but also the mechanisms by which metals are mixed and distributed from their stellar birthplaces.
Several papers (Cowie & Songaila 1998; Lu et al. 1998; Ellison et al. 1999a, hereafter Paper I) have addressed this question by attempting to probe the low density regions of the IGM where the difference in metallicity predicted by in-situ and Pop III enrichment may be most marked. The detection of C IV $`\lambda \lambda `$1548,1550 associated with low column density Ly$`\alpha `$ clouds (log $`N`$(H I) $`<`$ 14.0) is observationally challenging, even with the capabilities of Keck, due to the extreme weakness of the metal lines. Analysis techniques have therefore been developed to effectively enhance the sensitivity of the data beyond the normal equivalent width limits of the spectra. Two in particular have received recent attention, namely the production of a stacked C IV spectrum (Lu et al. 1998) and the use of individual pixel optical depths of Ly $`\alpha `$ and C IV (Cowie & Songaila 1998). In Paper I we attempted to reconcile the apparently conflicting results obtained from these two techniques with an analysis of a very high signal-to-noise ratio (S/N) spectrum of the ultra-luminous BAL quasar APM 08279+5255. Rigorous testing of the analysis procedures revealed that both methods suffered from hitherto unrecognised limitations and it was concluded that the question of whether or not the low density regions of the IGM have been enriched remains unanswered.
In order to probe deeper into the low density IGM, we have obtained an exceptionally high S/N spectrum of the well-known lensed QSO, Q1422+231, using the HIRES instrument (Vogt 1994) on the KeckI telescope. Although this QSO has been well studied in the past (e.g. Songaila & Cowie 1996; Songaila 1998), we have roughly doubled the exposure time of earlier spectra, allowing us to probe the metallicity of the Ly$`\alpha `$ forest to more sensitive levels than has previously been achieved. The spectrum is much better suited to the present work than the the much more complex BAL quasar APM 08279+5255 (Ellison et al. 1999a). We present a careful and extensive analysis of the C IV systems in order to determine the extent of metal enrichment in the IGM. This paper is organised as follows. In ยง2 we describe the observations, the data reduction procedures and the Voigt profile fitting process used to determine column densities. We briefly discuss in ยง3 the suitability of Q1422+231 for this work and define the redshift interval over which we will perform the analysis. In ยง4, we determine the column densities of the 34 detected C IV absorption systems in the spectrum and investigate the column density distribution of these absorbers. Finally, we critically assess the two methods of analysis, described in sections 5 and 6, which have been developed to probe low column density Ly$`\alpha `$ clouds to very sensitive levels. We utilize a suite of simulation techniques to fully test these methods and quantify the potential inaccuracies in our analysis.
We adopt $`\mathrm{\Omega }_M`$ = 1.0 throughout.
## 2. Data Acquisition and Reduction
Observations of Q1422+231 were made with the HIRES spectrograph (Vogt 1994) on the KeckI telescope in February 1999, using the $`1.14\times 7`$ arcsec slit. The resultant resolution is $`R=37,000`$. Individual exposure times varied between 30 min and 40 min, for a total exposure of 630 min in two (overlapping) grating settings to provide total wavelength coverage between the quasarโs Ly$`\alpha `$ and C IV emission lines. The data were reduced as described in Songaila (1998) and added to the spectrum of Q1422+231 obtained under similar conditions by Songaila and Cowie (1996). The total integration time for the two datasets is 1130 minutes and has a S/N of 200 โ 300 redward of the QSOโs Ly$`\alpha `$ emission, a quality superior to all previously published data.
After extraction and sky subtraction, we noted that the cores of the saturated absorption lines contain a very small systematic residual flux. This error in the zero level is possibly due to light from a foreground source or a small underestimate of the background sky. The correction required to bring the line cores to zero was found to vary slightly with wavelength, ranging from 0.8% for the bluest absorption lines ($`\lambda <4850`$ ร
) to 0.2% at $`\lambda 5500`$ ร
. No correction could be estimated redward of the QSOโs Ly$`\alpha `$ emission since the absence of saturated lines gave us no basis for estimating the adjustment required. However, since the correction factor required appears to diminish with increasing wavelength and is already very small at 5500 ร
, and given that such an adjustment will make little difference to the relatively weak C IV lines studied in the red, we consider any residual eror to be unimportant for the analysis presented here.
The continuum fit was achieved using the STARLINK package DIPSO with a cubic spline polynomial applied to windows of spectrum judged to be free from absorption. Voigt profiles were fitted to the entire normalized Ly$`\alpha `$ forest using VPFIT (Webb 1987) to decompose the complex absorption into individual components defined by a column density ($`N`$(H I)), a $`b`$-value (Doppler width) and a redshift. This task, though time consuming, is an important feature of our analysis since synthetic spectra can be simulated based on the line list of Voigt profile parameters. For example, additional C IV associated with the fitted H I can be included to these synthetic spectra according to any desired enrichment recipe. As will be discussed later in this paper, this is an essential step for testing the analysis techniques used here.
## 3. Q1422+231 โ Suitability and Sample Definition
Q1422+231 ($`z_{em}=3.625`$) is a well-studied quasar and actually consists of four closely spaced lensed images with separations of 0.5 โ 1.3 arcsec; the lensing galaxy is at $`z_{lens}`$= 0.34 (Patnaik et al. 1992; Kundic et al 1997; Tonry 1998). Gravitational lenses enhance the emission from high redshift QSOs making them more powerful probes of the intergalactic and interstellar medium, but obviously the sight lines will sample different spatial regions of the intervening material. Its redshift and luminosity (V=16.5) have made Q1422+231 an ideal candidate with which to probe intervening material through closely spaced sightlines (Petry et al. 1998; Rauch et al. 1999). The observations reported here are of the closely spaced A and B components. In our observations the sightlines are unresolved. It is well established that multiple lines of sight through quasar pairs separated by several arcsecs show coherence between Ly$`\alpha `$ clouds on scales $`>`$ 100 kpc (e.g. Bechtold et al. 1994; Dinshaw et al. 1997). This is consistent with the scenario that has emerged from hydrodynamical simulations that portrays structure in the IGM not as discrete localized clouds but as a smoothly fluctuating medium. For metal line systems, whilst there may be slight differences in the individual components that constitute C IV complexes, Rauch et al. (1998) have shown that the total system column density remains largely unchanged over $``$ 10 kpc. The small linear separations probed by Q1422+231 ($`<0.14`$ kpc h<sup>-1</sup> for the C IV systems in the $`z`$ range considered here) are therefore unlikely to give rise to by line of sight differences so large as to compromise our column density determinations.
In order to avoid confusion between Ly$`\alpha `$ and higher order Lyman lines, we restrict our analysis of the Ly$`\alpha `$ forest in Q1422+231 to the interval 4740 $`<\lambda `$(Ly$`\alpha `$)$`<5520`$ ร
which corresponds to a redshift range of 2.90 $`<z_{abs}<`$ 3.54. The lower limit of this interval is determined by the onset of Ly$`\beta `$ absorption and the upper limit is enforced in order to avoid effects due to quasar proximity (e.g. Lu, Wolfe & Turnshek 1991) and corresponds to a velocity separation of $`5500`$ km s<sup>-1</sup> relative to the emission redshift, $`z_{em}`$ = 3.625. In order to improve the statistics and S/N of stacked data (ยง5), we have, in some of the work presented here, included the spectrum of APM 08279+5255, recently analysed in Paper I. APM 08279+5255 is an ultra-luminous Broad Absorption Line (BAL) quasar with a systemic redshift $`z_{em}`$ = 3.911 and a broad band magnitude R = 15.2. The data obtained for this quasar (as presented in Paper I) are also of excellent quality, though the broad absorption line makes some portions of the spectrum less ideal for this problem than the spectrum of Q1422+231, and taken together with Q1422+231 they represent a premium data set for this work. The wavelength region used for APM 08279+5255 is 4995 $`<\lambda `$(Ly$`\alpha `$)$`<5720`$ ร
(3.11 $`<z_{abs}<`$ 3.70) which in addition to the selection criteria applied to Q1422+231, takes into account the BAL nature of this quasar. Finally, a region contaminated by atmospheric absorption from 6865 $`<\lambda `$(C IV) $`<`$ 6940 ร
(corresponding to $`3.43<z<3.48`$ for C IV) was excluded in both spectra.
## 4. Detected C IV Systems
Q1422+231 has been the target of several observing campaigns over the years, which have resulted in spectra of various S/N ratios. Whilst the primary motivation for obtaining such a high S/N spectrum was our scientific goal of probing the low column density Ly$`\alpha `$ forest, continued focus on this target has produced important results for IGM enrichment on a range of column density scales. In this section we present all of the detected C IV systems within the redshift range defined in ยง3, several of which have not been detected in previous spectra and we show that the power law column density distribution established for strong absorbers with $`N`$(C IV)$`>`$ 12.75 (Songaila 1997) continues to significantly lower column densities.
We detect 34 C IV systems within our defined wavelength interval, 29 of which are associated with saturated Ly$`\alpha `$ clouds (log $`N`$(H I) $`\begin{array}{c}>\\ \end{array}`$ 14.5). These systems exhibit a wide variety of column densities and complexity (i.e. number of components) as can be seen in Figure 1. The column densities, $`b`$-values and redshifts determined for each C IV component using the Voigt profile fitting program VPFIT, are presented in Table 1<sup>1</sup><sup>1</sup>1The 1548 ร
doublet component of system C11 is blended with the 1550 ร
component of C10. Both of these components fall at 6302 ร
which coincides with a strong sky line, whose subtraction has left some residuals in our spectrum. However, the C11 system is also observed in the spectrum of Boksenberg, Sargent and Rauch (in preparation) where clearly more absorption is required at 6302 ร
than is accounted for by the 1550 ร
line of C10. The Voigt profile for C11 presented here is based only on the 1550 ร
line which itself is found in a region of weak atmospheric absorption and therefore is likely to be less reliable than our other fits. Examples of C IV systems detected in our spectrum which have not been observed in other published spectra (Songaila & Cowie 1996; Songaila 1998) are the weak systems at $`z_{abs}=3.276,3.317`$ and 3.518 (C24, C25 and C33 in Figure 1).
### 4.1. Column Density Distribution of C IV Systems
Previous studies of C IV absorbers in a variety of quasar sightlines has established a power law column density distribution (with index $`\alpha 1.5`$, see eqn. 1) complete down to log $`N`$(C IV) $``$ 12.75 at $`z>3`$ and 12.25 at $`z<3`$.(Songaila 1997). The column density distribution function is defined as
$$f(N)dN=BN^\alpha dN$$
(1)
where, $`f(N)`$ is the number of systems per column density interval per unit redshift path. The redshift path (used instead of $`z`$ in order to account for co-moving distances) is given by $`X(z)=\frac{3}{2}[(1+z)^{\frac{3}{2}}1]`$ for our adopted cosmology. An important question is how much of the iceberg has been exposed? To what limit does this power law distribution continue? As spectral data have improved, early determinations of the column density distribution of C IV absorbers have been shown to be incomplete at low column densities (compare for example Petitjean & Bergeron 1994 and Songaila 1997). The superb quality of this single spectrum can address whether the established power law continues to lower $`N`$(C IV), in which case the apparent fall-off towards low $`N`$(C IV) seen previously is due to incompleteness, or whether there is a real turnover in number density. In Figure 4 we show the column density distribution of C IV derived from the systems in Q1422+231, assumed to be a power law of the form of eqn. 1. A maximum likelihood fit to the data (binned in Figure 4 for display purposes only) gives a power law index of $`\alpha =1.44\pm 0.05`$, consistent with other recent estimates. This high quality spectrum, however, clearly uncovers more of the โicebergโ than previous studies and the power law continues down to log $`N`$(C IV) $``$ 12.3 (Figure 4, solid points). Below this column density, $`f(N)`$ shows an apparent departure from the power law which may be due to incompleteness or alternatively reflect a real turnover in the $`N`$(C IV) distribution. The formal 5$`\sigma `$ detection limit for C IV in our spectrum is log $`N`$(C IV) = 11.6, for the median $`b`$-value of the observed C IV absorption lines, $`b_{median}=13`$ km s<sup>-1</sup>. This detection limit is for one C IV line and is based on the $`\lambda `$1548 ร
doublet component. However, identification of a suspected C IV system is dependent on confirmation from the weaker C IV $`\lambda `$1550 ร
line whose oscillator strength is only half of the $`f`$-value of the 1548 ร
line. This effectively reduces our sensitivity for detecting C IV systems by a factor of two to a 5$`\sigma `$ detection limit of log $`N`$(C IV) = 11.9. Moreover, lines with $`b`$-values significantly larger than the median $`b`$=13 km s<sup>-1</sup>, may not be detected. In order to estimate the incompleteness for systems with log $`N`$(C IV) $`<`$ 12.3 due to large $`b`$, 40 C IV doublets were simulated for the two bins in Figure 4 which show a departure from the established power law, i.e. $`N`$(C IV) = 12.05 and $`N`$(C IV) = 11.75. The $`b`$-value for each line was drawn at random from the real distribution of Doppler widths and noise was added at the appropriate level. Each simulated line was then inspected to assess whether it would have been identified in the original spectrum so that a correction factor could be determined to estimate incompleteness. For $`N`$(C IV) = 12.05 a correction factor of 2.4 was determined (17/40 test C IV systems detected in the incompleteness trial), with the largest $`b`$-value in the lines detected being 13 km s<sup>-1</sup>. In the lowest column density bin, $`N`$(C IV) = 11.75 only 9 out of 40 C IV lines (with $`b<`$ 8 km s<sup>-1</sup>) were identified, corresponding to a correction factor of 4.4. Clearly, these correction factors assume that the $`b`$-value distribution does not significantly change with decreasing C IV column density. The two data points adjusted for incompleteness are represented by open circles in Figure 4. Thus, it appears that the power law distribution of column densities continues down to at least log $`N`$(C IV) = 11.75.
### 4.2. Homogeneous or Variable Metallicity?
Recently, hydrodynamical simulations have been used to predict the expected scatter in individual C IV/H I ratios within strong Ly$`\alpha `$ absorbers for a fixed \[C/H\] (Hellsten et al. 1997; Rauch, Haehnelt & Steinmetz 1997, Davรฉ et al. 1998). When compared with observations, these models show that the data are consistent with a mean IGM metallicity \[C/H\] = $`2.5`$ with variations of up to a factor of 10 around this average value. In these models, the C IV systems associated with high column density Ly$`\alpha `$ clouds are dominated by metals from in-situ star formation, since simulations show that these metal enriched clouds are found within a few tens of kpc from collapsed, dense clumps at $`z3`$ (e.g. Haehnelt 1998 and references therein). We would then expect the metallicity of such clouds to be variable, being dependent on the local star formation history, and therefore the scatter amongst individual C IV/H I ratios at a given $`N`$(H I) to be larger than that predicted for a homogeneous \[C/H\].
We can calculate C IV/H I ratios for 19 systems in our spectrum of Q1422+231. Of the 34 detected C IV systems in the spectrum presented here, 29 are associated with saturated Ly$`\alpha `$ lines. Although it is not possible to determine accurate H I column densities from these lines alone, for 14 of the C IV systems higher order Lyman lines are both accessible and suitably uncontaminated by blends so that an accurate $`N`$(H I) can in fact be derived. In addition, 5 of the systems in Q1422+231 are associated with clouds with log $`N`$(H I) $`<`$ 14.5 which are not saturated so that the column densities can be determined from the Voigt profile fit of Ly$`\alpha `$ only. The C IV and H I column densities determined for these 19 systems are presented in Table 3.
Davรฉ et al. (1998) have also extensively studied the C IV/H I ratios of Q1422+231, as determined from the spectrum of Songaila & Cowie (1996). They constructed a mock spectrum of Q1422+231 from hydrodynamic simulations for comparison with the data. Having both simulated and observed spectra at their disposal, they measured C IV/H I ratios in both datasets in a consistent manner (using the AutoVP Voigt profile fitter, Davรฉ et al. 1997), and found that an intrinsic scatter of approximately 0.5 dex in metallicity is required to fit the data. However, the observed range of C IV/H I ratios is the result of many complex effects, which include not only spatial variations in the temperature-density relationship, but also, for example, fluctuations in the ionizing background. These complex effects have not yet been fully incorporated into models and therefore we should be mindful that the scatter in measured values of C IV/H I is an upper limit to the variations in \[C/H\] when compared with homogeneously enriched simulations with a uniform ionizing background.
Davรฉ et al. (1998) also found that the most robust diagnostic for determining the mean carbon abundance in detected C IV systems is $``$log $`N`$(C IV)$``$, although this statistic is clearly dependent on the sensitivity of the data. As a consistency check, we determine the $``$log $`N`$(C IV)$``$ for those C IV systems identified in our spectrum of Q1422+231 whose column densities are above the detection limit of the Songaila & Cowie (1996) spectrum (estimated to be log $`N`$(C IV $`>`$ 12.0). We calculate $``$log $`N`$(C IV)$``$ = 12.77 which compares well with the value of 12.72 determined by Davรฉ et al. (1998), considering the errors associated with column density determinations of individual C IV systems. This demonstrates that the different line finding and fitting procedures used on different spectra reproduce the same answer, even when the same systems are observed with a higher S/N. For Voigt profile fitting, this is an important point to make in a process that is sometimes not clear-cut. It is also an illustration that the detection limit of the spectrum used by Davรฉ et al. (1998) was well determined and did not suffer from serious incompleteness. Ideally, one would like to see how the mean \[C/H\] varies as we detect progressively weaker C IV systems. However, a simple extrapolation to lower column densities is not possible, since the relation between $``$log $`N`$(C IV)$``$ and \[C/H\] presented in Davรฉ et al. (1998) is tailored specifically for the detection limits of their data. In order to apply this technique to the full sample of C IV systems detected in the new Q1422+231 spectrum presented here, it would be necessary to re-create a model spectrum to match our data. However, we remind the reader that the caveats mentioned in relation to the scatter of C IV/H I are also applicable here, a fact of which one should be aware when comparing observational and simulated results and interpreting the conversion to metallicity.
We can summarise the results from this section with the following conclusions. By obtaining an ultra-high S/N spectrum of the $`z=3.625`$ quasar Q1422+231, we have detected 34 C IV absorption systems with 2.91 $`<z_{abs}<`$ 3.54. Some hitherto undetected weak C IV systems are reported whose column densities are consistent with the established $`f(N)`$ distribution showing that this power law ($`\alpha =1.44\pm 0.05`$) continues down to at least log $`N`$(C IV) = 11.75, a factor of 10 deeper than the previous determination at $`z>3`$ in Songaila (1997). By fitting down the Lyman series, we are able to determine accurate C IV/H I ratios for 19 of the identified systems, rather than relying on a statistical estimate of the median metallicity as has often been done. Davรฉ et al. (1998) have previously determined that the scatter in C IV/H I in the Q1422+231 spectrum of Songaila & Cowie (1996) requires an intrinsic scatter in the \[C/H\] of a factor of $``$ 3. Finally, by considering only the C IV systems above the detection limit, we obtain the same $``$log $`N`$(C IV)$``$ as Davรฉ et al. (1998), a statistic used to infer that the mean \[C/H\] = $`2.5`$. However, we stress that these simulations do not take into account the full complexity of the physical processes determining the C IV/H I ratio in the IGM so that both the mean \[C/H\] and its scatter may change with future more detailed work.
## 5. Stacking
As discussed in previous sections, the C IV/H I ratio in low $`N`$(H I) systems may hold the key to understanding metal enrichment mechanisms in the IGM. Since one of the major limitations in detecting weak C IV lines is S/N, stacking many sections of the spectrum is one way to circumvent this problem (e.g. Norris, Peterson & Hartwick 1983; Tytler et al. 1995).
In our application of this technique, we select Ly$`\alpha `$ lines with 13.5 $`<`$ log $`N`$(H I) $`<`$ 14.0 whose corresponding C IV spectrum shows no obvious metal lines or contamination from other absorption features. Each section is de-redshifted to the rest frame, re-binned to the dispersion of the lowest $`z`$ system, weighted according to its S/N, stacked and finally re-normalized. The optimal weighting used to co-add the C IV sections is given by
$$w_j=\frac{1/\sigma _j^2}{_j1/\sigma _j^2}$$
(2)
where $`\sigma _j^2`$ is the variance of the data. To achieve maximum sensitivity, we use low column density H I lines from both the APM 08279+5255 and Q1422+231 spectra. Within the ranges defined in ยง3, a total of 67 low $`N`$(H I) lines were identified in the two QSO sightlines. The resulting composite spectrum has a S/N = 1250 and shows no absorption at the rest wavelength of C IV, $`\lambda _0=1548.2`$, as can be seen from the top panel of Figure 5.
In order to place a significance limit on this non-detection, a synthetic spectrum was created, re-producing the Ly$`\alpha `$ forest of the 2 quasars from the fitted H I line lists. C IV was included for H I lines with log $`N`$(H I) $`<`$ 14.5 assuming a fixed C IV/H I ratio and $`b`$(C IV) = $`\frac{1}{2}b`$(Ly$`\alpha `$) (representing a combination of thermal and bulk motion, as found in Paper I). Figure 5 shows the results of stacking synthetic spectra with log C IV/H I = $`3.1`$. The resultant absorption in this stack has an equivalent width of 0.15 mร
and represents a 4$`\sigma `$ feature which we therefore adopt as the detection limit for our stacked data. Several analyses (e.g. Songaila & Cowie 1996; Paper I) have determined the C IV/H I ratio in high column density Ly$`\alpha `$ clouds to be in the range from log C IV/H I = $`2.9`$ to $`2.6`$. The detection limit of the synthetic stack is almost factor of two lower than the metal-poor limit of this range, and thus it would appear to indicate a drop in the C IV/H I ratio at lower H I column densities. It must be understood, however, that this technique has several potential problems which may compromise its efficiency in detecting weak absorption features. We now discuss these problems in turn.
* In order to co-add each section, the data must be re-binned, usually to the dispersion of the lowest redshift system. This could smooth out a weak feature, although the scale of smoothing is very small compared with the width of the expected line so that this is not likely to be a major effect.
* Even for a fixed carbon abundance there will be a scatter in the values of C IV/H I (as discussed in the previous section). Therefore, if absorption is detected in the composite spectrum it will be averaged over a range of column densities that can not be recovered individually. Moreover, depending on the scatter of C IV/H I values, the absorption may be dominated by the strong tail end of this metallicity range. In fact, it is conceivable that residual absorption could be caused by only a few relatively strong lines, since this method relies on an average which is very sensitive to a non-Gaussian tail. Interpreting a residual signal in the composite spectrum is therefore not straightforward, although in the analysis presented here we find that there is no absorption in the stacked data and so we can determine a useful detection limit.
More serious problems for the present analysis are:
* The sensitivity of this method to errors in continuum fitting, anomalous pixels and other forms of contamination. The usual procedure is to visually inspect each section before adding it to the stack in order to ensure that it is โcleanโ. This will filter out major contamination by, for example, uncorrected cosmic ray events or absorption due to systems other than C IV. However, small errors in the continuum fit or deviant pixels could seriously compromise the efficiency of the stack. In addition, a re-normalization of the composite spectrum is usually required, since the small errors in the original continuum fit have now been compounded by stacking.
* Stacking the individual sections of spectrum in order to build up a signal is pivotal upon centering the absorption feature in the composite. If there is a significant error in the stack center, i.e. if an offset exists between the redshift of the parent Ly$`\alpha `$ line and its associated C IV complex, the composite signal would be smeared out. Depending on the magnitude of the offset, this effect may lead us to under-estimate the amount of absorbing material. This problem could be exacerbated by the afore-mentioned need for a re-normalization because it may be difficult, if not impossible, to distinguish a weak smeared signal from the compound continuum errors.
### 5.1. Redshift Offset Between C IV and Ly$`\alpha `$
In Paper I, we showed that there is indeed a random redshift offset ($`\mathrm{\Delta }z`$) between the measured position of the Ly$`\alpha `$ and corresponding C IV lines. In that work it was found that the redshift offset had a distribution with $`\sigma _z`$ 4 $`\times 10^4`$ equivalent to a velocity difference of 27 km s<sup>-1</sup>, although this statistic was based on a relatively small number of systems. The addition of a second quasar to the analysis has improved the statistics and for a total of 56 C IV systems in the two quasars we now determine a $`\sigma _z`$=2.6 $`\times 10^4`$ ($`\sigma _v`$17 km s<sup>-1</sup> for an average redshift of 3.45). This redshift offset was determined in two different ways. In both cases the redshift of the C IV was estimated by taking the centroid of the system since it is often complex, consisting of several blended components. For most of the corresponding Ly$`\alpha `$ lines, the absorption is well represented by a single component and so in the first instance the redshift of the H I cloud was obtained from the Voigt profile fit. For comparison, the redshift of Ly$`\alpha `$ was also determined using the same centroid method applied to the C IV systems and it was found that both determinations yielded almost identical results for the distribution of $`\mathrm{\Delta }z`$, shown in Figure 6. Since it is the strongest component in each C IV complex that first emerges from the noise in the weak systems, this will be the feature enhanced by the stacking procedure. Therefore, we also investigated the distribution of offsets between the fitted $`z`$ of Ly$`\alpha `$ and the redshift of the deepest C IV trough and once again the value $`\sigma _v`$ 17 km s<sup>-1</sup> was found.
We investigated how seriously the composite spectrum would be affected by this offset which effectively shifts the expected positions of the C IV lines by a random amount. Again, synthetic spectra were produced, this time including a redshift offset applied to the position of the C IV line with $`\mathrm{\Delta }z`$ drawn at random from a Gaussian distribution with $`\sigma _z`$=2.6$`\times 10^4`$. We find that in order to reproduce a 4$`\sigma `$ detection in the presence of a redshift offset, twice as much C IV must be included in the Ly$`\alpha `$ forest clouds, i.e. log (C IV/H I) = $`2.8`$, see the bottom panel of Figure 5. This is consistent with the measured C IV/H I in log $`N`$(H I)$`>`$14.5 lines and is therefore not a sufficiently sensitive limit to establish whether the low column density Ly$`\alpha `$ clouds contain significantly less C IV that the stronger lines. We conclude that a factor of at least two improvement in sensitivity is required in order to show conclusively whether the low column density absorbers are more metal deficient than their high column density counterparts. Alternatively, it must be shown that the redshift offset in these lines is $``$ 17 km s<sup>-1</sup> and dilution of absorption from smearing the stack is unimportant.
The redshift offset in saturated lines could have two possible explanations, one physical and the other an observational artefact. There could be an intrinsic redshift difference between the C IV and Ly$`\alpha `$ absorbers, caused by, for example, ionization effects or outflows. An additional effect could be a redshift offset caused by the blending of strong Ly$`\alpha `$ lines which, when saturated, can not be distinguished into separate components. Two examples of this are shown in Figure 7. Plotted in velocity space are two examples of Ly$`\alpha `$ absorbers (black solid line) and corresponding C IV (dashed line). Both saturated (systems โbโ and โcโ) and weak (systems โaโ and โdโ) Ly$`\alpha `$ clouds are shown and in the top panel the Ly$`\beta `$ (in gray) is also included. The C IV system โcโ is associated with an apparently monolithic Ly$`\alpha `$ absorber centered at $`v`$ = 0 and exhibits a redshift offset of approximately 25 km s<sup>-1</sup>. However, as seen from the Ly$`\beta `$ absorption, this system is made up of more than one component and the C IV clearly has a redshift that more closely matches the strongest of these. The unsaturated Ly$`\alpha `$ absorber associated with the C IV system โaโ does not break down into a multi-component system in Ly$`\beta `$ and the redshift offset is correspondingly small ($`\mathrm{\Delta }z<`$ 1 km s<sup>-1</sup>). However, there are other examples of C IV systems associated with weak Ly$`\alpha `$ that do show a significant offset, for example system โdโ in the lower panel of Figure 7. Unfortunately, higher order Lyman lines are not available for this particular case. We re-measured the redshift offset for each of the 19 C IV systems in Table 3 for which an accurate $`N`$(H I) had been obtained by tracing down the Lyman series. For 16 of these systems, the C IV appeared to be associated with a single H I component (i.e. not obviously blended), as seen from Ly$`\beta `$ and/or Ly$`\gamma `$. The offsets associated with the sub-cloud whose redshift most closely matches that of the C IV are shown in the lower panel of Figure 6. Given the small number of systems for which the Lyman series can be traced, at least to Ly$`\beta `$, the statistics of the offset distribution are not very meaningful. However, the scatter is now clearly smaller, although the high $`N`$(H I) Ly$`\alpha `$ lines at the top end of the column density distribution may still have several sub-components that are unresolved in our data and the offset may be further reduced if the Lyman series could be traced down to subsequent transitions.
### 5.2. An Improvement on the Stacking Method?
We investigated a possible solution to the problem of an unknown offset from the predicted position of C IV systems discussed above. The technique involves scanning for the maximum absorption around the predicted position of C IV absorption and re-centering the stack at this wavelength. In practice, each (de-redshifted) data section is scanned $`\pm 2\sigma _z`$ from the predicted C IV line center (i.e. for $`\mathrm{\Delta }z`$ = 0) and the section re-centered on the pixel with the maximum optical depth ($`\tau _{max}`$) prior to stacking. Figure 8 shows an example of how this technique is clearly able to improve upon a direct stack in the presence of a $`\mathrm{\Delta }z`$ for a synthetic spectrum containing lines with log C IV/H I =$``$2.0. It can be seen that without re-centering, whilst the overall level of the continuum is below unity for the โsmearedโ stack, all profile information is lost. Such a broad depression may be mistaken for compound errors in the continuum level and lost in the subsequent renormalization of the stacked spectrum. Executing a re-center on the $`\tau _{max}`$ pixel before stacking recovers almost all of the original signal in this simulation, the characteristic line profile is prominent and will not be lost when the post-stacking continuum is re-fitted.
However, as one attempts to detect progressively weaker absorption, the likelihood of centering on a noise pixel becomes higher and this technique is no longer efficient. In order to determine whether a โre-centerโ is a viable improvement to the stacking method, given the S/N of the data and the low column density of the targeted features, we performed a feasibility study in which test absorption lines were created over a range of C IV/H I ratios. The results of this study are shown in Figure 9 where we plot the distribution of $`\sigma _z`$ caused by line center misidentification (i.e. where $`\tau _{max}`$ is a noise feature rather than line trough) as a function of C IV/H I ratio. Clearly, when the $`\sigma `$ of the offset caused by trough misidentification exceeds the value of the intrinsic offset we are trying to overcome, this re-centering technique is no longer useful. However, these are conservative estimates for how well the centering would perform on real data, since we would expect some trough misidentification from contaminating (i.e. anything but C IV) lines in addition to the effects of noise investigated here. Since we need to determine whether the low column density Ly$`\alpha `$ lines have the same metallicity as their high $`N`$(H I) counterparts, we must ideally reach levels of sensitivity deeper than log C IV/H I=$`2.9`$. Figure 9 shows that for log C IV/H I $`<`$2.6, this technique no longer compensates adequately for the intrinsic offset and is therefore unable to improve the stacking technique in the search for weak lines. Smoothing the spectrum over several pixels before locating $`\tau _{max}`$ improves the pre-shifting slightly as shown by the gray squares in Figure 9. However, even with smoothing, this re-centering procedure can not compensate for a $`\sigma _z`$ = 2.6 $`\times 10^4`$ below a log C IV/H I = $`2.75`$.
In summary, by stacking together the C IV regions associated with low column density Ly$`\alpha `$ lines in Q1422+231 and APM 08279+5255, we have obtained a S/N = 1250 composite spectrum which shows no residual C IV absorption. Using simulated spectra we find that for log C IV/H I = $`3.1`$ in the weak Ly$`\alpha `$ lines we would expect a 4$`\sigma `$ detection of the composite absorption. The shortcomings of this technique are discussed and we focus on an observed redshift offset between the position of Ly$`\alpha `$ and its associated C IV. The redshift offset determined from the detected C IV systems has a dispersion $`\sigma _z`$ = 2.6$`\times 10^4`$ which corresponds to a $`\sigma _v`$ = 17 km s<sup>-1</sup>. We have investigated how this offset would affect the stacking procedure by using simulated spectra and find that a $`\sigma _v`$ = 17 km s<sup>-1</sup> will reduce the sensitivity of this method by a factor of two and that a metallicity of log C IV/H I = $`2.8`$ is now required to achieve a 4$`\sigma `$ detection. We are unable to know whether the same redshift offset persists in the low $`N`$(H I) clouds targeted by the stacking technique, but a large range of offsets are measured over the full column density range of detected C IV systems, so one must clearly take into account the possible repercussions when interpreting the stacked data.
## 6. Analysis of Optical Depths
A potentially more robust way of measuring weak absorption in high S/N spectra is to analyse the optical depths in each Ly$`\alpha `$ forest pixel and its corresponding C IV pixel (Cowie & Songaila 1998). This technique also has the advantage that by tracing down the Lyman series, one can determine the C IV/H I ratio in each pixel over a very large range of $`\tau `$(Ly$`\alpha `$). In Paper I we investigated the potential of this method with APM 08279+5255 and critically analysed its performance on synthetic spectra. Specifically, we investigated the effect of including a redshift offset and how the results depended on the $`b`$-value of C IV. We concluded that, despite the excellent quality of the spectrum, the results from the spectrum of APM 08279+5255 were inconclusive and could not determine whether the metallicity of the IGM is constant or diminished at low values of H I, because at the lowest values of $`\tau `$(Ly$`\alpha `$) both scenarios were consistent with the observations. The spectrum of Q1422+231 presented here not only has a significantly higher S/N than that of APM 08279+5255, but has many of its higher order Lyman lines accessible for analysis and considerably less contaminating absorption by, for example, Mg II systems (Ellison et al. 1999b). This spectrum therefore represents the best data yet obtained for this analysis of the low column density Ly$`\alpha `$ forest.
### 6.1. The Analysis Procedure
Briefly, the optical depth technique consists of stepping through the Ly$`\alpha `$ forest measuring $`\tau `$(Ly$`\alpha `$) for each pixel. The noise ($`\sigma `$) array is used to determine which pixels are included in the analysis via a series of optical depth criteria to account for effects such as saturation. Using the values in the noise arrays rather than fixing the rejection criteria provides maximum flexibility for this technique so that it can be readily applied to spectra of different S/N ratios.
For pixels with a residual flux ($`F=e^\tau `$ for a normalised spectrum) $`F<3\sigma `$ above the zero level, we trace down the Lyman series since there is too little residual flux in these saturated pixels to determine an accurate optical depth from Ly$`\alpha `$ alone. However, the danger here is that higher order lines may be contaminated by lower redshift Ly$`\alpha `$. We therefore use $`\tau `$(Ly$`\alpha `$) = minimum($`\tau `$(Ly$`n`$)$`f_{\mathrm{Ly}\alpha }\lambda _{\mathrm{Ly}\alpha }`$ / $`f_{\mathrm{Lyn}}\lambda _{\mathrm{Lyn}}`$) over all observed higher order lines using the higher order pixel if $`3\sigma <F`$(Ly$`n)<(13\sigma )`$. This minimizes the effect of contamination and maximizes the number of usable pixels and range of $`\tau `$(Ly$`\alpha `$) which can be considered for analysis. If no $`3\sigma <F`$(Ly$`n)<(13\sigma )`$ pixels are found, then the pixel is discarded. The position of the associated C IV $`\lambda \lambda `$1548, 1550 lines are calculated and, again to avoid contamination, we use $`\tau `$(1548)= minimum($`\tau (1548)`$,$`2\tau `$(1550))<sup>2</sup><sup>2</sup>2The ratio of the line $`f`$-values (oscillator strengths) in the C IV doublet is 2:1 if the flux in the 1550 ร
component $`<13\sigma `$, otherwise only $`\tau `$(1548) is considered. The C IV optical depths are then binned according to their corresponding $`\tau `$(Ly$`\alpha `$) and the median determined for each interval. Taking the median of a large number of pixel optical depths not only provides a statistical advantage over considering a relatively small number of lines (as for the stacking method), but also is much less susceptible to non-Gaussian effects. In order to estimate 1 $`\sigma `$ errors for the optical depth determinations, we used bootstrap re-sampling with 2ร
sections of the Ly$`\alpha `$ forest and corresponding C IV, i.e. we drew $`n`$ random sections from the complete set of $`n`$ sections (in this case $`n=361`$) that comprise the original data, with replacement. This procedure was repeated 250 times and the 1$`\sigma `$ error was taken to be the dispersion of these 250 realizations.
### 6.2. Results
Figure 10 shows the results obtained for the optical depth analysis of Q1422+231. The shape of the optical depth distribution appears consistent with a constant level of C IV/H I (i.e. parallel to the dashed line) over optical depths from $`\tau `$(Ly$`\alpha `$) $``$ 100 down to $``$ 2 โ 3, below which $`\tau `$(C IV) flattens off to an approximately constant value. Each of the low optical depth bins ($`\tau `$(Ly$`\alpha `$) $`<`$ 3) contains approximately 10% of the H I pixels. This percentage decreases with increasing $`\tau `$(Ly$`\alpha `$) with only approximately 2% of pixels in the $`\tau `$(Ly$`\alpha `$) = 10 bin. In addition, we determine optical depths in pixel pairs separated by the C IV$`\lambda \lambda `$1548, 1550 doublet ratio regardless of $`\tau `$(Ly$`\alpha `$) over the entire range considered for C IV absorption. This is done using the same method of doublet comparison as before to eliminate contamination. The median of this reference distribution is plotted in Figure 10 as a dotted line and represents the median absorption for an effectively random set of pixel pairs separated by C IV doublet ratio and will include the effects of noise and low level contamination expected to affect our results. This median optical depth is higher than the observed $`\tau `$(C IV) for all bins with $`\tau `$(Ly$`\alpha `$)$`\stackrel{<}{}`$ 1. This is a significant result which indicates that the C IV absorption for these optical depths is less than that expected by selecting random pixel pairs. The key question here is whether the signal in the low optical depth H I pixels is due to C IV absorption in low density regions of the IGM or whether the flattening of the data is caused by some limiting factor in our analysis such as contamination or noise.
### 6.3. Testing the Results
As we saw in the previous section, thorough simulations of the analysis technique are vital for interpreting our results. Here, we have taken the Ly$`\alpha `$ forest directly from the data and artificially enriched it with C IV to test our methods and address several questions. In Paper I we explored the effect of $`b`$-values, redshift offset and noise on synthetic spectra in order to determine whether a break in the C IV/H I distribution could be distinguished from a constant ratio in the spectrum of APM 08279+5255. To review these findings and to give a visual impression of the optical depth analysis, we show the results of this technique on four synthetic spectra in Figure 11 and compare them with the results from Q1422+231 (solid points). The top panel (โAโ) of this figure shows the optical depth analysis of 2 spectra, one of which has a constant log C IV/H I = $`2.6`$ in all Ly$`\alpha `$ lines (shown as a dotted line) and a second spectrum which has log C IV/H I = $`2.6`$ only in log $`N`$(H I) $`>`$ 14.5 lines (dot-dashed line), with no noise added to either spectrum redward of Ly$`\alpha `$. The bottom panel shows the same spectra as panel โAโ except that noise has now been added to both spectra, based on the error array of the actual data (typically S/N = 200 redward of Ly$`\alpha `$). All four spectra use the real Ly$`\alpha `$ forest of Q1422+231 and have had C IV added to the synthetic spectrum based on the fitted H I linelist. In addition, all four synthetic spectra in Figure 11 have a $`\mathrm{\Delta }z`$ = 2.6 $`\times 10^4`$ and $`b`$(C IV) = $`\frac{1}{2}b`$(Ly$`\alpha `$). The dashed line indicates constant log C IV/H I = $`2.6`$.
There are several points to note here. Firstly, in the absence of noise, there is a clear drop in $`\tau `$(C IV) below $`\tau `$(Ly$`\alpha `$) $`3`$ if no C IV is added in Ly$`\alpha `$ lines with log $`N`$(H I) $`<`$ 14.5. However, this steep decline is much less drastic when noise is included, and although $`\tau `$(C IV) shows a steady decrease down to $`\tau `$(Ly$`\alpha `$) $``$ 1, below this value it flattens off to an approximately constant value. The inclusion of noise also causes the same apparent flattening in the spectrum of constant C IV/H I, although the $`\tau `$(C IV) in this spectrum is consistently above the value measured for the dot-dashed line. There is also a small contribution to this flattening from line blending which changes the overall slope of the expected C IV/H I, an effect that is exacerbated for larger $`b`$(C IV). This shows that even in relatively high S/N spectra, distinguishing a break in the C IV/H I distribution is very difficult. Instead, as we shall see later in this section, one of the main uses of this method is to determine whether the optical depths measured with this technique can be adequately accounted for with the detected C IV systems or whether there must be significant amount of C IV still below our current detection limit. We also note that for large $`\tau `$(Ly$`\alpha `$) the measured $`\tau `$(C IV) is always less than expected from the dashed line, given the input ratio of C IV/H I. This is due to two effects. Firstly, since this technique considers $`\tau `$(Ly$`\alpha `$) as the minimum value obtained by tracing down the Lyman series when a line is saturated, if contamination is successfully removed we will tend to under-estimate $`\tau `$(Ly$`\alpha `$) at these optical depths due to noise (remember that the real Ly$`\alpha `$ forest is used in the simulated spectra so that this will be an effect even in panel โAโ of Figure 11). Secondly, the C IV that is included in the synthetic spectra is based on a linelist of fitted H I values which will probably not be accurate for saturated lines. These effects also account for the drop of $`\tau `$(C IV) at very large $`\tau `$(Ly$`\alpha `$) which contain $`<`$ 1% of the total H I pixels and are therefore very uncertain.
Clearly, the exact results of the optical depth analysis are sensitive to the combination of several factors including blending, $`b`$-values and noise, which we will discuss further later in this section. Therefore, rather than attempting to directly fit the observed distribution of optical depths in Q1422+231 we aim to determine whether the optical depths measured in the spectrum can be explained solely by the relatively strong C IV absorbers detected directly and presented in ยง4 or whether the results from this analysis are indicative of additional C IV. If the latter is true, are these metals in the low density IGM? We must also investigate whether our results can be explained in terms of contamination, noise or some other limiting factor in the data. Finally, we consider the effect of scatter and redshift offset ($`\mathrm{\Delta }z`$), which have been found to compromise the efficiency of the stacking method.
Three synthetic spectra were produced with Ly$`\alpha `$ forest absorption taken directly from the data (i.e. not reconstructed from a linelist) and the 34 detected C IV systems re-produced from the Voigt profile parameters in Table 1. Spectrum โAโ has had no further metals added and therefore shows the results expected of an optical depth analysis if we had already uncovered all of the C IV in the spectrum. Spectra โBโ and โCโ have both been enriched with additional metals. In spectrum โBโ, C IV is included in all strong (log $`N`$(H I) $`>`$ 14.5) lines with $`N`$(C IV) = 12.0, i.e. the detection limit for directly identifiable C IV systems. This spectrum therefore represents the maximum amount of C IV that could be โhiddenโ in high column density Ly$`\alpha `$ lines. In addition to this โhiddenโ C IV in strong H I lines and the fitted C IV in Table 1, spectrum โCโ has been enriched with a constant log C IV/H I = $`2.6`$ in weak Ly$`\alpha `$ lines ($`N`$(H I) $`<`$ 14.5). The results from analysis of these three synthetic spectra are compared with the data in Figure 12. For all three spectra, we include noise taken from the $`\sigma `$ error array, a random redshift offset in the position of C IV ($`\sigma _z=2.6\times 10^4`$) and take $`b`$(C IV) = $`\frac{1}{2}`$ $`b`$(Ly$`\alpha `$).
The conclusions that can be drawn from Figure 12 are as follows. First, there is clearly more C IV in the data than we have directly identified in ยง4, since the dotted line of the synthetic spectrum in the top panel is well below the solid line at all but the very highest $`\tau `$(Ly$`\alpha `$) points. For $`\tau `$(Ly$`\alpha `$) $`\begin{array}{c}>\\ \end{array}`$ 3, the C IV optical depths can be recovered when additional metals are added into strong Ly$`\alpha `$ absorbers, but below our current detection limit (spectrum โBโ). This supports the results of ยง4 where we determine that $`f(N)`$ is consistent with a power law distribution that continues down to log $`N`$(C IV) = 11.75, with no evidence for a turnover in column density distribution. However, adding C IV at the limit of our detection in log $`N`$(H I)$`>`$ 14.5 lines with no extra metals in weaker Ly$`\alpha `$ clouds can not reproduce the measured $`\tau `$(C IV) in $`\tau `$(Ly$`\alpha `$)$`<`$ 3 pixels. The results from spectrum โCโ in the bottom panel of Figure 12 which include log C IV/H I = $`2.6`$ in weak Ly$`\alpha `$ lines show that the observed C IV optical depths in the low $`\tau `$(Ly$`\alpha `$) pixels are consistent with the presence of some metals in significantly lower column density H I clouds, although for $`\tau `$(Ly$`\alpha `$) $`>`$ 1 $`\tau `$(C IV) exceeds the observed value. From Figure 11 we have also seen that a constant log C IV/H I = $`2.6`$ in all Ly$`\alpha `$ lines is a good approximation to the data. We stress here that these simulations are not an attempt to fit the observed distribution of optical depths, rather they are tests to determine whether the C IV systems in Table 1 can account for the measured absorption. If not, then the objective is to investigate in which H I column density regime additional metals could be added in order to achieve the observed quota.
The results from these simulations show that the analysis of optical depths in the spectrum of Q1422+231 is consistent with hitherto undetected C IV in both the strong and weak Ly$`\alpha `$ lines. Whilst the results from our determination of the column density distribution from detected C IV systems is consistent with a power law function $`f(N)`$ that continues down to at least log $`N`$(C IV) = 11.7, these optical depth results provide direct evidence that there are more metals lurking below our current detection limit.
We investigate the effect that contaminating lines may have on this result. The analysis may be affected not only by low level contamination from other metal lines such as Si IV or telluric absorption (strong absorbers will have been rejected by the C IV doublet strength comparison discussed in ยง6.1), but also effects due the non-uniformity of ionizing sources and noise/fluctuations in the continuum level due to fitting errors. This latter effect will cause fluctuations around the true continuum which will be important only if they are greater than the level of the noise. An additional (possibly systematic) error may be present if the low order polynomial fit to the continuum is consistently over or under-estimated. To test the continuum fit of the C IV regions, sections of the data redward of Ly$`\alpha `$ deemed to contain no obvious absorption were examined. From a total of $``$ 5000 pixels, a mean flux of 1.0002 and a median of 1.0001 were determined, vaules one order of magnitude lower than $`\tau `$(C IV) $`10^3`$. It therefore seems unlikely that a systematic error in the continuum fit is the cause of the observed flattening of C IV optical depths seen in Figure 10. The distribution of noise pixels would also suggest that there is no significant effect from weak telluric lines. Errors in the continuum fit blueward of Ly$`\alpha `$ emission, however, are likely to be a more serious effect, firstly because selection of absorption-free zones is difficult and secondly because the S/N is lower (50 โ 150). The result of continuum fitting errors in the forest in this analysis will be to classify pixels with small $`\tau `$(Ly$`\alpha `$) into the wrong optical depth bin. This will effectively associate C IV with the wrong $`\tau `$(Ly$`\alpha `$) and averaging out the metal absorption by this โmis-binningโ could explain the observed flattening of the measured optical depths. The effect of a small continuum error is therefore similar to the effect of noise, affecting only those pixels with optical depths smaller than the fitting error. We also find that including a large number of weak contaminating lines such as Si IV or Mg II could increase the measured $`\tau `$(C IV) at low $`\tau `$(Ly$`\alpha `$) to the constant level determined. However, whilst realistic column densities do have a small effect on the optical depth result, it is unlikely that weak contaminating lines account for all of the observed absorption.
As discussed previously, a scatter is expected in the observed C IV/H I values even if \[C/H\] and the ionizing background remain uniform. To investigate the importance of this effect, we simulated a โcontrolโ spectrum by re-producing directly the Ly$`\alpha `$ forest of Q1422+231 and using $`b`$(C IV) = $`\frac{1}{2}b`$(Ly$`\alpha `$) and an input log C IV/H I = $`2.6`$ with no scatter and no redshift offsets. For comparison, we then produced a second spectrum in the same way, but rather than adopting a constant C IV/H I, a Gaussian scatter was introduced with $`\sigma `$(C IV/H I) = $`1\times 10^3`$. The comparison between the two simulated spectra is presented in Figure 13 where the โcontrolโ spectrum is shown as open squares connected with a solid line and the spectrum containing a scatter of metallicities is shown as solid diamonds connected with a dotted line. The points have been offset from one another in the figure for clarity. The points at high $`\tau `$(Ly$`\alpha `$) are again very uncertain due to the reasons previously discussed with regards to Figure 11. The error bars are estimated using the same bootstrap technique as employed for the data. The results from the 2 spectra are entirely consistent with one another and indicate that a Gaussian scatter will therefore not affect the overall median of $`\tau `$(C IV). Whilst the extent and magnitude of the scatter in the distribution of C IV/H I values has not yet been fully investigated in simulations down to column densities below log $`N`$(H I) $``$ 14.0 โ 13.5, this result should not be significantly affected unless the scatter becomes highly non-Gaussian at low $`N`$(H I). Also plotted in Figure 13 are the results of the optical depth analysis if a redshift offset is included in the position of the C IV line. As before, metals are added to all Ly$`\alpha `$ lines with a metallicity log C IV/H I = $`2.6`$. There is no scatter in these values, but an offset in the position of C IV has been included, drawn at random from a Gaussian distribution with $`\sigma _z=2.6\times 10^4`$ ($`\sigma _v`$ = 17 km s<sup>-1</sup>). These points are also consistent with the control spectrum, so we can conclude from this simulation batch that neither scatter in the metallicity nor redshift offset will have a large effect on the outcome of our optical depth analysis.
Finally we note that more spectra are required in order to provide a representative analysis of the high redshift IGM, since our view of the enrichment of the Ly$`\alpha `$ forest probed with a single sightline is clearly blinkered. From the list of C IV systems in Table 1 it can be seen that these absorbers are not uniformly distributed in redshift. For example, splitting this spectrum in two by redshift produces vastly different optical depth distributions due to a relative dearth of C IV systems that spans over 300 ร
in this spectrum in the range $`3.13<z_abs<3.33`$. With more spectra, there is the possibility of not only developing a more representative study of these high redshift absorbers but also investigating the redshift evolution of C IV systems.
## 7. Conclusions
In this paper, we have addressed the enrichment history of the IGM by studying the Ly$`\alpha `$ forest and its associated C IV systems in a very high S/N ($`200)`$, high resolution ($``$ 8 km s<sup>-1</sup>) spectrum of the well-known lensed quasar, Q1422+231 obtained with Keck/HIRES. The numerous C IV systems associated with high column density Ly$`\alpha `$ absorbers are fitted with Voigt profiles defined by a redshift, $`b`$-value and column density for each component line. We investigate the C IV column density distribution, $`f(N)`$, to very sensitive levels and detect several weak C IV systems which had not been previously identified in lower quality spectra of the same object. We determine a power law index $`\alpha =1.44\pm 0.05`$ which continues down to log $`N`$(C IV) = 12.2 before starting to turnover. By simulating synthetic absorption lines with $`b`$-values taken at random from the observed distribution, we estimate a correction factor to account for incompleteness and find that the corrected data points now indicate that the power law continues down to at least log $`N`$(C IV) = 11.75, a factor of ten more sensitive than previous measurements (e.g. Songaila 1997). This shows that even at these low column densities there is no evidence for a flattening of the power law and therefore there are probably many more C IV systems that lie below the current detection limit.
We investigate two methods with which it may be possible to recover these weak C IV systems. Firstly, we select 67 Ly$`\alpha `$ lines with 13.5 $`<`$ $`N`$(H I) $`<`$ 14.0 in Q1422+231 and APM 08279+5255 and produce a stacked spectrum centered on the predicted position of C IV $`\lambda `$1548. The composite stack has a S/N = 1250 and shows no residual absorption; we use synthetic stacked spectra to determine a 4$`\sigma `$ upper limit of log C IV/H I = $`3.1`$. We critically assess the accuracy of this method by performing the stacking procedure on a suite of simulated, synthetic spectra and identify several associated problems. We investigate, in particular, the effect of a redshift offset between the position of the Ly$`\alpha `$ line and its associated C IV. With improved statistics, we refine the redshift offset determined in Paper I and find that a $`\sigma _v`$=17 km s<sup>-1</sup> is present in the C IV systems which we detect directly. By including a random redshift offset drawn from a Gaussian distribution with $`\sigma _v`$=17 km s<sup>-1</sup> in our stacking simulations, we find that log C IV/H I = $`2.8`$ is now required to achieve a 4$`\sigma `$ detection. This limit is still consistent with current measured metallicities in higher column density Ly$`\alpha `$ clouds and is therefore not sufficiently sensitive to determine whether the C IV/H I ratio drops in low $`N`$(H I) lines. A feasibility study is performed to assess the effectiveness of a โre-centerโ on the maximum optical depth pixel prior to stacking for removing the effect of an unknown offset. We find that this technique can not improve the quality of the stack result in the C IV/H I regime that we are targeting. It is not yet clear whether the observed redshift offset persists in the low column density Ly$`\alpha `$ clouds, but it must be considered a factor. Moreover, the effects of contamination, continuum fitting errors and anomolous pixels also pose problems for this technique, although in simulations the redshift offset appears to be a major effect. Therefore, if this technique is to be pursued further to reach a meaningful detection limit, an improvement in S/N by at least a factor of two is required.
The optical depth technique introduced by Cowie & Songaila (1998) is considered as an alternative approach. This technique is shown to exhibit several advantages over the stacking method such as its insensitivity to redshift offsets and its ability to exclude contamination from other absorption features. We develop this technique as a method that can be used to test whether the detected C IV systems represent the full tally of absorbers. The data have optical depths consistent with an almost constant log C IV/H I $`3`$ down to $`\tau `$(Ly$`\alpha `$)$``$ 2โ3, below which $`\tau `$(C IV) flattens off to an approximately constant value. It is unlikely that this flattening is real and is most probably caused by the effect of noise and/or continuum errors, even at such high S/N ratios as have been achieved in this spectrum. Given the many effects that may alter the measured $`\tau `$ values, such as blending and noise, we do not attempt to fit the observed distribution of optical depths. Instead, our strategy is to test whether the detected C IV systems are sufficient to reproduce the measured $`\tau `$(C IV) and if not, determine how much additonal C IV may be present below our current detection limit. By simulating synthetic spectra with different enrichment recipes, we have shown that the C IV systems detected directly in the spectrum are not sufficient to reproduce the results of the optical depth analysis of Q1422+231. This is in agreement with the conclusions drawn from the column density distribution of C IV, i.e. that the data are consistent with a continuous power law $`f(N)`$ down to at least $`N`$(C IV) = 11.75 and that there is therefore likely to be a large number of weak metal lines not yet directly detected. This agrees with the conclusions of Cowie & Songaila (1998).
In order to interpret the results from the optical depth method, we have simulated synthetic spectra with a range of input C IV/H I ratios. We find that including C IV associated with strong Ly$`\alpha `$ lines ($`N`$(H I) $`>`$14.5) but below the current detection limit, in addition to the 34 identified C IV systems, can reproduce the optical depths measured in the observed spectrum for $`\tau `$(Ly$`\alpha `$) $`>`$ 3. For smaller values of $`\tau `$(Ly$`\alpha `$), some additional metals are required and we find that including C IV in low column density H I lines ($`N`$(H I) $`<`$ 14.5) with log C IV/H I $`=2.6`$ produces optical depth results consistent with those measured in the data. However, determining the precise C IV/H I in the low $`N`$(H I) Ly$`\alpha `$ clouds and the density to which the enrichment persists is still uncertain due to effects such as noise and continuum fluctuations and it is therefore not possible to say whether the low column density forest is pristine. Nevertheless, we find that even in the high optical depth H I pixels (which will not be seriously affected by noise or small continuum errors) the detected C IV systems are not sufficient to cause all the measured absorption and clearly there are more metals in the IGM than we can currently detect.
The authors would like to thank Bob Carswell, Alberto Fernandez-Soto, Matthias Steinmetz and the anonymous referee for discussions and useful suggestions. We are grateful to Len Cowie for his support and enthusiasm towards this project. SLE and JS are supported by a PPARC postgraduate award and JS acknowledges additional funding from the Isaac Newton Trust. We are also grateful to the expert assistance of the staff at the Keck I telescope for their help in obtaining the spectra presented here.
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# Focal Loci of Algebraic Varieties I
## 1 Introduction
The goal of the present paper is to introduce a general theory of focal loci of algebraic varieties in Euclidean space.
The theory of focal loci was classically considered only for plane curves and surfaces in 3-space ( cf. \[Coolidge\] , \[Salmon-Fiedler\]), and Hilbert himself lectured in the Winter Semester 1893-94 at the University of Gรถttingen on the focal loci of curves and surfaces of degree two in 3-space.
Recently the theory was considered in (\[Fantechi\], \[Trifogli\]) for the respective cases of plane curves and hypersurfaces.
We would like to first briefly present the relevant concepts.
Usually the focal locus of a submanifold $`X`$ ( cf. \[Milnor\], 6, pp. 32-38, or also \[D-F-N\], vol. II 11, sections 2-3) is defined in Euclidean differential geometry as either the locus of centres of principal curvatures, or, more geometrically, as the locus where the infinitely near normal spaces intersect each other. Equivalently, the focal locus can also be defined as the complement of the set of points $`p`$ such that the square of the distance function from $`p`$ induces a local Morse function on $`X`$, or also as the union of the singular points of the parallel varieties to $`X`$.
To make the definition algebraic, one picks up the second geometrical definition, where the notion of length is not needed, just the notion of orthogonality is sufficient.
To explain this in more detail, let us consider (complex) affine space as the complement of a hyperplane ( the โhyperplane at infinityโ) in projective space. In the hyperplane at infinity $`๐_{\mathrm{}}`$, we give a non degenerate quadric $`Q_{\mathrm{}}`$.
These data allow , for each projective linear subspace $`L`$ , to define the orthogonal $`L^_x`$ to $`L`$ in a point $`x`$ as the join of $`x`$ with the โorthogonal directionโ to $`L`$ ( this is the subspace of $`๐_{\mathrm{}}`$ given by the polar of $`L๐_{\mathrm{}}`$ with respect to $`Q_{\mathrm{}}`$).
Given now an irreducible algebraic variety $`X_d^n๐^๐ฆ`$, of dimension $`n`$ and degree $`d`$ and not contained in the hyperplane at infinity , for each smooth point $`xX๐_{\mathrm{}}`$ we define the normal space $`N_x(X)`$ as the orthogonal in $`x`$ of the projective tangent space to $`X`$ at $`x`$. The condition that $`x`$ is a point in affine space ensures that $`N_x(X)`$ has the correct dimension $`mn`$.
The normal variety $`N_X`$ is then defined as the irreducible algebraic set in $`๐^m\times ๐^m`$, closure of the set $`N_X^{good}`$ consisting of the pairs $`(x,y)`$ where $`x`$ is a smooth point of $`X`$, $`xX๐_{\mathrm{}}`$ and $`yN_x(X)`$.
Clearly, $`N_X`$ is a projective variety of dimension $`m`$ and the second projection induces a map $`\pi `$ whose image is the closure of the union of the normal spaces to the smooth points of $`X๐_{\mathrm{}}`$. Observe moreover that $`N_X^{good}`$ is a projective bundle over $`X๐_{\mathrm{}}Sing(X)`$, in particular $`N_X^{good}`$ is smooth of dimension $`m`$ : therefore we can consider the ramification locus $`Y_X^{good}`$ of $`\pi :N_X^{good}๐^m`$, and we define the ordinary ramification locus as the closure $`Y_X`$ of $`Y_X^{good}`$.
Defining the good focal locus as $`\mathrm{\Sigma }_X^{good}=\pi (Y_X^{good})`$, and the focal locus $`\mathrm{\Sigma }_X`$ as the closure of $`\mathrm{\Sigma }_X^{good}`$ ( thus $`\mathrm{\Sigma }_X`$ is contained in the branch locus of $`\pi :N_X๐^m`$), we have a priori at least four cases:
* 1) $`\pi :N_X๐^m`$ is not dominant : in this case we say that the variety $`X`$ is isotropically focally degenerate (for short : isotropically degenerate), and observe that the focal locus $`\mathrm{\Sigma }_X`$ of $`X`$ is then simply the image of $`\pi `$ ( whence, $`\mathrm{\Sigma }_X`$ is an irreducible variety in this case!).
* 2) $`\pi :N_X๐^m`$ is dominant , but the focal locus $`\mathrm{\Sigma }_X`$ (respectively, the branch locus of $`\pi :N_X๐^m`$) has dimension at most $`m2`$ : in this case we say that $`X`$ is strongly focally degenerate ( respectively, completely strongly focally degenerate).
* 3) $`\pi :N_X๐^m`$ is dominant , whence surjective , and the focal locus $`\mathrm{\Sigma }_X`$ โis not a hypersurfaceโ, in the sense that not every component $`Z`$ of the ordinary ramification divisor $`Y_X`$ (closure of $`Y_X^{good}`$) maps to a hypersurface. In this case we shall say that $`X`$ is weakly focally degenerate. We shall moreover say that we have the vertical case if $`Z`$ does not dominate $`X`$.
* 4) When none of the above occurs, in particular $`\pi :N_X๐^m`$ is surjective , and the focal locus $`\mathrm{\Sigma }_X`$ is a hypersurface, we shall say that $`X`$ is focally non degenerate. In this case, defining the focal hypersurface as a divisor, consisting as the image of the ramification divisor $`Y_X`$ with multiplicities (if $`Y_X=\mathrm{\Sigma }_{i=1,..k}n_iY_i`$ , and $`d_i:=degree(Y_i\pi (Y_i)`$, then, setting $`\mathrm{\Sigma }_i:=\pi (Y_i)`$, we get $`\mathrm{\Sigma }_X:=\mathrm{\Sigma }_{i=1,..k}d_in_i\mathrm{\Sigma }_i`$), the main problem is to describe $`\mathrm{\Sigma }_X`$.
The first main result of this paper consists in calculating the degree (with multiplicity) of the focal hypersurface under a certain hypothesis upon $`X`$, which we call of being โorthogonally generalโ, and which ensures that $`X`$ is focally non strongly degenerate if it is not a linear subspace. This concept is important because, if $`X`$ is smooth and not a linear subspace, then for a general projectivity $`g`$ the translate $`g(X)`$ of $`X`$ by $`g`$ satisfies this condition whence it is not focally strongly degenerate and we have a divisor $`\mathrm{\Sigma }_X`$. The hypothesis that $`X`$ be โorthogonally generalโ is indeed very easy to verify since it simply amounts to three requirements: the smoothness of $`X`$, plus the two general position properties that $`X`$ be transversal to $`๐_{\mathrm{}}`$, respectively to $`Q_{\mathrm{}}`$.
More precisely, we have the following Theorem
Theorem 1 Let $`X๐^m`$ be a variety of dimension $`n1`$ which is orthogonally general. Then $`dim\mathrm{\Sigma }_X<m1X`$ is a linear space. If $`X`$ is a linear space, $`\mathrm{\Sigma }_X`$ is a linear space of dimension equal to $`codimX1`$.
One can ask in the above theorem whether one can replace the condition $`dim\mathrm{\Sigma }_X<m1`$ (i.e., that $`X`$ be strongly focally degenerate) by the weaker condition that $`X`$ be focally degenerate.
As a corollary of the full description given in Theorem 3 of the focally degenerate varieties, it turns out that if $`X`$ is an orthogonally general and focally degenerate variety, then either $`X`$ or $`X_{\mathrm{}}`$ should be a developable variety rather explicitly described, but we have not yet had the time to look at the existence question for such very special varieties.
It is rather clear ( e.g., from the case of plane curves) that the condition of being orthogonally general is a sufficient but not necessary condition in order that $`X`$ be non focally strongly degenerate. When $`X`$ is non orthogonally general, but focally non degenerate, what happens is that the degree of the focal divisor can drop ( in this case, for plane curves we have Plรผcker type formulae, cf. \[Fantechi\]).
Naturally, what we have said insofar opens a series of problems. To some of them we give an answer in the present paper, to some others we hope to return in a sequel to this paper :
* 1) Try to completely classify the focally isotropically degenerate varieties. In section 7 we give a structure Theorem ( Theorem 4) stating that the isotropically focally degenerate hypersurfaces are exactly the isotropically developable hypersurfaces. We observe thus that there are plenty of intriguing examples already in the case of surfaces in 3-space: these are obtained as the tangential developable surface of any space curve whose tangent direction is always an isotropic vector. We give moreover a description in section 8, Theorem 5, of the general case, in terms of the inverse focal construction applied to the focal variety $`\mathrm{\Sigma }`$ and to an algebraic function $`r`$ on $`\mathrm{\Sigma }`$. We get thus an implicit classification of these varieties as developable varieties, but for this we need to start with a variety $`\mathrm{\Sigma }`$ whose normal spaces are totally isotropic, and the function $`r`$ must also satisfy a suitable condition.
* 2) Try to classify the weakly and the strongly focally degenerate varieties. In section 6 we give a complete classification for the weakly focally degenerate varieties, showing in Theorem 3 how they can be divided into someโprimitiveโ classes ( cases 1), 2), 6), 7)) and some โderivedโ classes (cases 3),4),5)), related for instance by some tangential conditions to some primitive focally degenerate varieties. The primitive focally degenerate varieties can be described starting from fibrations in spheres or in affine spaces โaroundโ the degenerate component $`\mathrm{\Sigma }`$ of the focal locus.
The question of classifying the strongly focally degenerate varieties seems harder.
* 3) Determine whether for a general projective deformation of $`X`$ the focal hypersurface is reduced of degree equal to the virtual degree, and moreover answer more specific questions such as :
* 3a) can we also obtain that for a general deformation the focal hypersurface has generic Lagrangian singularities ?
* 3b) can we obtain the above good properties for the focal hypersurface $`\mathrm{\Sigma }_{gX}`$ of a general translate $`gX`$ of $`X`$ by a general projectivity $`g`$ ?
Concerning the first problem, the situation seems to us rather hard (although quite interesting) as soon as the dimension of the ambient space grows: for instance, whereas a focally isotropically degenerate plane curve $`C`$ is necessarily a line through a cyclic point ( these are the two points of $`Q_{\mathrm{}}`$ , satisfying the equations $`z=x^2+y^2=0`$ ), in the case of a surface in 3-space we obtain the tangential developable of a space curve $`C`$ which is โisotropicโ in the following sense: $`C`$ is just a curve such that any of its tangent lines $`L`$ has the property that $`L`$ intersects $`๐_{\mathrm{}}`$ in a point of $`Q_{\mathrm{}}`$. Therefore, if we write the point of the curve $`C`$ as a vector function $`x(t)`$ of a parameter $`t`$, we just have to solve the differential equation
$`Q_{\mathrm{}}(dx/dt)=0`$.
Thus such a curve $`C`$ yields a curve $`\mathrm{\Gamma }`$ in $`Q_{\mathrm{}}`$ parametrizing the projective tangent lines to $`C`$, and the question reduces to: for which $`\mathrm{\Gamma }`$ can one find an algebraic integral ? (however, since the ring of polynomials in $`t`$ is stable by $`d/dt`$ , the above observation easily allows us to construct a lot of focally isotropically degenerate surfaces, which are tangential surfaces of rational space curves, cf. Example 10).
In higher dimension, as we already remarked, Theorem 5 partially reduces the quest to the search of varieties with totally isotropic normal spaces.
Turning to the other problems, the situation is clear for the plane curves (cf. \[Fantechi\]) : the only focally degenerate plane curves, which are not lines, are the circles (conics through the two cyclic points), and moreover, for an irreducible plane curve $`C`$ the map of $`C`$ to the focal curve $`\mathrm{\Sigma }_C`$ is non birational exactly for a well classified class of curves (by the way, Fantechi shows that this class is non empty, contrary to a statement made in \[Coolidge\]).
As we said, we characterize (cf. Theorem 3) the weakly focally degenerate varieties, distinguishing six essentially different cases :
* two vertical cases, where the exceptional component $`R`$ of $`Y_X`$ does not dominate $`X`$, but is instead the restriction of the normal bundle $`N_X`$ to a divisor $`X\mathrm{"}`$. In both cases, $`X\mathrm{"}`$ is focally degenerate, and the focal degeneracy of $`X`$ is determined by the first order neighbourhood of $`X`$ along $`X\mathrm{"}`$ (see Theorem 3 for more details).
* the case where $`X`$ consists of a family of $`(m1a)`$-dimensional spheres parametrized by the $`a`$-dimensional degenerate component $`\mathrm{\Sigma }`$ of the focal locus: this family is moving according to a simple differential equation which can be explicitly solved, and it turns out that we get a family of spheres each obtained as the intersection of the big sphere with centre $`O\mathrm{\Sigma }`$ with an affine subspace orthogonal to the tangent space to $`\mathrm{\Sigma }`$ in $`O`$.
* The case where $`X`$ is a โtransversalโ divisor in a focally isotropically degenerate variety.
* The asymptotic case, i.e., the case where $`\mathrm{\Sigma }`$ lies at infinity, and then $`X`$ is a developable variety whose intersection $`X_{\mathrm{}}`$ with the hyperplane at infinity โisโ the dual variety of $`\mathrm{\Sigma }`$ in $`๐_{\mathrm{}}`$. In this case there is another simple process, called the โasymptotic inverse focal constructionโ, describing $`X`$ in terms of the data of $`\mathrm{\Sigma }`$ and of an algebraic function $`r(s)`$ on $`\mathrm{\Sigma }`$.
* The isotropically asymptotic case, where $`\mathrm{\Sigma }`$ lies at infinity, and a component $`\mathrm{\Delta }`$ of $`X_{\mathrm{}}`$ is projectively isotropically degenerate. This case is characterized by the property that $`\mathrm{\Delta }X_{\mathrm{}}`$ be obtained via the isotropic projective inverse focal construction, starting from $`\mathrm{\Sigma },r(s)`$ satisfying suitable conditions.
The characterization given in Theorem 3 (where also the case of the focally isotropically degenerate varieties is considered) is expressed in terms of the โinverse focal constructionโ, which, starting from a variety $`\mathrm{\Sigma }`$ of dimension $`a`$, and an algebraic function $`r(s)`$ on $`\mathrm{\Sigma }`$, considers the union $`X^{}`$ of the family of spheres each obtained as the intersection of the big sphere with centre $`O\mathrm{\Sigma }`$ and radius equal to the square root of $`r(s)`$ with an affine subspace orthogonal to the tangent space to $`\mathrm{\Sigma }`$ in $`O`$, and whose position is determined by the differential of the function $`r(s)`$.
It turns out that for the focally isotropically degenerate varieties the above spheres degenerate to affine spaces,and $`X`$ equals $`X^{}`$, whereas in the case where these spheres have the right dimension $`m1a`$ $`X^{}`$ is focally degenerate.
For hypersurfaces in higher dimensions the second author ( \[Trifogli\]) showed that the focal hypersurface of a general hypersurface is reduced (indeed that this holds for a general diagonal hypersurface, i.e., for a translate of the Fermat hypersurface by a projectivity in the diagonal torus).
Concerning problem 3a), this is a global problem which is however related to a local problem which has been extensively studied: the theory of Lagrangian singularities. In fact the Normal variety $`N_X`$ is a Lagrangian variety for the symplectic form on the product $`๐^๐ฆ\times ๐^๐ฆ`$ which is associated to $`Q_{\mathrm{}}`$, namely $`{}_{}{}^{t}xQ_{\mathrm{}}y^tyQ_{\mathrm{}}x`$, and the second projection is also Lagrangian ( cf. \[Arnold et al.\]).
Partial results concerning problem 3a) have been obtained by the second author for surfaces in 3- space( \[Trif2\]).
## 2 Notation
$`V^{}`$ := a fixed vector space of dimension $`m`$
$`V`$ := the vector space $`V=V^{}๐`$
$`๐(V)=๐^m`$ := the projective space whose points correspond to the 1-dimensional vector subspaces of $`V`$
$`๐(V^{})=๐_{\mathrm{}}๐(V)`$ $`(๐^{๐ฆ\mathrm{๐}})`$ the complement of the affine space $`๐(V)๐_{\mathrm{}}V^{}`$.
$`X_d^n๐^๐ฆ`$ a quasi-projective algebraic variety of dimension $`n`$ and degree $`d`$ which does not lie at infinity , i.e.,
$`X_d^n๐_{\mathrm{}}`$
$`Q_{\mathrm{}}`$ := a non degenerate quadratic form on $`V^{}`$, yielding an isomorphism
$`Q:V^{}\stackrel{}{}(V^{})^{}`$.
By slight abuse of notation, the corresponding quadric
$`Q_{\mathrm{}}๐_{\mathrm{}}`$.
$`W`$ := a vector subspace of $`V^{}`$,
$`Ann(W)`$ := the vector subspace of $`V^{}`$ which is the orthogonal space of $`W`$ with respect to the quadratic form $`Q_{\mathrm{}}`$
$`L^{}`$ := $`๐(W)`$ a linear subspace at infinity
$`L_{}^{}{}_{}{}^{}`$ := $`๐(Ann(W))`$ , the polar subspace of $`L^{}`$ , also called the orthogonal direction to $`L^{}`$
$`L๐(V)`$ := a projective linear subspace , $`L๐_{\mathrm{}}`$ the direction of $`L`$
$`L^_x`$ := the orthogonal to $`L`$ in $`x`$, defined as the smallest linear subspace containing $`x`$ and the orthogonal direction of $`L`$ ( i.e., the polar of $`L๐_{\mathrm{}}`$) .
## 3 โNormal Bundleโ in Euclidean Setting
In this section, we shall consider a smooth quasi- projective variety $`X_d^n๐^๐ฆ`$ and we shall define its projective normal variety $`N_X๐^m\times ๐^m`$, and its Euclidean Normal sheaf $`๐ฉ_X`$.
Under some assumptions that we are going to specify, the first projection of the normal variety $`N_X`$ to $`X`$ yields a projective bundle over $`X`$, which is the projectivization of the Euclidean Normal sheaf :
$`N_X=๐(๐ฉ_X)๐(V๐ช_X)๐(V๐ช_{๐^m})=๐^m\times ๐^m`$.
Start from the Euler sequence
$`(1)`$ $`0๐ช_๐(1)V๐ช_๐T_๐(1)0`$ :
setting $`=๐ช_X(1)`$, the restriction to $`X`$ of ยฟthe Euler sequence and the inclusion of the tangent bundle of $`X`$ in the restricted tangent bundle of $`๐^m`$ define the bundle $`\widehat{T}_X(1)`$ whose projectivization is the projective tangent bundle to $`X`$.
We get thus two exact sequences, the second included into the first:
$`(2)`$ $`\begin{array}{c}0V๐ช_XT_๐(1)๐ช_X0\hfill \\ 0\widehat{T}_X(1)T_X(1)0\hfill \end{array}`$
$`\underset{ยฏ}{Assumption\mathrm{\hspace{0.17em}0}=smoothness}`$: $`X`$ is smooth, whence $`T_X`$ and $`\widehat{T}_X`$ are subbundles.
Recalling that $`V=V^{}๐`$,we state the further
$`\underset{ยฏ}{Assumption\mathrm{\hspace{0.17em}1}}`$ ( = transversality of the intersection $`X๐_{\mathrm{}}`$ with the hyperplane at infinity) : $`\overline{T}_x:=V^{}\widehat{T}_x`$ is a hyperplane in $`\widehat{T}_x`$ $`xX`$.
This means that we have two more exact sequences
$`(3)`$ $`\begin{array}{c}0V^{}๐ช_XV๐ช_X๐ช_X0\hfill \\ 0\overline{T}_X(1)\widehat{T}_X(1)๐ช_X0\hfill \end{array}`$
At this stage we can define the bundle of normal directions $`๐ฉ_X^{}`$ as a twist of the annihilator of $`\overline{T}_X`$.
We define it through the exact sequence
$`(4)`$ $`0๐ฉ_X^{}(1)V^{}๐ช_X(V^{})^{}๐ช_X(\overline{T}_X(1))^{}0`$
In order to obtain a projective normal bundle from the bundle of normal directions we need a last
$`\underset{ยฏ}{Assumption\mathrm{\hspace{0.17em}2}}`$ ( = transversality of $`X`$ with $`Q_{\mathrm{}}`$) : The natural map
$`๐ฉ_X^{}(1)V๐ช_X`$ is a bundle embedding, thus its image $`๐ฉ_X(1)`$ is a subbundle of $`V๐ช_X`$, isomorphic to $`๐ฉ_X^{}(1)`$
We notice thus that if assumption 2) holds, then $`๐ฉ_X๐ช_X๐ฉ_X^{}`$
###### Definition 1
$`X`$ is said to be orthogonally general if it satisfies Assumptions $`02`$ above.
###### Remark 1
For every algebraic variety $`X๐^๐ฆ`$ which is not contained in the hyperplane at infinity there is a maximal nonempty Zariski open set $`U`$ of $`X`$ which is orthogonally general ($`U`$ obviously contains the open set $`X๐_{\mathrm{}}Sing(X)`$).
###### Remark 2
The situation can be slightly generalized as follows : let $`Z`$ be a singular projective variety, let $`Z^{}`$ be its normalization, and let $`X`$ be the open set of $`Z^{}Sing(Z^{})`$ where the natural morphism to $`๐^m`$ has maximal rank. In this case, restrictions of bundles have to be understood as pull backs. If instead one wants to generalize to the case where $`X`$ is the resolution of $`Z`$, many things change substantially because one does not get bundle maps any longer.
Thus we can give the following definition
###### Definition 2
Let $`X`$ be an algebraic variety, not contained in the hyperplane at infinity, $`U`$ a Zariski open set of $`X`$ which is orthogonally general, and $`N_U`$ the projective normal bundle of $`U`$. Then the projective normal variety $`N_X`$ of $`X`$ is defined as the Zariski closure of $`N_U`$.
We can easily verify that the above definition is indeed independent of the choice of $`U`$.
Assume now that $`X`$ is orthogonally general: in particular, $`X`$ is smooth and we have a vector bundle (locally free sheaf) $`๐ฉ_X`$ on $`X`$, which is called the EUCLIDEAN NORMAL BUNDLE of $`X`$.
###### Remark 3
The Euclidean Normal Bundle differs from the usual Normal Bundle (of a smooth subvariety $`X๐^๐ฆ`$) defined in algebraic geometry( cf. \[Hartshorne\]): the reader may in fact notice that their respective ranks differ first of all by $`1`$. However, as we shall shortly see in the forthcoming example, they are somehow related to each other.
We can therefore compute now the total Chern class of $`๐ฉ_X`$:
$`c(๐ฉ_X)=c(๐ฉ_X^{})`$ and $`c(๐ฉ_X(1))=c()c(\overline{T}_X(1)^{})^1`$ by $`(4)`$.
But $`(3)`$ yields $`c(\overline{T}_X(1))=c(\widehat{T}_X(1))`$ which by $`(2)`$ equals $`c()c(T_X(1))`$. Thus
$`c(๐ฉ_X(1))=c()c()^1c(\mathrm{\Omega }_X^1(1))^1`$
Let us verify this formula for a hypersurface of degree $`d`$. Then we have
$`0๐ช_X(1d)\mathrm{\Omega }_๐^1(1)๐ช_X\mathrm{\Omega }_X^1(1)0`$
and $`c(\mathrm{\Omega }_P^1(1))=c(๐ช(1))^1`$.
So, for a hypersurface, the rank $`2`$ bundle $`๐ฉ_X`$ has
$`c(๐ฉ_X(1))=c()c()^1c()c(๐ช_X((d1)))=(1H)(1(d1)H)`$
( indeed, the previous formulae show $`๐ฉ_X๐ช_X๐ช_X((d2))`$).
In general we have an exact sequence
$`0๐ฉ_X^{}(1)\mathrm{\Omega }_๐^1(1)๐ช_X\mathrm{\Omega }_X^1(1)0`$,
where $`๐ฉ_X^{}`$ is the usual conormal bundle of $`X`$.
Hence, $`c(๐ฉ_X(1))=c()c()^1c(๐ฉ_X^{}(1))c()`$, and we obtain the
$`\underset{ยฏ}{FINALFORMULA}`$: $`c(๐ฉ_X)=c(๐ฉ_X^{}(2))`$.
###### Corollary 1
If $`X`$ is a general complete intersection of degrees $`d_1,\mathrm{}d_{mn}`$, then $`c(๐ฉ_X)=_i(1(d_i2)H)`$, where $`H`$ is the hyperplane divisor.
We recall once more the definition of the Focal Locus $`\mathrm{\Sigma }_X`$ of $`X`$.
###### Definition 3
Continue to assume that $`X`$ is orthogonally general, let $`N_X๐^m\times ๐^m`$ be the projectivization of the Euclidean Normal Bundle, and let $`\pi =p_2:N_X๐^m`$ be the second projection. Denote then by $`Y_X`$ the ramification locus of $`\pi `$ (recall: $`N_X`$ is smooth and dim $`N_X=m`$). Clearly, if $`X`$ is projective, $`Y_X\mathrm{}`$, since $`rkPic(N_X)2`$, and therefore $`\pi `$ cannot be an isomorphism. We define in general the focal locus as $`\mathrm{\Sigma }_X:=\pi (Y_X)`$.
###### Definition 4
Let now $`Z`$ be any projective variety, possibly singular. Let $`X`$ be a maximal orthogonally general open set of the normalization $`Z^{}`$ of $`Z`$ (cf. remark 2): then the focal locus $`\mathrm{\Sigma }_Z`$ is defined as the closure of $`\mathrm{\Sigma }_X`$.
###### Remark 4
In order to verify whether the definition would be the same when one would replace $`X`$ by any orthogonally general open set of $`Z`$,i.e., independent of the chosen open set $`X`$, we may observe:
* For $`X`$ orthogonally general, the projective normal bundle $`N_X`$ has a canonical section, provided by the diagonal of $`X`$, and corresponding to the tautological sheaf $`๐ฉ_X(1)`$.
* In a neighbourhood of the canonical section, the morphism $`\pi `$ is of maximal rank if and only if $`๐ฉ_X^{}(1)`$ and $`(\overline{T}_X(1))`$ yield a direct sum, i.e., $`(\overline{T}_X(1))`$ contains no isotropic vectors.
We shall say that a point $`xX`$ is totally non isotropic if the above situation occurs.
It follows that, in the open set of totally non isotropic points, the ramification divisor cannot contain the fibre of the projection to $`X`$. Therefore, in this locus, the ramification divisor is the closure of its restriction to the inverse image of an open set in $`X`$.
Instead, when there is a divisor $`D`$ of isotropic points of $`X`$, the inverse image of $`D`$ may yield a component of the ramification divisor, as happens in the following example.
Consider the plane curve $`C`$ given, in a standard system of Euclidean coordinates, by the parametrization $`(t,it+t^3)`$.
Then the normal vector is proportional to the vector $`(i+3t^2,1)`$ and the endpoint map $`\pi `$ associates to $`(t,\lambda )`$ the point
$`x=t+\lambda (i+3t^2)`$
$`y=it+t^3\lambda `$,
and the Jacobian determinant equals
$`J=(1+6t\lambda )(i+3t^2)^2=3t(2\lambda +2it+3t^3)`$.
Thus the focal locus consists of the evolute $``$ ( image of the curve $`\lambda =it3/2t^3`$) and of the isotropic line $`\{(x,y)|ixy=0\}`$.
$``$ is here the parametrical curve $`(2t9/2it^39/2t^5,2it+5/2t^3)`$.
The previous remark and example justify the following
###### Definition 5
Let now $`Z`$ be any projective variety, possibly singular. We say that an open set $`X`$ of $`Z`$ is excellent if $`X`$ is orthogonally general and $`X`$ is contained in the set of totally non isotropic points . If there exists an excellent open set $`X`$, then the strict focal locus $`\mathrm{\Sigma }_Z^s`$ is defined as the closure of $`\mathrm{\Sigma }_X`$.
We define instead the large focal locus $`\mathrm{\Sigma }_Z^L`$ as the branch locus of the second projection $`\pi `$ from $`N_Z^{}Z^{}\times ๐^m`$ to $`๐^m`$, where as before $`Z^{}`$ is the normalization of $`Z`$.
Obviously one has inclusions $`\mathrm{\Sigma }_Z^s\mathrm{\Sigma }_Z\mathrm{\Sigma }_Z^L`$.
###### Example 1
In the case of a plane curve $`C`$, the strict focal locus is precisely the evolute of the curve $`C`$, as in \[Fantechi\]. Whereas, even if all the points are totally non isotropic, the large focal locus can be larger, as we shall now see in the case where the curve has as a singularity a higher order cusp.
Let our curve $`C`$ be locally given by $`(t^2,t^5)`$, with respect to some standard Euclidean coordinates; then the normal vector is , for $`t0`$, proportional to $`(5t^4,2t)`$, i.e., to $`(5t^3,2)`$, and thus the large focal locus is provided by the image of the jacobian determinant of the map
$`x=t^25t^3\lambda `$,
$`y=t^5+2\lambda `$.
The equation of the Jacobian determinant equals therefore
$`t(430t\lambda 25t^6)=0`$, whence the large focal
locus consists of the evolute plus the line obtained for $`t=0`$, namely the $`y`$ axis.
###### Remark 5
Assume now that $`Z`$ is any projective variety and assume that there is a non empty excellent open set $`XZ`$. If $`\mathrm{\Sigma }_Z^L`$ has dimension $`m2`$, then $`\pi `$ is a birational morphism, since then $`_1(๐^m\mathrm{\Sigma }_Z^L)=\{1\}`$. Thus if $`Z`$ is not isotropically focally degenerate and $`dim\mathrm{\Sigma }_Z^L<m1N_X`$ is rational $`Z`$ is unirational, and indeed stably rational.
###### Example 2
If $`X`$ is a smooth hypersurface of degree $`d`$ and $`\mathrm{\Sigma }_X^L`$ has dimension $`m2`$, then $`dm`$.
###### Example 3
If $`X`$ is a smooth complete intersection of multidegree $`(d_1,\mathrm{},d_{mn})`$, and $`\mathrm{\Sigma }_X^L`$ has dimension $`m2`$, then $`d_im`$.
###### Remark 6
Let $`X^{}๐^m`$ be a smooth variety not necessarily satisfying the non degeneracy conditions, i.e., Assumptions $`1`$ and $`2`$. Then $`g\mathrm{๐๐๐}(m+1)`$ such that $`X=gX^{}`$ satisfies the non degeneracy conditions.
Proof
The non degeneracy conditions are equivalent to $`(1^{})`$ $`X`$ is transversal to $`๐_{\mathrm{}}`$ and $`(2^{})`$ $`X_{\mathrm{}}:=X๐_{\mathrm{}}`$ is transversal to $`Q_{\mathrm{}}`$. By Bertiniโs theorem, we can find a hyperplane $`H`$ and a smooth quadric $`QH`$ such that $`X^{}`$ is transversal to $`H`$ and $`X^{}H`$ is transversal to $`Q`$. Then choose $`h,k\mathrm{๐๐๐}(m+1)`$ such that $`hH=๐_{\mathrm{}}`$ and $`k๐_{\mathrm{}}=๐_{\mathrm{}}`$, $`khQ=Q_{\mathrm{}}`$ and set $`g=kh`$. $`\mathrm{}`$
Let us continue now to assume that $`X`$ is orthogonally general. Moreover, we shall from now on assume that $`X`$ is indeed projective. Then we can calculate $`deg\mathrm{\Sigma }_Xdeg\pi _{|Y_X}`$ (notice that $`\pi `$ is a morphism) by working in the Chow (or cohomology) ring of $`N_X`$.
Observe that, by the Leray-Hirsch theorem, the cohomology algebra of the projective normal bundle is generated by $`H^{}(X)`$ and the relative hyperplane divisor $`H_2`$, and holds
$`H^{}(N_X)H^{}(X)[H_2]/(c_i(๐ฉ_X(1))H_2^{mn+1i})`$
We denote by $`\mathrm{\Pi }`$ the first projection $`\mathrm{\Pi }:N_XX`$, and for commodity we also set $`p:=\pi `$.
Let $`H_1=\mathrm{\Pi }^{}(hyperplane)`$, and observe that, since $`๐ฉ_X(1)`$ is a subbundle of $`V๐ช_X`$, we have $`H_2=p^{}(hyperplane)`$.
Moreover, setting $`N=N_X`$, we have also the ramification formula
$`Y=K_Np^{}(K_{๐^๐ฆ})=K_N+(m+1)H_2`$.
In order to determine the canonical divisor $`K_N`$ of $`N=N_X`$, we write as usual
$`K_N=K_{N|X}+\mathrm{\Pi }^{}(K_X)`$,
where $`K_{N|X}`$ can be calculated through the Euler exact sequence for the relative tangent bundle $`T_{N|X}`$ of $`N`$
$`0๐ช_N(H_2)\mathrm{\Pi }^{}(๐ฉ_X(1))T_{N|X}(H_2)0`$,
whence
$`K_{N|X}=c_1(T_{N|X})=[c_1(T_{N|X}(H_2)+(mn)H_2)]=(mn+1)H_2\mathrm{\Pi }^{}(c_1(๐ฉ_X(1)))=(mn+1)H_2+H_1\mathrm{\Pi }^{}(c_1(๐ฉ_X^{}(1)))=(mn+1)H_2+H_1c_1(๐ฉ_X^{})(mn)H_1=(mn+1)H_2+H_1+(m+1)H_1+K_X(mn)H_1`$
In the end we obtain:
$`K_N=2\mathrm{\Pi }^{}K_X(mn+1)H_2+(n+2)H_1`$
thus we get the
$`\underset{ยฏ}{CLASSFORMULA}`$: $`Y_X=2\mathrm{\Pi }^{}K_X+nH_2+(n+2)H_1`$,
and the
$`\underset{ยฏ}{DEGREEFORMULA}`$:
$`(deg\mathrm{\Sigma })(degp_{|Y})=H_2^{m1}(2\mathrm{\Pi }^{}K_X+nH_2+(n+2)H_1)`$.
In the sequel ( section 5) we shall see how the above cited Leray-Hirsch Theorem allows to make the degree formula more explicit.
## 4 Non degeneracy of Focal Loci
Throughout this section we assume that $`X`$ is projective and orthogonally general, i.e., the non degeneracy conditions $`02`$ above are satisfied, in particular we have that $`N=N_X`$ is a bundle . Our aim is then to determine for which $`X`$ it is possible that $`\mathrm{\Sigma }_X`$ is degenerate, that is, has dimension strictly less than $`m1`$. It is easy to see that, if $`X`$ is a linear space, then $`\mathrm{\Sigma }_X`$ is degenerate and is a linear space of dimension equal to $`codimX1`$. In what follows, we shall prove that if $`X`$ is orthogonally general also the converse holds, i.e., if $`\mathrm{\Sigma }_X`$ is strongly degenerate, then $`X`$ is a linear space.
###### Notation 1
Let $`๐^๐ฆ=๐^๐ฆ๐_{\mathrm{}}`$, $`N_{\mathrm{}}=p^1(๐_{\mathrm{}})`$, $`N_a=p^1(๐^๐ฆ)`$.
We have
###### Lemma 1
$`dimp^1(y)=0`$ $`y๐^๐ฆ`$
Proof
After identifying $`p^1(y)`$ with the set $`\mathrm{\Gamma }=\{xX:yN_x\}`$, it is easy to see that $`\mathrm{\Gamma }`$ has empty intersection with the hyperplane $`๐_{\mathrm{}}`$. Indeed, if $`xX_{\mathrm{}}`$, then $`N_x๐_{\mathrm{}}`$ $`\mathrm{}`$
###### Corollary 2
If $`\mathrm{\Sigma }`$ is a component of the focal locus, image of a component $`Y`$ of the divisor $`Y_X`$, and $`dim\mathrm{\Sigma }<m1`$, then
$`(i)`$ $`YN_{\mathrm{}}`$ (since $`y\mathrm{\Sigma }`$ dim $`p^1(y)>0`$)
$`(ii)`$ $`\mathrm{\Sigma }๐_{\mathrm{}}`$ (hence $`\mathrm{\Sigma }`$ is degenerate).
$`(iii)`$ if for every component $`\mathrm{\Sigma }`$ of the focal locus holds $`dim\mathrm{\Sigma }<m1`$, then
$`p:N_a๐^๐ฆ`$ is an isomorphism.
###### Remark 7
The divisor $`N_{\mathrm{}}`$ splits as $`N_{|X_{\mathrm{}}}N^{}`$, where $`N^{}=๐(๐ฉ_X^{})`$.
Let us first consider the case where $`X`$ is a curve ( for this case we shall give a different proof in the sequel, showing that then either $`C`$ is a line, or $`C`$ is a circle, what contradicts the hypothesis that $`C`$ be orthogonally general).
CASE: $`X=C`$ curve.
Let $`C`$ be an irreducible ( and orthogonally general) curve of degree $`d`$. Then $`N_{|C_{\mathrm{}}}`$ consists of $`d`$ distinct copies of $`๐_{\mathrm{}}`$, $`p:N_{|C_{\mathrm{}}}๐_{\mathrm{}}`$ is a finite map, and by the transversality of $`C`$ to $`๐_{\mathrm{}}`$, the divisor $`Y_C`$ does not contain any component of $`N_{|C_{\mathrm{}}}`$.
Therefore we get
###### Corollary 3
If $`C`$ is an irreducible ( and orthogonally general) curve and $`dim\mathrm{\Sigma }_C<m1`$, then $`Y_C=N^{}`$ (set-theoretically)
Proof
Indeed, $`Y_CN_{\mathrm{}}`$, but no component of $`N_{|C_{\mathrm{}}}`$ is contained in $`Y_C`$. Thus $`Y_CN^{}`$. We can conclude since $`dimY_C=dimN^{}`$ and $`N^{}`$ is irreducible (being a projective bundle on the curve $`C`$). $`\mathrm{}`$
###### Proposition 1
Assume again that $`C`$ is an irreducible ( and orthogonally general) curve . Then $`(1)`$ $`dim\mathrm{\Sigma }<m1`$ $`(2)`$ $`C`$ is a line.
$`(3)`$ In this case $`\mathrm{\Sigma }`$ is a linear space of dimension $`m2=codimC1.`$
Proof
(1) $``$ (2) being clear, letโs prove the other implication $`(1)`$ $``$ (2):
Let $`N_p^{}`$ be the fibre of $`N^{}`$ over $`pC`$, which is a hyperplane in $`๐_{\mathrm{}}`$. Now $`\mathrm{\Sigma }_C=p(N^{})`$ is irreducible, has dimension $`<m1`$ and contains $`N_p^{}`$, which has dimension equal to $`m2`$. Therefore, $`\mathrm{\Sigma }_C=N_p^{}`$ and $`N_p^{}=N_q^{}`$ $`p,qC`$. This implies $`T_pC๐_{\mathrm{}}=T_qC๐_{\mathrm{}}`$ $`p,qC`$.
This clearly implies that $`C`$ is a line, since then for each point $`pC`$ the projective tangent line $`T_pC`$ is the join of $`p`$ and of a fixed point $`p_{\mathrm{}}`$ (whence one can find then $`m1`$ independent linear forms vanishing on $`C`$).
$`\mathrm{}`$
CASE: $`dimX=n2`$
Since $`X_{\mathrm{}}`$ is smooth, by Bertiniโs theorem $`X_{\mathrm{}}`$ is irreducible. Therefore also $`N|_X_{\mathrm{}}`$ and $`N^{}`$ are irreducible.
We have
###### Lemma 2
If $`n2`$ and $`dim\mathrm{\Sigma }_X<m1`$, then
(i) $`Y_X=N^{}`$ set-theoretically.
(ii) $`[Y_X]=[nN^{}]`$ in $`Pic(N)`$.
(iii) $`2(K_X+(n+1)H)=0`$ in $`Pic(X)`$, where $`H=H_1`$ is the hyperplane divisor on $`X`$.
Proof
Since $`p`$ is surjective, we have one and only one of the following two cases: (a) $`p(N|_X_{\mathrm{}})๐_{\mathrm{}}`$; (b) $`p(N^{})๐_{\mathrm{}}`$. But (a) cannot hold. Indeed, since $`[H_1]=[N|_X_{\mathrm{}}]`$ in $`Pic(N)`$, (a) implies $`[Y]=[\alpha H_1]`$ for some $`\alpha >0`$. But then from the class formula $`()`$ $`Y_X=2\mathrm{\Pi }^{}K_X+nH_2+(n+2)H_1`$, it follows that $`2\mathrm{\Pi }^{}K_X+nH_2+(n+2\alpha )H_1=0`$, contradicting the Leray-Hirsch Theorem.
Therefore, (b) holds and hence $`Y=N^{}`$ set-theoretically. (ii) and (iii) follow immediately from the class formula $`()`$, because $`H_2=H_1+[N^{}]`$ in $`Pic(N)`$. $`\mathrm{}`$
ยฟFrom point $`(iii)`$ it follows that
###### Corollary 4
If $`dim\mathrm{\Sigma }_X<m1`$, then $`X`$ is a linear space.
Proof
Let $`C=XH_1\mathrm{}H_{n1}`$ be a smooth curve. By successive applications of the adjunction formula $`(iii)`$ yields $`2(K_C+2H)=0`$. Extracting degrees, we get $`2(2g(C)2+2deg(C))=0`$, which is equivalent to $`g(C)=0`$ and $`deg(C)=1`$. $`\mathrm{}`$
We can conclude
###### Theorem 1
Let $`X๐^m`$ be a projective variety of dimension $`n1`$ which is orthogonally general. Then $`dim\mathrm{\Sigma }_X<m1X`$ is a linear space. In this case, $`\mathrm{\Sigma }_X`$ is a linear space of dimension equal to $`codimX1`$.
## 5 The Degree of the Focal Locus of a Surface
Let $`X^2=S๐^m`$ be a surface and assume that $`S`$ satisfies the non-degeneracy conditions. Setting $`n=2`$ in the Degree-Formula given in Section $`1`$, we get ( recall $`H=H_1`$)
$`(F1)`$ $`deg\mathrm{\Sigma }_Sdegp_{|Y_S}=2H_2^{m1}(K_S+H_2+2H)`$
Our first aim in this section is to express the right-hand side of $`(F1)`$ in terms of the Chern classes $`c_1(S)`$, $`c_2(S)`$ and of the hyperplane divisor $`H`$ of $`S`$.
By the Leray-Hirsch theorem $`H_2^{m1}=c_1(๐ฉ_S(1))H_2^{m2}c_2(๐ฉ_S(1))H_2^{m3}`$.
Using this relation, the right-hand side of $`(F1)`$ becomes
$`()`$ $`2H_2^{m2}(c_1(๐ฉ_S(1))^2c_2(๐ฉ_S(1))K_Sc_1(๐ฉ_S(1))2Hc_1(๐ฉ_S(1)))`$
Recall that $`c(๐ฉ_S(1))=c()c(๐ฉ_S^{}(1))`$, where $`=๐ช_S(1)`$ and $`๐ฉ_S^{}`$ is the conormal bundle of $`S`$. Thus,
$`(1)\begin{array}{c}c_1(๐ฉ_S(1))=c_1(๐ฉ_S^{}(1))H=c_1(๐ฉ_S^{})+(m3)H\hfill \\ c_2(๐ฉ_S(1))=c_2(๐ฉ_S^{}(1))Hc_1(๐ฉ_S^{}(1))=\hfill \\ c_2(๐ฉ_S^{})+(m4)Hc_1(๐ฉ_S^{})+\frac{1}{2}(m2)(m5)H^2\hfill \end{array}`$
Using the normal-bundle sequence we get
$`(2)\begin{array}{c}c_1(๐ฉ_S^{})=c_1(S)(m+1)H\hfill \\ c_2(๐ฉ_S^{})=c_2(S)+\frac{1}{2}m(m+1)H^2+c_1(S)c_1(๐ฉ_S^{})=\hfill \\ c_2(S)+\frac{1}{2}m(m+1)H^2+c_1(S)^2(m+1)Hc_1(S)\hfill \end{array}`$
and substituting in $`(1)`$, we get
$`(3)\begin{array}{c}c_1(๐ฉ_S(1))=4H+c_1(S)\hfill \\ c_2(๐ฉ_S(1))=9H^25Hc_1(S)+c_1^2(S)c_2(S)\hfill \end{array}`$
Hence $`()`$ becomes
$`()`$ $`2H_2^{m2}(15H^2+c_1^2(S)+c_2(S)9Hc_1(S))`$
We recall that $`H_2=[N^{}]+H`$, so that $`()`$ can be rewritten as
$`()`$ $`2[N^{}]^{m2}(15H^2+c_1^2(S)+c_2(S)9Hc_1(S))`$
Finally, since $`[N^{}]_{|F}^{m2}=1`$, where $`F`$ is a fibre of ยฟ$`\pi :NS`$, we conclude
$`(DF)`$ $`deg\mathrm{\Sigma }_Sdegp_{|Y_S}=2(15H^2+c_1^2(S)+c_2(S)9Hc_1(S))=2(15d+c_1^2(S)+c_2(S)9Hc_1(S))`$,
where $`d=deg(S)`$.
By Noetherโs formula, we can also write
$`(DF^{})`$ $`deg\mathrm{\Sigma }_Sdegp_{|Y_S}=2(15d+12\chi (๐ช_S)9Hc_1(S))`$.
We can express also our formula in terms of the sectional genus $`\pi `$ of our surface $`S`$ ( recall that $`2\pi 2=H^2Hc_1(S)`$) as
$`(DF^{\prime \prime })`$ $`deg\mathrm{\Sigma }_Sdegp_{|Y_S}=2(18(\pi 1)+6d+12\chi (๐ช_S))`$
###### Example 4
For $`m=3`$, we have $`c_1(S)=(4d)H`$ and
$`c_2(S)=6H^2d(4d)H^2`$. Thus
$`deg\mathrm{\Sigma }_Sdegp_{|Y_S}=2d(d1)(2d1)`$
###### Example 5
For $`m=4`$, we have the formula $`c_2(S)=c_1(S)^25Hc_1(S)+10dd^2`$ $`[Hartshorne,p.434]`$, or , equivalently,
$`d^25d+2(6\chi (๐ช_S)c_1(S)^2)=10(\pi 1)`$
which gives
$`deg\mathrm{\Sigma }_Sdegp_{|Y_S}=2/5(9d^215d+168\chi (๐ช_S)18c_1(S)^2)`$
## 6 Weakly focally degenerate varieties
In this section we shall first consider the case of a hypersurface $`X`$ of dimension $`n`$ , and we shall characterize the case where $`X`$ is weakly focally degenerate. The characterization of the hypersurfaces $`X`$ which are isotropically focally degenerate will be given in the next section.
Later on in this section we shall deal with focally degenerate varieties of any codimension.
We shall essentially use very classical tools such as the implicit function theorem, dimension counts and the standard method of obtaining new equations by differentiating old ones .
Let $`F(x_1,\mathrm{}x_{n+1})=0`$ be the affine polynomial equation of a hypersurface $`X`$. We shall in this section be mostly interested about a birational description of $`X`$, whenceforth we might, by abuse of notation, not distinguish between a projective variety and its affine part (or any nonempty Zariski open set of it).
In this case the gradient $`F`$ of $`F`$ gives a trivialization of the Normal Bundle $`N_X`$ at the smooth points of $`X`$, and the second projection $`\pi :N_X๐^{๐ง+\mathrm{๐}}`$ coincides with the endpoint map
$`ฯต(x,\lambda )=x+\lambda F(x)`$,
where $`x=(x_1,\mathrm{}x_{n+1})`$ is a point of $`X`$, and $`\lambda `$ is a scalar coordinate $`=\lambda _1/\lambda _0`$ , $`(\lambda _0,\lambda _1)`$ being homogeneous coordinates on $`๐^\mathrm{๐}`$.
As a warm up, let us investigate when does it occur that the endpoint map is not finite. That is, let us assume that $`\mathrm{\Gamma }`$ is a curve in $`N_X`$ which is mapped to a point $`O`$ by the endpoint map $`ฯต`$, and that this point does not lie at infinity.
Choosing a parameter $`t`$ for $`\mathrm{\Gamma }`$, we have functions $`x(t),\lambda (t)`$ such that
1) $`F(x(t))0`$
1โ)$`x(t)+\lambda (t)F(x(t))O`$.
If $`x(t)`$ is a smooth point of $`X`$, then the gradient $`F(x(t))`$ does not vanish, whence $`x(t)`$ is not constant: thus at a general point of $`\mathrm{\Gamma }`$ we may assume that the derivative $`\dot{x}(t):=dx(t)/dt`$ is non vanishing.
Let us use the scalar product $`<,>`$ associated to the quadratic form $`Q_{\mathrm{}}`$, and let us choose affine coordinates such that $`<,>`$ is the standard scalar product ( i.e., the matrix of $`Q_{\mathrm{}}`$ is the identity matrix); since
2) $`x(t)O\lambda (t)F(x(t))`$ , and $`<F(x(t)),\dot{x}(t)>0`$ we infer that
3) $`<x(t)O,x(t)O>constant`$.
Therefore, the basis curve $`\gamma X=\{x|F(x)=0\}`$ is a curve contained in a sphere with centre the point $`O`$ ( note that the sphere may also have radius zero !).
Conversely, if we have such a spherical curve $`\gamma `$ meeting $`X`$ and with the property that the two vectors $`x(t)O,F(x(t))`$ are proportional, then we find $`\lambda (t)`$ so that 1โ), 1) hold, whence we find $`\mathrm{\Gamma }`$ which is mapped to the point $`O`$ by the endpoint map ( and moreover it follows from 1) that $`\gamma `$ is contained in $`X`$). Finally, since $`\mathrm{\Gamma }`$ is mapped to a point, it is obviously contained in the ramification divisor $`Y`$ of the endpoint map.
We have therefore the following
###### Proposition 2
Given a smooth affine hypersurface $`X`$ , the positive dimensional irreducible components of the fibres $`\pi ^1(O)`$ of the map to the affine part of the Focal Locus correspond exactly to the subvarieties $`\mathrm{\Phi }`$ contained in a sphere $`S`$ with centre $`O`$, and such that $`X`$ is everywhere tangent to $`S`$ along $`\mathrm{\Phi }`$.
Proof
Let $`\mathrm{\Psi }`$ be a component of the fibre $`\pi ^1(O)`$. Then consider that $`\mathrm{\Psi }`$ is the union of the curves $`\mathrm{\Gamma }`$ contained in it : each of these projects to $`\gamma X`$ contained in a sphere $`S_c`$ with centre $`O`$ and radius $`c`$. But the image of $`\mathrm{\Psi }`$, call it $`\mathrm{\Phi }`$, is irreducible, whence all the radii are equal, and we get the desired sphere $`S`$. Conversely, the tangency condition provides a rational function $`\lambda `$ on $`\mathrm{\Phi }`$ whose graph is the required variety $`\mathrm{\Psi }`$. $`\mathrm{}`$
It is rather clear that the previous proposition allows easily to construct examples where the map $`\pi :Y\mathrm{\Sigma }_X`$ is not finite.
###### Remark 8
If instead the point $`O`$ is at infinity, letโs identify it with one vector in $`V^{}`$, then we get the equation
$`O\lambda (t)F(x(t))`$, whence
$`<O,x(t)>constant`$ .
So, in this case, the positive dimensional irreducible components of the fibres $`\pi ^1(O)`$ correspond exactly to the subvarieties $`\mathrm{\Phi }`$ contained in a hyperplane $`H`$ with normal direction $`O`$, and such that $`X`$ is everywhere tangent to $`H`$ along $`\mathrm{\Phi }`$.
We can push the previous calculations to describe the weakly focally degenerate hypersurfaces.
Let us thus assume that $`X=\{x|F(x)=0\}`$ is weakly focally degenerate. This simply means that there is a component $`\mathrm{\Sigma }`$ of the focal locus which has dimension
$`dim\mathrm{\Sigma }=a<n`$.
Arguing as before, we notice that $`\mathrm{\Sigma }`$ will simply be any maximal irreducible variety such that its inverse image in $`N_X`$ has a dominating component $`Z`$ of dimension $`n`$. We can analogously treat the case where this dimension is bigger than $`n`$, i.e., when $`Z`$ = $`N_X`$, or equivalently $`X`$ is isotropically focally degenerate : in this case we may also allow $`dim\mathrm{\Sigma }=n`$.
We have thus an irreducible component $`Z`$ of the ramification divisor, with $`\pi (Z)=\mathrm{\Sigma }`$.
To start with, let us assume that $`\mathrm{\Sigma }๐_{\mathrm{}}`$.
Therefore, at the general point of $`Z`$ we can choose local coordinates
$`s=(s_1,\mathrm{}.,s_a)`$ and $`t=(t_1,..t_{\nu a})`$ ( $`\nu =n`$ or $`n+1`$)
such that the fibres of $`\pi `$ are locally given by setting $`s=constant`$, in other words we have functions
$`x(s,t),\lambda (s,t)`$ parametrizing the points of $`Z`$ ,
and a function $`O(s)`$ parametrizing the image $`\pi (Z)=\mathrm{\Sigma }`$ of the end-point map. This means that the following equations hold :
1โ) $`F(x(s,t))0`$
2โ) $`x(s,t)O(s)\lambda (s,t)F(x(s,t))`$ ,
differentiating 1โ) with respect to both sets of variables $`s,t`$, we infer that
$`<F(x(s,t)),(dx(s,t)/dt_i)>0`$ as well as
$`<F(x(s,t)),(dx(s,t)/ds_j)>0`$.
We argue as we did before :
since $`x(s,t)O(s)`$ is proportional to $`F(x(s,t))`$, we obtain that $`x(s,t)O(s)`$ is orthogonal to all the partial derivatives of $`x(s,t)`$.
Since however $`(dx(s,t)/dt_i)=(d(x(s,t)O(s))/dt_i)`$, it follows that there is a function $`r(s)`$ such that
3โ) $`<x(s,t)O(s),x(s,t)O(s)>r(s)`$.
What we have done insofar is to write down the family of spheres containing the projections $`X_s`$ to $`X`$ of the fibres over $`O(s)\mathrm{\Sigma }`$.
On the other hand, we can use the other partial derivatives $`(dx(s,t)/ds_j)`$ in order to obtain a complete description of $`X_s`$.
In fact, let us calculate the partial derivatives $`(r(s)/s_j)`$
They are $`=2<x(s,t)O(s),((x(s,t)O(s)/s_j)>`$
$`=2<x(s,t)O(s),(O(s)/s_j)>`$.
We have therefore established
4โ) $`(r(s)/s_j)=2<x(s,t)O(s),(O(s)/s_j)>`$,
whose geometric meaning is the following: if $`O(s)`$ is a smooth point of $`\mathrm{\Sigma }`$, whence all the partial derivatives $`(dO(s)/ds_j)`$ are linearly independent, then $`X_s`$ is contained in the intersection of the sphere given by 3โ) with the codimension $`a`$ affine subspace given by 4โ).
If this intersection has the expected dimension $`na`$, then it has the same dimension as $`X_s`$ and if it is moreover irreducible it will coincide with $`X_s`$.
###### Lemma 3
Consider an affine subspace $`L=\{x|<xO,v_j>=c_j`$ for $`j=1,..a\}`$ of codimension $`a`$ and assume that $`L`$ is contained in the sphere $`S(O,r^{1/2})=\{x|<xO,xO>=r\}`$. Then
(\*) the direction $`W`$ of $`L`$ is an isotropic subspace for $`<,>`$, and there exists $`x_0L`$ such that $`x_0O`$ is orthogonal to $`W`$ ( equivalently, $`W`$ is isotropic and $`LO+W^{}`$).
Observe moreover that the orthogonal $`W^{}`$ is the vector space $`U`$ generated by the vectors $`v_j`$.
Also the converse holds, in the sense that if (\*) is verified, then there exists a constant $`R`$ such that $`L`$ is contained in the sphere $`S(O,R^{1/2})`$.
Proof
Let $`x_0L`$ and write $`L=x_0+W`$. Since $`<xO,xO>r`$ for $`xL`$, we get
$`<x_0O,x_0O>+2<w,x_0O>+<w,w>r`$ for each vector $`wW`$.
Thus the quadratic polynomial $`<w,w>`$ is identically zero on $`W`$ , what amounts to say that the subspace $`W`$ is isotropic; the vanishing of the linear form yields the desired orthogonality of $`x_0O`$ to $`W`$.
Conversely , $`<xO,xO><x_0O,x_0O>`$ and $`L`$ is contained in the sphere
$`\{x|<xO,xO>=R\}`$ once we set $`R=<x_0O,x_0O>`$. $`\mathrm{}`$
###### Lemma 4
Consider an affine subspace $`L=\{x|<xO,v_j>=c_j\}`$ as in the previous lemma $`3`$, and assume that the affine quadric $`LS(O,r^{1/2})`$ is reducible. Then either
(i) $`dim(W/WW^{})=1`$ and there exists $`x_0L`$ such that $`x_0O`$ is orthogonal to $`W`$ and $`<x_0O,x_0O>r`$ or
(ii) $`dim(W/WW^{})=2`$ and there exists $`x_0L`$ such that $`x_0OW^{}`$ and $`<x_0O,x_0O>=r`$
Proof
As before , for each choice of $`x_0L`$ we can write $`L=x_0+W`$. Since the equation of our affine quadric is
$`<x_0O,x_0O>+2<w,x_0O>+<w,w>=r`$ for each vector $`wW`$, and we impose the condition that the quadric be the union of two affine hyperplanes, it follows that the quadratic form $`<w,w>`$ on $`W`$ has rank either $`1`$ or $`2`$.
In the latter case, since the rank of the complete quadric equals the rank of the quadratic form, acting with a translation on $`W`$, we can kill the terms of lower degree.
In the former case, if the linear part of the equation would not belong to the image under $`Q_{\mathrm{}}`$ of $`W/WW^{}`$, the rank would be at least $`3`$. Whence, acting with a translation on $`W`$, we may kill the linear part and then the constant must be non zero.
$`\mathrm{}`$
We have therefore found that the projection of $`Z`$ is contained in the locus $`X^{}`$ given by
3โโ) $`\{x|s,<xO(s),xO(s)>r(s)`$
4โโ) $`(dr(s)/ds_j)=2<xO(s),(O(s)/s_j)>\}`$.
If moreover $`Z`$ surjects onto $`X`$ and $`X^{}`$ is irreducible, then $`X^{}`$ equals $`X`$ unless we are in the exceptional case where (cf. Lemma 3) for each point $`O(s)`$ the (vector) tangent space $`V_s`$ to $`\mathrm{\Sigma }`$ at $`O(s)`$ satisfies the condition that $`V_s`$ contains its orthogonal $`W_s:=V_s^{}`$, and moreover then $`(x(s,t)O(s))`$, for each $`t`$ belongs to the subspace $`V_s:=W_s^{}`$.
The locus $`X^{}`$, as written, is the projection of the locus
$`Z^{}๐^m\times \mathrm{\Sigma }`$ defined as
3โโ) $`\{(x,s)|<xO(s),xO(s)>r(s)`$
4โโ) $`(r(s)/ds_j)=2<xO(s),(O(s)/s_j)>\}`$.
If we calculate the tangent space to $`Z^{}`$ at the point $`(x,s)`$ we obtain that it is contained in the hyperplane:
5โโ) $`\{(\xi ,\sigma )|2<xO(s),\xi >2<xO(s),(O(s)/\sigma )>(r(s)/\sigma )=0\}`$ = $`\{(\xi ,\sigma )|2<xO(s),\xi >=0\}`$
(since $`(x,s)`$ is a point of $`Z^{}`$).
By Sardโs lemma, $`X^{}`$ has dimension at most $`n`$: whence, if we assume that the component $`Z`$ dominates $`X`$, and thus $`XX^{}`$, we conclude that $`X=X^{}`$ ( in the exceptional case, or if $`X^{}`$ is irreducible) or at least that $`X`$ is a component of $`X^{}`$.
We are now in the position to explain the main constructions which are underlying the characterization of the focally degenerate varieties.
###### Definition 6
THE INVERSE CONSTRUCTION TO FOCAL DEGENERACY.
Start from the following data :
i) Let $`\mathrm{\Sigma }`$ be an irreducible affine variety of dimension $`a`$ , and let $`\mathrm{\Sigma }^{}`$ be an irreducible subvariety of the product $`\mathrm{\Sigma }\times ๐`$ which is the graph of an algebraic function $`r`$ on $`\mathrm{\Sigma }`$.
Proceed constructing an algebraic set $`X^{}`$ as follows:
ii) The subvariety $`\mathrm{\Sigma }^{}`$ defines a family of spheres
$`Z^{}๐^m\times \mathrm{\Sigma }\times ๐`$ defined as
3โโ) $`\{(x,O,r)|(O,r)\mathrm{\Sigma }^{},<xO,xO>=r\}`$.
iii) Consider in $`๐^m\times ๐`$ the tangent space to $`\mathrm{\Sigma }^{}`$ at a point $`(O,r)`$, and its orthogonal with respect to the quadratic form $`Q_{\mathrm{}}1`$: under the embedding of $`๐^m`$ in $`๐^m\times ๐`$ sending $`x`$ to $`(xO,1/2)`$, its pull back is precisely an affine space given by an equation as 4โโ). We can in this way define a bundle (if $`\mathrm{\Sigma }`$ is smooth) of affine spaces
$`A^{}๐^m\times \mathrm{\Sigma }^{}`$,
$`A^{}=\{(x,O,r)|(O,r)\mathrm{\Sigma }^{},(xO,1/2)T\mathrm{\Sigma }_{(O,r)}^{}{}_{}{}^{}\}`$.
iv) define $`Z^{}`$ as the intersection $`Z^{}A^{}`$ ( a divisor in $`A^{}`$)
v) define $`X^{}`$ as the projection of $`Z^{}`$ on the first factor $`๐^m`$;
vi) observe that, by the argument we gave above, $`dimX^{}m1`$.
vii) assume finally that $`\mathrm{\Sigma },r`$ are $`\mathrm{๐๐๐ฆ๐ข๐ฌ๐ฌ๐ข๐๐ฅ๐}`$, which amounts to the requirement that $`Z^{}`$ dominate $`\mathrm{\Sigma }`$.
###### Remark 9
The condition that $`\mathrm{\Sigma },r`$ be $`\mathrm{๐๐๐ฆ๐ข๐ฌ๐ฌ๐ข๐๐ฅ๐}`$ is obviously satisfied unless $`Z^{}`$ is a union of fibres of the projection $`A^{}\mathrm{\Sigma }`$. This means, unless the quadratic function $`<xO(s),xO(s)>`$ is constant on the affine spaces $`A_s^{}`$. Therefore, the pair $`\mathrm{\Sigma },r`$ is $`\mathrm{๐๐๐ฆ๐ข๐ฌ๐ฌ๐ข๐๐ฅ๐}`$ unless we are in the situation of Lemma 3 , whence $`<xO(s),xO(s)>R(s)`$ on $`A_s^{}`$, but $`R(s)r(s)`$.
There remains however to see what happens in the case where $`\mathrm{\Sigma }`$ lies at infinity .
In this case, we derive (cf. remark 8) the following equations, where $`O(s)`$ is a $`V^{}`$ -valued function leading to a parametrization of $`\mathrm{\Sigma }`$ :
6) $`<x(s,t),O(s)>r(s)`$
7) $`(r(s)/s_j)<x(s,t),(O(s)/s_j)>,forj=1,\mathrm{}a`$.
In this case, if $`O(s)`$ is a smooth point of $`\mathrm{\Sigma }`$, then the $`a+1`$ vectors $`O(s),O(s)/s_j`$ are linearly independent and 6) and 7) imply that $`X_s`$ is contained in the affine space
8) $`X_s^{}=\{x|<x,O(s)>r(s),(r(s)/s_j)<x,(O(s)/s_j)>,forj=1,\mathrm{}a\}`$.
Since $`\mathrm{\Sigma }`$ lies at infinity , $`X`$ is not isotropically focally degenerate, whence $`Z`$ has dimension $`m1`$ : it follows that $`X_s`$, $`X_s^{}`$ have the same dimension $`m1a`$, whence they coincide.
Moreover, $`Z`$ must dominate $`X`$, else a whole fibre of $`N_XX`$ is contained in $`Z`$, and therefore its projection cannot lie at infinity (remember that $`X`$ is here supposed to be affine).
Therefore, it follows that $`X`$ equals $`X^{}`$, the closure of the union of the $`X_s^{}`$.
We are therefore led to the following
###### Definition 7
THE ASYMPTOTIC INVERSE CONSTRUCTION TO WEAK FOCAL DEGENERACY.
Start from the following data :
i) Let $`\mathrm{\Sigma }`$ be an irreducible variety of dimension $`a`$ , contained in $`๐_{\mathrm{}}`$, and let $`\mathrm{\Sigma }^{}`$ be an irreducible subvariety of the product $`\mathrm{\Sigma }\times ๐`$ which is the graph of an algebraic section $`r`$ of $`๐ช_\mathrm{\Sigma }(1)`$.
Then consider the algebraic set $`X^{}`$ which is the closure of the union of the family of affine spaces $`X_s^{}`$ defined by 8).
###### Remark 10
The attentive reader will find a slight abuse of notation above, which can be explained as follows : in the case where $`\mathrm{\Sigma }`$ does not lie at infinity, since we have a privileged affine chart = $`๐^๐ฆ๐_{\mathrm{}}`$, we consider $`r(s)`$ just as an algebraic function on $`\mathrm{\Sigma }`$. If however $`\mathrm{\Sigma }๐_{\mathrm{}}`$, then there is no favourite standard affine chart and we make clear that $`r`$ is not really a function, but a section of $`๐ช_\mathrm{\Sigma }(1)`$ (possibly multivalued and with poles!).
###### Remark 11
Consider a variety $`X=X^{}`$ obtained from the asymptotic inverse focal construction.
Then its part $`X_{\mathrm{}}=X๐_{\mathrm{}}`$ consists of the points
$`\{x๐_{\mathrm{}}|<x,O(s)>r(s),(r(s)/s_j)=<x,(O(s)/s_j)>,forj=1,\mathrm{}a\}`$.
If we therefore identify $`๐_{\mathrm{}}`$ with its dual space via the quadratic form $`Q_{\mathrm{}}`$, it follows that $`X_{\mathrm{}}`$ is the dual variety of $`\mathrm{\Sigma }`$ !
Observe moreover, that if $`X`$ is a linear space, then $`\mathrm{\Sigma }_X`$ equals $`X_{\mathrm{}}^{}`$.
In this case the section $`r(s)`$ is just induced by a linear form on $`\mathrm{\Sigma }_X`$ (i.e., a vector in $`(V^{})^{}`$).
We observe now that we have insofar proved the following
###### Theorem 2
Let $`X`$ be a focally degenerate hypersurface in $`๐^{n+1}`$ and let $`\mathrm{\Sigma }`$ be a component of the strict focal locus ( i.e., we are in the non vertical case and the corresponding component $`Z`$ of $`Y_X`$ projects onto $`X`$). Then
* either $`\mathrm{\Sigma }`$ is contained in $`๐_{\mathrm{}}`$ and $`X`$ is obtained from $`\mathrm{\Sigma },r`$ via the asymptotic inverse focal construction associated to an algebraic section $`r`$ of $`๐ช_\mathrm{\Sigma }(1)`$
* or, $`\mathrm{\Sigma }`$ is not contained in $`๐_{\mathrm{}}`$ and there is an algebraic function $`r(s)`$ on $`\mathrm{\Sigma }`$ such that, applying the inverse construction to focal degeneracy, we get a hypersurface $`X^{}`$ such that $`X`$ is a component of $`X^{}`$.
Conversely, start from any admissible pair $`(\mathrm{\Sigma },r)`$, and assume that an irreducible hypersurface is a component of the algebraic set $`X^{}`$ obtained from the inverse construction or from the asymptotic inverse construction : then $`X`$ is a focally degenerate hypersurface.
Proof
There remains only to show that if $`X`$ is an irreducible hypersurface, component of the algebraic set $`X^{}`$ obtained from an inverse construction : then $`\mathrm{\Sigma }`$ is a component of the focal locus of $`X`$. This follows since, by 5โโ), $`xO(s)`$ is a normal vector to $`X^{}`$, respectively since $`O(s)`$ is a normal vector to $`X^{}`$; moreover, $`Z^{}`$ dominates $`\mathrm{\Sigma }`$ by the assumption that $`r`$ be admissible. $`\mathrm{}`$
However, the inverse constructions, as we are going to see, work more generally also in the case where $`X^{}`$ has smaller dimension than the expected dimension $`m1`$.
We have in fact the following
###### Theorem 3
Let $`X`$ be a focally degenerate variety of dimension $`n`$ in $`๐^m`$ and let $`\mathrm{\Sigma }`$ be a component of the focal locus of dimension $`am1`$, projection of a component $`Z`$ of $`Y_X`$. Then $`\mathrm{\Sigma }`$ determines birationally an irreducible subvariety $`\mathrm{\Sigma }^{}`$ of $`\mathrm{\Sigma }\times ๐`$ corresponding to an algebraic section $`r(s)`$ of $`๐ช_\mathrm{\Sigma }(1)`$ and, applying the appropriate inverse construction to focal degeneracy, we get an algebraic set $`X^{}`$ which is focally degenerate, and indeed isotropically focally degenerate in the case where $`\mathrm{\Sigma }`$ is not contained in the hyperplane at infinity $`๐_{\mathrm{}}`$ and $`dimZ^{}=m`$ (in this case the fibres $`X_s^{}`$ of $`Z^{}\mathrm{\Sigma }`$ are affine spaces).
There are seven cases :
* 1) $`X`$ is isotropically focally degenerate : then $`X=X^{}`$, $`dimZ^{}=m`$ and the fibres $`X_s`$ of $`N_X\mathrm{\Sigma }`$ are affine spaces. Moreover, here $`\mathrm{\Sigma }`$ is not contained in $`๐_{\mathrm{}}`$.
* 2) $`\mathrm{\Sigma }`$ is not contained in $`๐_{\mathrm{}}`$ , $`Z`$ projects onto $`X`$ and $`X^{}`$ is not isotropically focally degenerate: then $`X`$ is a component of $`X^{}`$
* 3) $`Z`$ projects onto $`X`$, $`X^{}`$ is isotropically focally degenerate, but $`X`$ is not isotropically focally degenerate: then $`XX^{}`$ is a divisor, $`Z`$ is the restriction to $`X`$ of the normal bundle $`N_X^{}`$, and $`\mathrm{\Sigma }`$ is the focal locus of $`X^{}`$ (again here $`\mathrm{\Sigma }`$ is not contained in $`๐_{\mathrm{}}`$)
* 4) $`\mathrm{\Sigma }`$ is not contained in $`๐_{\mathrm{}}`$ , $`Z`$ projects onto a divisor $`X\mathrm{"}X`$, $`X\mathrm{"}`$ is a component of $`X^{}`$, $`X\mathrm{"}`$ is focally degenerate, with a component $`Z\mathrm{"}`$ of the ramification locus $`Y_{X\mathrm{"}}`$ which is a subbundle of $`N_{X\mathrm{"}}`$ : then the tangent bundle to $`X`$ around $`X\mathrm{"}`$ is annihilated by the given subbundle $`Z\mathrm{"}`$.
* 5) $`Z`$ projects onto a divisor $`X\mathrm{"}X`$ which is focally degenerate, $`X\mathrm{"}`$ is a divisor of $`X^{}`$ and $`X^{}`$ is isotropically focally degenerate (again here $`\mathrm{\Sigma }`$ is not contained in $`๐_{\mathrm{}}`$). Then $`X`$ and $`X^{}`$ are tangent along $`X\mathrm{"}`$.
* 6) $`\mathrm{\Sigma }`$ is contained in $`๐_{\mathrm{}}`$, $`Z`$ projects onto some affine point of $`X`$ , whence it dominates $`X`$ and $`X=X^{}`$ is obtained via the asymptotic inverse focal construction.
* 7) $`\mathrm{\Sigma }`$ is contained in $`๐_{\mathrm{}}`$, $`Z`$ projects onto a component $`\mathrm{\Delta }`$ of $`X_{\mathrm{}}`$. In this case $`Z`$ is the restriction of $`N_X`$ to $`\mathrm{\Delta }`$, the second projection to $`๐_{\mathrm{}}`$ is not surjective. This case is characterized by the property that $`\mathrm{\Delta }X_{\mathrm{}}`$ be projectively isotropically degenerate, which is equivalent to the property that $`\mathrm{\Delta }`$ be obtained via the isotropic projective inverse focal construction (this case will be treated separately in the next proposition).
Conversely, start from any admissible pair of a variety $`\mathrm{\Sigma }`$ not contained in $`๐_{\mathrm{}}`$, and of an algebraic section $`r(s)`$. Consider the algebraic set $`X^{}`$ obtained from the inverse construction : then $`X^{}`$ is focally degenerate ( if it has two components, this means that each of them is focally degenerate) and isotropically focally degenerate iff the fibres of $`Z^{}\mathrm{\Sigma }`$ are affine spaces of dimension $`mdim\mathrm{\Sigma }`$ (then $`Z^{}=N_X^{}`$) .
All the isotropically focally degenerate varieties $`X`$ are gotten by the inverse construction as such an $`X^{}`$.
In case 6), where $`\mathrm{\Sigma }`$ is a component of the strict focal locus contained in $`๐_{\mathrm{}}`$, all such weakly focally degenerate varieties are obtained from the asymptotic inverse focal construction.
Let us consider the remaining cases where $`\mathrm{\Sigma }`$ is not contained in $`๐_{\mathrm{}}`$.
Then the weakly focally degenerate varieties in the non vertical case (i.e., when $`Z`$ dominates $`X`$) are gotten either
(i) as a component of such an $`X^{}`$, or
(ii) as a divisor in an isotropically focally degenerate variety $`X^{}`$, which is transversal to the general fibres $`X_s^{}`$ of $`N_X^{}\mathrm{\Sigma }_X^{}`$ and where $`dim\mathrm{\Sigma }_X^{}m2`$.
Instead, in the vertical case, the weakly focally degenerate varieties are given as varieties containing a focally degenerate divisor $`X\mathrm{"}`$ such that either
(i) $`X\mathrm{"}`$ is a component of $`X^{}`$, with a component $`Z\mathrm{"}`$ of the ramification locus $`Y_{X\mathrm{"}}`$ which is a subbundle of $`N_{X\mathrm{"}}`$, and such that the tangent bundle of $`X`$ along $`X\mathrm{"}`$ is given by the annihilator of the subbundle $`Z\mathrm{"}`$ or
(ii) $`X\mathrm{"}`$ is a divisor of $`X^{}`$, $`X^{}`$ is isotropically focally degenerate with $`dim\mathrm{\Sigma }_X^{}m2`$, $`X\mathrm{"}`$ is transversal to the fibres $`X_s^{}`$ of $`N_X^{}\mathrm{\Sigma }_X^{}`$, and $`X`$ and $`X^{}`$ are tangent along $`X\mathrm{"}`$.
Proof
We discuss first of all the case where $`\mathrm{\Sigma }`$ is not contained in $`๐_{\mathrm{}}`$ (whence $`Z`$ does not project to $`๐_{\mathrm{}}`$ under the first projection) .
Around each smooth point of $`X`$ there are a Zariski open set $`U`$ of $`๐^m`$ and polynomials $`F_1(x),\mathrm{}.F_{mn}(x)`$ such that $`XU`$ = $`\{xU|F_1(x)=\mathrm{}.F_{mn}(x)=0\}`$ and such that $`XU`$ consists of smooth points.
Therefore, the gradients of the polynomials $`F_1(x),\mathrm{}.F_{mn}(x)`$ yield a framing of the Euclidean normal bundle on $`XU`$, and the endpoint map is locally given by
$`ฯต(x,\lambda _1,..\lambda _{mn})=x+\mathrm{\Sigma }_{i=1,..mn}\lambda _iF_i(x)`$.
We choose as we did before a component $`Z`$ of the ramification locus $`Y_X`$ which maps onto an irreducible variety $`\mathrm{\Sigma }`$ of dimension $`am2`$ (respectively $`am1`$ in the focally isotropically degenerate case) and local coordinates
$`s=(s_1,\mathrm{}.,s_a)`$ for the points of $`\mathrm{\Sigma }`$ and $`t=(t_1,..t_{\nu a})`$ for the fibres of $`\pi `$, where $`\nu `$ equals $`m1`$ in the non focally isotropically degenerate case, otherwise $`m=\nu `$ and $`Z=N_X`$.
Whence, we have local functions $`x(s,t),\lambda (s,t)`$ parametrizing the points of $`Z`$ ,
and a function $`O(s)`$ parametrizing the image $`\pi (Z)=\mathrm{\Sigma }`$ such that
1โ) $`F_i(x(s,t))0`$ $`i`$
2โ) $`x(s,t)O(s)\mathrm{\Sigma }_{i=1,..mn}\lambda _i(s,t)F_i(x(s,t))`$.
Differentiating 1โ) with respect to both sets of variables $`s,t`$, we infer that
$`<F_j(x(s,t)),(x(s,t)/t_i)><F_j(x(s,t)),(x(s,t)/s_h)>0i,j,h`$.
By 2โ) $`x(s,t)O(s)`$ is a normal vector, whence
3โ) $`<x(s,t)O(s),x(s,t)O(s)>r(s)`$.
and
4โ) $`(r(s)/s_j)=2<x(s,t)O(s),(O(s)/s_j)>`$.
Therefore, for fixed $`s`$, the projection $`X_s`$ of the fibre $`Z_s`$ (generally a manifold of dimension $`\nu a`$) is contained in the intersection $`X_s^{}`$ of a sphere $`S_s`$ of centre $`O(s)`$ and radius $`r(s)^{1/2}`$ with an affine space $`\mathrm{\Pi }_s`$ of codimension $`a`$ (since at the general point we can assume $`O(s)/s_1),\mathrm{}O(s)/s_a)`$ to be linearly independent).
Thus, the manifold $`X_s^{}`$ has dimension either $`m1a`$ or $`ma`$ ( but in the latter case, by Lemma 3, the orthogonal to $`T\mathrm{\Sigma }_{O(s)}`$ is contained in $`T\mathrm{\Sigma }_{O(s)})`$.
Consider as before the locus $`X^{}`$ given as the projection of the locus
$`Z^{}๐^m\times \mathrm{\Sigma }`$ defined as
3โโ) $`\{(x,s)|<xO(s),xO(s)>r(s)`$
4โโ) $`(r(s)/s_j)=2<xO(s),(O(s)/s_j)>\}`$.
###### Lemma 5
$`Z^{}N_X^{}`$
Proof
We must prove that the vector $`xO(s)`$ is normal to $`X^{}`$. This follows ยฟfrom the calculation of the tangent space to $`Z^{}`$ at the point $`(x,s)`$ that we have done above (cfr. 5โโ). $`\mathrm{}`$
###### Corollary 5
Each component of $`X^{}`$ is focally degenerate and indeed isotropically focally degenerate iff $`X_s^{}=\mathrm{\Pi }_s`$ (whence, in the latter case, $`X^{}`$ is also irreducible).
Proof
If $`X_s^{}=\mathrm{\Pi }_s`$, then $`Z^{}`$ is irreducible and $`dimZ^{}=dimN_X^{}=m`$ so that $`Z^{}=N_X^{}`$ and $`X^{}`$ is irreducible and isotropically focally degenerate.
If $`dimX_s^{}=m1a`$ (in this case $`am2`$), then $`Z^{}`$ is a divisor in $`N_X^{}`$ and hence $`\mathrm{\Sigma }`$ is contained in $`\mathrm{\Sigma }_X^{}`$. Either $`\mathrm{\Sigma }`$ is a component of $`\mathrm{\Sigma }_{X^0}`$, for each component $`X^0`$ of $`X^{}`$, or there is a component $`X^0`$ of $`X^{}`$ which is focally isotropically degenerate.
Assume that the latter holds: then, for general $`O(s)\mathrm{\Sigma }`$, $`X_s^{}`$ is a divisor of the fibre of $`N_{X^0}\mathrm{\Sigma }_{X^0}`$, whence by dimension reasons $`\mathrm{\Sigma }=\mathrm{\Sigma }_{X^0}`$.
Since the direction of $`\mathrm{\Pi }_s`$ is the vector subspace $`W=T\mathrm{\Sigma }_{O(s)}^{}`$, and $`\mathrm{\Sigma }=\mathrm{\Sigma }_{X^0}`$, it follows that $`N_{X_s^0}=\mathrm{\Pi }_s`$.
Moreover, being $`X^0`$ isotropically focally degenerate, by lemma $`3`$ follows that $`W`$ is totally isotropic, whence the quadratic function $`<xO(s),xO(s)>`$ is then constant on $`\mathrm{\Pi }_s`$, contradicting the fact that for general $`s`$ $`X_s^{}`$ is a nonempty and proper divisor in $`\mathrm{\Pi }_s`$. $`\mathrm{}`$
If $`X`$ is focally isotropically degenerate, the projection $`X_s`$ of $`Z_s`$ has dimension $`ma`$, whence it equals $`X_s^{}`$, and it follows immediately that $`X`$ equals $`X^{}`$.
Suppose then that $`X`$ is not isotropically focally degenerate, and let $`X^{\prime \prime }`$ be the projection of $`Z`$, that is the closure of $`_sX_s`$. Thus $`X\mathrm{"}X`$ and $`X\mathrm{"}X^{}`$.
Assume first that $`dimX_s^{}=dimX_s=m1a`$. Therefore, $`X\mathrm{"}`$ is a component of $`X^{}`$ and $`Z`$ equals a component $`Z\mathrm{"}`$ of $`Z^{}`$. It follows that either $`X\mathrm{"}=X`$ and case 2) of the theorem occurs, or $`X\mathrm{"}X`$ would be a divisor and $`Z`$ would be the restriction to $`X\mathrm{"}`$ of the normal bundle $`N_X`$, a subbundle of the normal bundle $`N_{X\mathrm{"}}`$.
Whence, $`X\mathrm{"}`$ is focally degenerate, with a component $`Z=Z\mathrm{"}`$ of the ramification locus which is a projective subbundle of ยฟ$`N_{X\mathrm{"}}`$, and case 4) occurs. Any variety $`M`$ containing $`X\mathrm{"}`$ as a divisor, and with tangent bundle annihilated by the given subbundle would be a weakly focally degenerate variety with $`\mathrm{\Sigma }`$ in the focal locus.
In other words, in the vertical case, the inverse focal construction can by no means reconstruct $`X`$, but only the first order neighbourhood of $`X`$ along $`X\mathrm{"}`$.
Finally, there remains the case where $`dimX_s=m1a,dimX_s^{}=ma`$, in which case $`X^{}`$ is isotropically focally degenerate. Then $`Z`$ is a divisor in $`Z^{}=N_X^{}`$.
Assume $`X\mathrm{"}=X^{}`$ : since then $`X^{}X`$, but $`X^{}X`$ since $`X`$ is not isotropically focally degenerate, we get that $`Z=N_X|_X^{}Z^{}=N_X^{}`$, and we are again in case 4).
Thus we may consider the remaining cases where $`X\mathrm{"}`$ is a divisor of $`X^{}`$. Furthermore, either $`X\mathrm{"}=X`$ or $`X\mathrm{"}`$ is a divisor in $`X`$. If $`X=X\mathrm{"}`$, then $`Z`$ is the restriction of $`N_X^{}`$ to $`X`$, and case 3) occurs. If $`X\mathrm{"}`$ is a divisor of $`X`$, we have that $`Z=N_X|_{X\mathrm{"}}=N_X^{}|_{X\mathrm{"}}`$ so that $`X`$ and $`X^{}`$ are tangent along $`X\mathrm{"}`$ and case 5) occurs. Conversely, let $`X^{}`$ be a isotropically focally degenerate variety and let $`X`$ be a divisor inside $`X^{}`$; since $`N_X^{}|_XN_X`$ is a divisor, it follows immediately that, setting $`Z=N_X^{}|_X`$, the image of $`Z`$ is contained in $`\mathrm{\Sigma }_X^{}`$. If moreover, as it should be, the divisor $`X`$ is transversal to the fibres $`X_s^{}`$, then its image equals $`\mathrm{\Sigma }_X^{}`$, whence $`Z`$ will make $`X`$ weakly focally degenerate if and only if $`dim\mathrm{\Sigma }_X^{}m2`$. More generally, if $`M`$ is any variety containing $`X`$ as a divisor and such that $`M`$ and $`X^{}`$ are tangent along $`X`$, then $`M`$ is weakly focally degenerate.
Let us then consider case 6) : then, analogously to the case of hypersurfaces we can find a parametrization $`O(s)`$ of $`\mathrm{\Sigma }`$ in homogeneous coordinates such that
$`O(s)\mathrm{\Sigma }_{i=1,..mn}\lambda _i(s,t)F_i(x(s,t))`$.
Then $`<x(s,t)/t_i,O(s)><x(s,t)/s_j,O(s)>0`$.
ยฟFrom the first equalities we conclude that there exists a local function $`r(s)`$ such that
6) $`<x(s,t),O(s)>r(s)`$.
One momentโs reflection, since the vector $`O(s)`$ gives homogeneous coordinates for $`\mathrm{\Sigma }`$, shows that indeed $`r(s)`$ globalizes to a (possibly multivalued and with poles) section of $`๐ช_\mathrm{\Sigma }(1)`$.
ยฟFrom the second equalities follows also
7) $`(r(s)/s_j)<x(s,t),(O(s)/s_j)>,forj=1,\mathrm{}a`$.
Thus an entirely similar argument yields that $`X`$ is gotten ยฟfrom the asymptotic inverse focal construction, and conversely if $`X`$ is obtained in this way then $`X`$ is weakly focally degenerate and we are in case 6).
$`\mathrm{}`$
Let us discuss case 7), where the whole condition of degeneracy bears on $`X_{\mathrm{}}`$, and tells that, $`O(s)`$ being the $`V^{}`$\- vector valued function giving local homogeneous coordinates around a smooth point of $`\mathrm{\Sigma }`$ as usual, there is a local function $`\lambda (s,t)`$ and a local parametrization $`x(s,t)`$, of $`X_{\mathrm{}}`$ this time, and giving homogeneous coordinates, such that
$`\lambda (s,t)x(s,t)O(s)`$ is a normal vector to $`X_{\mathrm{}}`$ at the point $`x(s,t)`$, in the sense that
$`<\lambda (s,t)x(s,t)O(s),x(s,t)><x(s,t)/t_i,\lambda (s,t)x(s,t)O(s)><x(s,t)/s_j,\lambda (s,t)x(s,t)O(s)>0`$.
At the points where $`\lambda (s,t)`$ is not vanishing we can replace the parametrization $`x(s,t)`$ by $`\lambda (s,t)x(s,t)`$, so with these new homogeneous coordinates we have
$`I)<x(s,t)O(s),x(s,t)>II)<x(s,t)/t_i,x(s,t)O(s)>III)<x(s,t)/s_j,x(s,t)O(s)>0`$.
Deriving equation I) with respect to $`/t_i`$, and using II) we obtain
$`IV)<x(s,t)/t_i,x(s,t)>0`$
whereas, applying $`/s_j`$ to I) and using III) we get
$`V)<x(s,t)/s_j,x(s,t)><O(s)/s_j,x(s,t)>.`$
IV) yields
$`A)<x(s,t),x(s,t)><O(s),x(s,t)>r(s)`$ which implies, together with V) :
$`B)<O(s)/s_j,x(s,t)>1/2r(s)/s_j`$.
Since we chose a smooth point of $`\mathrm{\Sigma }`$ the $`a+1`$ vectors
$`O(s),O(s)/s_1,\mathrm{}O(s)/s_a`$ are linearly independent, and it follows that the vectors $`x(s,t)`$ , for $`s`$ fixed, vary in an affine space $`X\mathrm{"}_s`$ of dimension $`m1a`$.
Since however $`X_s`$ is assumed to have dimension exactly equal to $`m1a`$, it follows that $`X_s=X\mathrm{"}_s`$, where $`X\mathrm{"}_s`$ is defined by the equations
$`A^{})<O(s),x>r(s)`$
$`B^{})<O(s)/s_j,x>1/2r(s)/s_j`$.
However, also the equality $`<x,x>r(s)`$ must be satisfied on $`X_s=X\mathrm{"}_s`$, thus by Lemma 3 we get an affine linear subspace with direction $`W`$ which is totally isotropic, and is contained in the orthogonal $`W^{}`$ to $`W`$.
The conclusion is that the projective tangent space to $`\mathrm{\Sigma }`$ at any smooth point has a totally isotropic annihilator .
###### Definition 8
Let $`\mathrm{\Sigma }`$ be a projective subvariety of the projective space $`๐(V^{})=๐_{\mathrm{}}`$ associated to a vector space $`V^{}`$ of dimension $`m`$ endowed with a non degenerate quadratic form $`Q_{\mathrm{}}`$, such that any point $`O(s)`$ of $`\mathrm{\Sigma }`$ the projective tangent space to $`\mathrm{\Sigma }`$ at $`O(s)`$ (a vector subspace of $`V^{}`$) has a totally isotropic annihilator.
Let $`r(s)`$ be an algebraic section of $`๐ช_\mathrm{\Sigma }(1)`$ and consider the developable variety $`X\mathrm{"}`$ defined by the union of the subspaces $`X\mathrm{"}_s`$ defined by the equations Aโ) and Bโ).
Assume moreover that $`r`$ be $`\mathrm{๐๐๐ฆ๐ข๐ฌ๐ฌ๐ข๐๐ฅ๐}`$ in the sense that the local function (constant on $`X\mathrm{"}_s`$)
$`<x(s,t),x(s,t)>R(s)`$ be equal to $`r(s)`$.
Then we shall say that $`X\mathrm{"}`$ is projectively isotropically degenerate and that $`X\mathrm{"}`$ is obtained via the isotropic projective inverse focal construction from the admissible pair $`(\mathrm{\Sigma },r)`$.
###### Proposition 3
Assume that $`X`$ is weakly focally degenerate and that a component $`\mathrm{\Sigma }`$ of the focal locus is contained in $`๐_{\mathrm{}}`$, with the corresponding component $`Z`$ of $`Y_X`$ projecting onto a component $`\mathrm{\Delta }`$ of $`X_{\mathrm{}}`$ ( case 7) of theorem 3). In this case $`Z`$ is the restriction of $`N_X`$ to $`\mathrm{\Delta }`$, $`\mathrm{\Delta }`$ is projectively isotropically degenerate. Conversely, if $`\mathrm{\Delta }`$ is obtained via the isotropic projective inverse focal construction , then $`X`$ is weakly focally degenerate and we are in case 7) of theorem 3.
Proof
If $`X`$ is as in case 7) of theorem 3, then we have already seen that $`\mathrm{\Delta }X_{\mathrm{}}`$ is projectively isotropically degenerate.
It remains to prove the converse, which follows since Aโ), Bโ) and our assumption $`R(s)=r(s)`$ imply A), B) by which immediately follow I), II), and III), whence $`x(s,t)O(s)`$ is a normal vector to $`X_{\mathrm{}}`$. Since $`X\mathrm{"}=\mathrm{\Delta }`$ and $`X\mathrm{"}_s`$ has dimension $`m1a`$ we get a component $`Z`$ of dimension $`m1`$ projecting onto the $`a`$-dimensional variety $`\mathrm{\Sigma }`$ contained in the hyperplane at infinity and we are done.
$`\mathrm{}`$
###### Remark 12
It follows from the previous theorem that any variety $`\mathrm{\Sigma }`$ is a component of some focal locus.
Moreover, in the asymptotic inverse construction, we see immediately that the tangent space at a point of $`X_s`$ depends only upon $`s`$ , so that then our $`X`$ is developable.
In particular, if $`m=3`$ we get either a linear subspace or a developable, whence singular, surface.
Observe finally that if $`X`$ is orthogonally general and projective, then only cases 6) or 7) can a priori occur.
For case 6), start choosing $`X_{\mathrm{}}`$ as a smooth and transversal variety to $`Q_{\mathrm{}}`$, apply then the asymptotic inverse focal construction : then we get a variety $`X`$ which will be orthogonally general exactly iff $`X`$ is smooth. But the smoothness of $`X`$, as we have just seen, is the main obstruction.
###### Example 6
Let $`m=3`$, and let $`\mathrm{\Sigma }`$ be the line at infinity parametrized as
$`O(s)=(0,0,1,s)`$, and set, in these affine coordinates, $`r(s)=s^2/2`$.
Then an easy computation for the asymptotic inverse focal construction yields the quadric cone
$`X=\{x|x_0x_2x_3^2/2=0\}`$, whose vertex lies at infinity.
If instead we choose $`O(s)=(0,1,s,s^2)`$, and $`r(s)1`$, $`X`$ will be the quadric cone
$`X=\{x|4(x_1x_0)x_3x_2^2=0\}`$, whose vertex does not lie at infinity.
###### Example 7
Let us now consider the most classical example, namely the rotational torus $`X`$ obtained rotating a circle of radius, say, $`1`$ around the point with coordinates $`(2,0)`$. This is the example of a strongly focally degenerate variety.
The equation $`F`$ of $`X`$, in affine coordinates $`(x,y,z)`$ for which $`Q_{\mathrm{}}`$ yields the standard Euclidean scalar product, is then given, setting
$`q(x,y,z)=(x^2+y^2+z^2+3)`$, or , in homogeneous coordinates $`(x,y,z,w)`$, $`q(x,y,z,w)=(x^2+y^2+z^2+3w^2)`$, by
(\*) $`q^216(x^2+y^2)w^2`$.
The intersection with the plane at infinity is precisely our conic $`Q_{\mathrm{}}=\{q=w=0\}`$, which is a double curve for the quartic surface $`X`$. Moreover, $`Sing(X)`$ consists of $`Q_{\mathrm{}}`$ and of the two points $`\{P,P^{}\}=\{q=x=y=0\}=\{(z^2+3w^2)=x=y=0\}`$.
Now, a classical and easy formula for a rotation surface of a curve $`C`$ =$`r(s),z(s)`$ parametrized by arclength,
$`x(s,\theta )=r(s)cos(\theta )`$
$`y(s,\theta )=r(s)sin(\theta )`$
$`z(s,\theta )=z(s)`$
is that the two principal curvatures equal $`k(s)`$ , $`z^{}(s)/r(s)`$.
In this case, $`r(s),z(s)=(2+cos(s),sin(s))`$, whence $`k1`$ and $`z^{}(s)/r(s)`$ = $`12/r(s)`$.
These formulae are easily rationalized on our surface $`X`$ since $`q^2=16(x^2+y^2)`$, whence $`r=q/4`$. Therefore the critical points are obtained by taking the multiples of the unit normal by the opposites of their inverses, i.e., $`1`$ and $`q/(q8)`$. Finally, the unit normal is obtained by the gradient of $`F`$
$`F=(4x(q8),4y(q8),4qz)`$ upon dividing by its norm, which equals
$`|F|=4((q8)^2(x^2+y^2)+q^2z^2)^{1/2}`$ =
$`(16(q8)^2(x^2+y^2)+16q^2z^2)^{1/2}`$= $`q((q8)^2+16z^2)^{1/2}`$. But since
$`z^2=q3q^2/16`$ , we get $`q(6448)^{1/2}=4q`$,
and the focal locus is obtained for the values $`\lambda =1/4q,\lambda =1/4(q8)`$ as the image of the endpoint map $`(x,y,z)+\lambda F(x,y,z)`$.
For $`\lambda =1/4(q8)`$ we get the points $`(0,0,z(8/q8))`$ , for $`\lambda =1/4q`$ we get the points $`(8x/q,8y/q,0)`$.
The conclusion is that the focal locus consists of the $`z`$-axis and of the circle $`z=0,x^2+y^2=4`$. That is, our surface is strongly focally degenerate, and we can indeed see geometrically the two families of circles corresponding to the two components of the focal locus.
We end this protracted example by observing that the rotation surface is clearly a rational surface. Indeed, we can say more, since a smooth model is obtained by blowing up the singular conic $`Q_{\mathrm{}}`$ and the two points $`P,P^{}`$.
Let $`R`$ and $`E,E^{}`$ be the respective exceptional divisors in the blow-up $`\stackrel{~}{๐}`$ of $`๐^3`$: the first is a ruled surface $`๐(๐ช_{๐^\mathrm{๐}}๐ช_{๐^\mathrm{๐}}(2))`$ , the other are two $`๐^2`$โs.
Let $`S`$ be the strict transform of $`X`$: it belongs to the linear system $`|4H2R2E2E^{}|`$, whereas the canonical system of $`\stackrel{~}{๐}`$ equals $`|4H+R+2E+2E^{}|`$. Thus $`S`$ belongs to $`|KR|`$ , and by the exact sequence
$`0๐ช_{\stackrel{~}{๐}}`$ $`(S)๐ช_{\stackrel{~}{๐}}`$ $`๐ช_S0`$
we infer $`h^i(๐ช_S)=h^{i+1}(๐ช_{\stackrel{~}{๐}}`$ $`(K+R)=h^{2i}(๐ช_{\stackrel{~}{๐}}`$ $`(R)=0`$, since $`R`$ is irreducible. $`S`$ is clearly then rational, and the anticanonical effective divisor has self-intersection $`4`$.
###### Example 8
More generally, for a rotation surface $`(r(s)cos(\theta ),r(s)sin(\theta ),z(s))`$
the unit normal is given by $`(z^{}(s)cos(\theta ),z^{}(s)sin(\theta ),r^{}(s))`$ , therefore we see easily that the focal locus consists of the $`z`$-axis and of the rotation surface obtained by rotating the evolute of the plane curve $`C=\{r(s),z(s)\}`$ we were starting with.
Therefore, general rotation surfaces provide examples of weakly but not strongly focally degenerate varieties.
###### Example 9
This last example shows the important role of the algebraic function r(s).
Let $`\mathrm{\Sigma }`$ be the line $`\{(0,0,s)๐^3\}`$ : then if we take the function $`r(s)R`$, where $`R๐`$ is a constant, the inverse construction yields a cylinder $`X^{}`$. Instead, if we choose $`r(s)R+s^2`$, we obtain as $`X^{}`$ simply a circle in the plane $`z=0`$.
## 7 Isotropically focally degenerate hypersurfaces
In the preceding section we gave a characterization, in terms of the inverse focal construction, of the focally isotropically degenerate varieties. However, in general such a construction yields a hypersurface, which is only weakly degenerate, and although in the next section we shall write down conditions which characterize the focally isotropically degenerate case, in the case of hypersurfaces, we can give an easier characterization for the isotropically degenerate case with a direct proof.
Let thus $`F(x_1,\mathrm{}x_{n+1})=0`$ be the polynomial equation of an affine hypersurface $`X`$, which we may assume, without loss of generality, to be irreducible.
Again the gradient $`F`$ of $`F`$ gives a map of the Normal Bundle $`N_X`$, $`\pi :N_X๐^{๐ง+\mathrm{๐}}`$ which we will also call the endpoint map
$`ฯต(x,\lambda )=x+\lambda F(x)`$, where $`x`$ is a point of $`X`$ ( thus, for $`\lambda =0`$ we reobtain the points of $`X`$).
###### Proposition 4
Let $`X`$ be a projective hypersurface : then $`X`$ is focally isotropically degenerate if and only if $`X`$ coincides with its focal locus $`\mathrm{\Sigma }_X`$.
Proof
In this case the focal locus equals the image $`\mathrm{\Sigma }_X`$ of the map $`\pi :N_X๐^{n+1}`$, and since $`X`$ may be assumed to be irreducible, $`N_X`$ is irreducible, whence $`\mathrm{\Sigma }_X`$ is also irreducible. But $`X`$ is contained in $`\mathrm{\Sigma }_X`$ and has not lesser dimension, thus equality holds. $`\mathrm{}`$
###### Remark 13
We derive thus the equality
$`F(x+\lambda F(x))0`$ $`\lambda `$.
In particular $`(d/d\lambda )F(x+\lambda F(x))0`$, and, for $`\lambda =0`$, we get
(I) $`<F(x),F(x)>0`$.
By the previous proposition the general fibre of $`\pi `$ has dimension $`1`$, and for each $`x_0X`$, $`\lambda _0๐`$ there exists a curve
(II) $`x(t),\lambda (t)`$ such that $`x(0)=x_0,\lambda (0)=\lambda _0`$, which is a fibre of $`\pi `$.
Since a fibre intersects a normal line $`x_0\times ๐`$ in at most one point, it follows that up to a birational transformation we can take $`(x_0,t)`$ as coordinates on $`N_X`$ by taking the curves $`x(x_0,t),\lambda (x_0,t)`$ satisfying (II) for $`\lambda _0=0`$, and assume that the curve $`x(x_0,t)`$ is a non constant curve in $`X`$ satisfying
(III) $`x(x_0,t)+\lambda (x_0,t)F(x(x_0,t))x_0`$.
We argue as in the preceding section :
$`x(x_0,t)x_0\lambda (x_0,t)F(x(x_0,t))`$,
thus by (I) our usual function $`r(x_0)0`$ and
(IV) $`<x(x_0,t)x_0,x(x_0,t)x_0>0`$.
In this case we also get, if $`s=(s_1,\mathrm{}.,s_n)`$ are local coordinates for $`x_0X`$, that $`(dr(s)/ds_j)0`$ and
(V) $`0=2<x(x_0(s),t)x_0(s),(dx_0(s)/ds_j)>`$ for each $`s,t`$.
Since the tangent space to $`X`$ at $`x_0`$ has dimension $`n`$, we infer that, fixing $`s`$ and varying $`t`$, we obtain a curve $`x(x_0,t)`$ which moves on the line through $`x_0`$ with direction $`F(x_0)`$.
We can thus write
(VI) $`x(x_0,t)=x_0+\mu (x_0,t)F(x_0)`$,
and then (VI) and (III) combine to yield
(VII) $`\lambda (x_0,t)F(x_0,t)\mu (x_0,t)F(x_0)`$.
Since the function $`\lambda (x_0,t)`$ is non zero, it follows that not only the line through $`x_0`$ with direction $`F(x_0)`$ is contained in $`X`$, but also that the normal direction stays constantly proportional to $`F(x_0)`$ on it.
We have thus proven the following
###### Theorem 4
A hypersurface $`X`$ is isotropically focally degenerate if and only if it is isotropically developable, i.e., for each point the normal line is contained in $`X`$, and along this line the tangent space to $`X`$ does not vary.
We would like now to give some examples and show where lies the difficulty in the fine classification of isotropically focally degenerate hypersurfaces.
It is classically known that in 3-space the analytical surfaces which are developable are only cones, cylinders, and tangential developable surfaces.
###### Proposition 5
Assume $`X`$ is an isotropically developable surface. Then , if $`X`$ a cylinder then $`X`$ is a plane. If $`X`$ is a cone , it is the cone over $`Q_{\mathrm{}}`$ with vertex in a point of affine space.
Proof
If $`X`$ is a cylinder, then the generatrices are the normal lines, therefore the normal direction is constant on the whole surface and the surface is a plane. If $`X`$ is a cone, with vertex, say, at the origin, then the vectors $`x`$ and $`F(x)`$ are proportional,
but the vector $`F(x)`$ is always isotropic, whence $`<x,x>0`$ on $`X`$, q.e.d. $`\mathrm{}`$
Let us now discuss the tangential surface $`X`$ of a curve $`C`$.
We write as usual $`X`$ parametrically as
$`x(s,t)=\alpha (s)+t\alpha ^{}(s)`$,
so that the tangent plane is generated by the two vectors
$`\alpha ^{}(s),\alpha ^{\prime \prime }(s)`$.
Up to local analytic reparametrization we can assume that one and only one of the following two possibilities occurs:
(I) $`<\alpha ^{}(s),\alpha ^{}(s)>0`$
(U) $`<\alpha ^{}(s),\alpha ^{}(s)>1`$.
In both cases follows that
(\*) $`<\alpha ^{}(s),\alpha ^{\prime \prime }(s)>0`$.
In the isotropic case (I), then clearly $`\alpha ^{}(s)`$ is a normal vector to $`X`$, constant on the generatrices, and our $`X`$ is thus isotropically developable.
We could stop our discussion here, since the isotropic ruling, in the situation we are interested in, is obtained by fixing $`s`$ and varying $`t`$, which means that we are in principle through with our discussion. Nevertheless, for curiosity, we analyze also the unitary case which we could avoid to consider in view of the assumption that our surface is not only developable, but also isotropically developable.
###### Lemma 6
The unitary case (U) occurs only if the curve $`C`$ is a plane curve, thus its tangential surface is a plane.
Proof
In the unitary case (U), the normal vector must be proportional to $`\alpha ^{\prime \prime }(s)`$, whence $`X`$ is isotropically degenerate if and only if
(\**) $`<\alpha ^{\prime \prime }(s),\alpha ^{\prime \prime }(s)>0`$. Now, by taking derivatives of (\*) and (\**), and using (\**), we obtain
(\***) $`<\alpha ^{\prime \prime }(s),\alpha ^{\prime \prime \prime }(s)>0`$
$`<\alpha ^{}(s),\alpha ^{\prime \prime \prime }(s)>0`$
ยฟfrom which it follows that $`\alpha ^{\prime \prime \prime }(s),\alpha ^{\prime \prime }(s)`$ are proportional vectors, whence also
$`<\alpha ^{\prime \prime \prime }(s),\alpha ^{\prime \prime \prime }(s)>0`$.
By induction, we show that for each integer n
(\*n\*) $`<\alpha ^{\prime \prime }(s),\alpha ^{(n)}(s)>0`$
$`<\alpha ^{}(s),\alpha ^{(n)}(s)>0`$
whence $`\alpha ^{\prime \prime }(s),\alpha ^{(n)}(s)`$ are proportional and thus also
$`<\alpha ^{(n)}(s),\alpha ^{(n)}(s)>0`$.
Consider now the Taylor development of $`\alpha (s)`$ at any point : from the fact that all higher derivative vectors are proportional follows that $`\alpha (s)`$ yields a plane curve.
But this means that its tangential surface is a plane.
$`\mathrm{}`$
It is now clear that in order to classify the non-trivial isotropically developable surfaces in 3-space we would need to classify the isotropic space curves $`C`$ ( i.e., those whose tangent vector is always an isotropic vector,that is, (I) holds).
Now, the condition that $`C`$ is algebraic is an obstacle!
Indeed, $`C`$ will be the birational image of a smooth curve $`B`$ , given through $`4`$ sections $`(s_0,s_1,s_2,s_3)`$ of a line bundle on $`B`$: the isotropicity condition amounts to the following equation ( where represents the derivative with respect to a local parameter on $`B`$)
($``$) $`\mathrm{\Sigma }_{i=1,2,3}(s_i^{}s_0s_0^{}s_i)^20`$.
###### Example 10
It is rather easy to give examples of rational curves which are isotropic.
It suffices, chosen an affine coordinate $`t`$ on $`๐`$, to set $`s_01`$ and let $`(s_1,s_2,s_3)`$ be polynomials in $`t`$ such that their derivatives satisfy
$`\mathrm{\Sigma }_{i=1,2,3}(s_i^{})^20`$.
In other words ,$`(s_1^{},s_2^{},s_3^{})`$ give a rational parametrization of the conic $`Q_{\mathrm{}}`$ and $`(s_1,s_2,s_3)`$ are taken to be the integrals of the three polynomials $`(s_1^{},s_2^{},s_3^{})`$.
In this way we see more generally that, up to translation, our curve $`C`$ is determined by our map $`BQ_{\mathrm{}}`$.
Using in our particular case of the rational curves fixed isomorphisms of $`๐^\mathrm{๐}`$ with $`B`$ and with $`Q_{\mathrm{}}`$, we obtain that our isotropic rational curves are parametrized by a pair of polynomials $`f_0(t),f_1(t)`$.
In concrete terms , we may take
$`(s_1^{},s_2^{},s_3^{})`$ = $`(f_0^2+f_1^2,if_1^2if_0^2,2if_0f_1)`$.
Assume now not only that the map $`f_0(t),f_1(t)`$ is of positive degree and is primitive (does not factor through an intermediate cover), e.g. it could be a cyclic Galois cover of prime order $`p`$.
If the map $`(s_1,s_2,s_3)`$ would not be birational onto its image, then the tangent map from $`C`$ to $`Q_{\mathrm{}}`$ would be a birational isomorphism.
But, in the example we gave above, $`f_0(t)=1,f_1(t)=t^p`$, we see immediately that $`(s_1,s_2,s_3)`$ are not polynomials in $`t^p`$.
## 8 Isotropically focally degenerate varieties and further examples
In the previous section we have given a classification, and concrete examples of isotropically focally degenerate hypersurfaces.
It is easy to obtain concrete examples in higher codimension by the following simple device : consider varieties $`M๐^๐ฆ`$, $`W๐^๐ฐ`$ and consider the product variety $`X=M\times W`$ in the orthogonal direct sum $`๐^๐ฆ^{}๐^๐ฐ=๐^{๐ฆ+๐ฐ}`$ .
It is immediate to see that in this case the normal bundle of $`X`$ is a product, likewise the endpoint map.
###### Remark 14
If thus $`M๐^๐ฆ`$ and $`W๐^๐ฐ`$ are isotropically focally degenerate, then $`X=M\times W๐^{๐ฆ+๐ฐ}`$ is also isotropically focally degenerate, and $`\mathrm{\Sigma }_X=\mathrm{\Sigma }_M\times \mathrm{\Sigma }_W`$. In particular, we obtain in this way $`\mathrm{\Sigma }_X`$ of arbitrary codimension.
We obtain also, by letting $`M`$ be an isotropically developable hypersurface, and $`W`$ general, an example of a variety $`X`$ of arbitrary codimension which is isotropically focally degenerate, and whose $`\mathrm{\Sigma }_X`$ is a hypersurface.
We now finally observe that the inverse focal construction gives a characterization of the isotropically focally degenerate varieties in terms of their focal variety $`\mathrm{\Sigma }_X`$.
###### Theorem 5
Let $`\mathrm{\Sigma }`$ be a projective variety of dimension $`a`$, and let $`r(s)`$ be an algebraic function on its affine part. Assume moreover that
1) at any point $`O(s)`$ of $`\mathrm{\Sigma }`$ the vector tangent space to $`\mathrm{\Sigma }`$ at $`O(s)`$ (a vector subspace of $`V^{}`$) has a totally isotropic annihilator.
Then, if $`X`$ is gotten from $`(\mathrm{\Sigma },r)`$ via the inverse focal construction, and moreover
2) the algebraic function $`r(s)`$ satisfies the conditions
2.1) $`dr(s)Im(T_{\mathrm{\Sigma },s}^Q_{\mathrm{}}T_{\mathrm{\Sigma },s}^{})`$
2.2) given $`\xi `$ with Im($`\xi `$) = $`df`$ , then $`1/4<\xi ,\xi >=r(s)`$
then $`X`$ is isotropically focally degenerate and $`\mathrm{\Sigma }=\mathrm{\Sigma }_X`$.
Proof
This follows immediately from Lemma 3, since conditions 1) and 2.1) imply that on the affine space given by equations 4โโ) the quadratic function $`Q_{\mathrm{}}`$ is constant, and 2.2) guarantees that this constant equals $`r(s)`$, whence also 3โโ) is satisfied and thus the sphere $`X_s`$, fibre over the point $`O(s)`$, is then an affine space of dimension $`ma`$.
$`\mathrm{}`$
###### Remark 15
The above theorem immediately implies the characterization given in the previous section of the isotropically focally degenerate hypersurfaces. Because in the case of hypersurfaces we noticed that $`X=\mathrm{\Sigma }_X`$, and then the tangential condition on $`\mathrm{\Sigma }_X`$ reads out as the condition that the normal vector is isotropic, moreover by the inverse focal construction $`X`$ is developable, and the fibre dimension equals $`ma=m(m1)=1`$. Thus $`X`$ is developable with the ruling by lines given by the normal direction.
We end by showing an explicit example of the situation considered in case 3) of Theorem 3.
###### Example 11
Consider first $`X^{}๐^\mathrm{๐}`$ obtained as the product $`X^{}=M\times W`$ of ยฟtwo isotropically developable surfaces :
thus $`X^{}`$ has a parametrization
$`(\alpha (s)+t\alpha ^{}(s),\beta (\sigma )+\tau \beta ^{}(\sigma ))`$.
Inside $`X^{}`$ we consider the divisor $`X`$ obtained by setting $`\tau =t`$.
Whence $`X`$ has a parametrization
$`(\alpha (s)+t\alpha ^{}(s),\beta (\sigma )+t\beta ^{}(\sigma ))`$,
and, remembering that $`\alpha ^{}(s),\beta ^{}(\sigma )`$ are isotropic vectors it follows that the normal space to $`X`$ is spanned, at the smooth points of $`X`$ , by the three vectors $`(\alpha ^{}(s),0)`$ , $`(0,\beta ^{}(\sigma ))`$ and
$`(<\beta ^{\prime \prime \prime }(\sigma ),\beta ^{}(\sigma )>[\alpha ^{\prime \prime }(s)+t\alpha ^{\prime \prime \prime }(s)],<\alpha ^{\prime \prime \prime }(s),\alpha ^{}(s)>[\beta ^{\prime \prime }(\sigma )+t\beta ^{\prime \prime \prime }(\sigma )])`$.
The endpoint map is given by
$`(\alpha (s)+t\alpha ^{}(s)+\lambda _1\alpha ^{}(s)\lambda _3<\beta ^{\prime \prime \prime }(\sigma ),\beta ^{}(\sigma )>[\alpha ^{\prime \prime }(s)+t\alpha ^{\prime \prime \prime }(s)],\beta (\sigma )+t\beta ^{}(\sigma )+\lambda _2\beta ^{}(\sigma )+\lambda _3<\alpha ^{\prime \prime \prime }(s),\alpha ^{}(s)>[\beta ^{\prime \prime }(\sigma )+t\beta ^{\prime \prime \prime }(\sigma )])`$
thus its image equals the image of the map
$`(\alpha (s)+\lambda _1\alpha ^{}(s)\lambda _3<\beta ^{\prime \prime \prime }(\sigma ),\beta ^{}(\sigma )>[\alpha ^{\prime \prime }(s)+t\alpha ^{\prime \prime \prime }(s)],\beta (\sigma )+\lambda _2\beta ^{}(\sigma )+\lambda _3<\alpha ^{\prime \prime \prime }(s),\alpha ^{}(s)>[\beta ^{\prime \prime }(\sigma )+t\beta ^{\prime \prime \prime }(\sigma )])`$.
To simplify the discussion we may assume $`<\alpha ^{\prime \prime }(s),\alpha ^{\prime \prime }(s)>1`$, and similarly $`<\beta ^{\prime \prime }(\sigma ),\beta ^{\prime \prime }(\sigma )>1`$, therefore our formula simplifies to
$`(\alpha (s)+\lambda _1\alpha ^{}(s)\lambda _3[\alpha ^{\prime \prime }(s)+t\alpha ^{\prime \prime \prime }(s)],\beta (\sigma )+\lambda _2\beta ^{}(\sigma )+\lambda _3[\beta ^{\prime \prime }(\sigma )+t\beta ^{\prime \prime \prime }(\sigma )])`$
and we see that the image of the normal bundle $`N_X`$ will in general be dominant.
Therefore $`X`$ is weakly focally degenerate, not focally isotropically degenerate, but the inverse focal construction reconstructs only the isotropically focally degenerate fourfold $`X^{}`$.
References
$`[Arnoldetal.]`$ V. Arnold, A. Varchenko, S. Goussein-Zadeโ, Singularites des Applications Differentiables. I : Classification des points critiques, des caustiques et des fronts dโ onde, Editions MIR, Moscow ( 1986), translation of the Russian 1982 edition.
$`[Coolidge]`$ J.L. Coolidge, A Treatise on Algebraic Plane Curves, Oxford University press (1931), reedited by Dover, New York ( 1959).
$`[DFN]`$ B.A. Dubrovin, A.T. Fomenko, S.P. Novikov Modern Geometry - Methods and Applications. Part 2: The geometry and topology of manifolds, reedited in Springer GTM 104, (1985).
$`[Fantechi]`$ B. Fantechi, The Evolute of a Plane Algebraic Curve, UTM 408, University of Trento, (1992).
$`[Hartshorne]`$ R. Hartshorne, Algebraic Geometry, Springer GTM 52 (1977).
$`[Milnor]`$ J. Milnor, Morse Theory, Princeton, Princeton Univ. Press, (1973).
$`[SalmonFiedler]`$ G. Salmon, W. Fiedler, Analytische Geometrie des Raumes, 2 vols., Leipzig, Teubner-Verlag, (1879-1880).
$`[Trifogli]`$ C. Trifogli, Focal Loci of Algebraic Hypersurfaces: a General Theory, Geom. Dedicata 70: 1-26 (1998).
$`[Trif2]`$ C. Trifogli, The Geometry of Focal Loci, Ph. D. Thesis, Universitaโ di Milano (in preparation).
Address of the authors:
Fabrizio Catanese,
Mathematisches Institut der Georg August Universitรคt Gรถttingen,
Bunsenstrasse 3-5, D 37073 Gรถttingen,
catanese@uni-math.gwdg.de
Cecilia Trifogli,
Dipartimento di Matematica F. Enriques, Universitaโ di Milano,
via C. Saldini 50,I 20133 Milano
trifogli@vmimat.mat.unimi.it and : All Souls College, Oxford OX1 4AL
cecilia.trifogli@all-souls.ox.ac.uk
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# Translational ChernโSimons Action and New Planar Particle Dynamics
## 1 Introduction
The aim of this paper is to study a new type of gravitational point particle interactions in two space dimensions. The classical and quantum dynamics of point sources interacting with standard (2+1)-dimensional gravity has been well known since the papers by Staruszkiewicz, Deser, Jackiw and tโHooft were published (see \[1-4\]).
In this paper we consider a new system of coupled two-dimensional torsion fields with point particles, derived in the following way:
i) We assume a nonstandard gravitational action described by a translational ChernโSimons term \[5-7\]
$$\text{S}_T=\frac{1}{\lambda }d^3x\epsilon ^{\mu \nu \rho }E_\mu ^{\underset{ยฏ}{\alpha }}T_{\nu \rho }^{\underset{ยฏ}{\beta }}\eta _{\underset{ยฏ}{\alpha }\underset{ยฏ}{\beta }},$$
(1.1)
where $`E_\mu ^{\underset{ยฏ}{\nu }}`$ ($`\overline{\nu }=0,1,2;\mu =0,1,2`$) describes (2+1)โdimensional dreibeins, and the Abelian strength
$$T_{\mu \nu }^{\underset{ยฏ}{\rho }}=_\mu E_\nu ^{\underset{ยฏ}{\rho }}_\nu E_\mu ^{\underset{ยฏ}{\rho }}$$
(1.2)
provides the components of a (2+1)-dimensional torsion tensor field. Note that a linear torsion Lagrangian taking the form of a translational Chern-Simons term can be introduced only in (2+1) dimensions.
Here we would like to recall that, following Witten , the standard gravitational Einstein action in (2+1) dimensions is described by another ChernโSimons theory
$$\text{S}_E=\frac{1}{\lambda }d^3x\epsilon ^{\mu \nu \rho }E_\mu ^{\underset{ยฏ}{\alpha }}\omega _{\nu \rho }^{\underset{ยฏ}{\beta }}\eta _{\underset{ยฏ}{\alpha }\underset{ยฏ}{\beta }},$$
(1.3)
where $`\omega _{\mu \nu }^{\underset{ยฏ}{\alpha }}`$ are the $`O(2,1)`$ spin connection components.
ii) In this paper we consider gravitational interactions of nonrelativistic $`D=2`$ particles described by the trajectory $`\stackrel{}{x}(t)=(x_1(t),x_2(t))`$, which are covariant under the following set of nonrelativistic local transformations:
$$x_l^{}=x_l^{}(\stackrel{}{x},t),t^{}=t+a,$$
(1.4)
which imply that
$$\dot{x}_i^{}=\frac{x_i^{}}{x_j}\dot{x}_j^{}+\frac{x_i^{}}{t},\frac{dt^{}}{dt}=1.$$
(1.5)
The set of nonrelativistic reparametrizationโinvariant velocities then takes the form $`(\underset{ยฏ}{a},\underset{ยฏ}{b}=1,2)`$:
$$\xi ^{\underset{ยฏ}{a}}=E_i^{\underset{ยฏ}{a}}\dot{x}_i+E_{0}^{\underset{ยฏ}{a}},$$
(1.6)
which can be obtained from the relativistic reparametrizationโinvariant velocities
$$\dot{x}^\mu \xi ^{\underset{ยฏ}{\mu }}=E_\nu ^{\underset{ยฏ}{\mu }}\dot{x}^\nu =(\xi ^{\underset{ยฏ}{0}},\xi ^{\underset{ยฏ}{a}}),$$
(1.7)
by putting $`x_0=t`$ and introducing a nonrelativistic gauge fixing condition $`E_{0}^{\underset{ยฏ}{0}}=1`$, $`E_i^{\underset{ยฏ}{0}}=0`$. In such a gauge we obtain the nonrelativistic form of the action (1.1), invariant under (1.4)
$$\text{S}_T^{(\mathrm{NR})}=\frac{1}{\lambda }d^3xe^{\mu \nu \rho }E_\mu ^{\underset{ยฏ}{a}}T_{\nu \rho }^{\underset{ยฏ}{a}}=\frac{1}{\lambda }๐td^2x\left(e^{\underset{ยฏ}{a}}B^{\underset{ยฏ}{a}}\epsilon _{jk}h_j^{\underset{ยฏ}{q}}_k^{\underset{ยฏ}{a}}\right),$$
(1.8)
where
$$B^{\underset{ยฏ}{a}}=\epsilon _{jk}_jh_k^{\underset{ยฏ}{a}}_k^{\underset{ยฏ}{a}}=_th_k^{\underset{ยฏ}{a}}_ke^{\underset{ยฏ}{a}}.$$
(1.9)
We denote $`h_\kappa ^{\underset{ยฏ}{a}}=E_\kappa ^{\underset{ยฏ}{a}}`$ and $`e^{\underset{ยฏ}{a}}=E_0^{\underset{ยฏ}{a}}`$. The fields $`B^{\underset{ยฏ}{a}}`$ and $`_k^{\underset{ยฏ}{a}}`$ play the role of the magnetic and electric fields, respectively, with the internal $`O(2)`$ index $`\underset{ยฏ}{a}`$ describing the $`D=2`$ nonrelativistic rotation group.
In Sect. 2 we discuss a coupled classical field + point particles system. Firstly, we present general considerations for the $`N`$โbody particle dynamics, and then we consider the twoโbody problem in classical theory. In this case we obtain, in terms of relative planar coordinates, a new type of dynamics which can be understood in two equivalent ways:
$``$ as a $`D=2`$ free particle system with nonstandard, symplectic structure
$``$ as a modified $`D=2`$ free particle system with an energy-dependent angular momentum $`\overline{l}(E)`$ and a canonical symplectic structure.
We also show that the classical solutions split into two classes for $`\lambda \overline{l}<0`$
$``$ for the planar particle distances
$$r<r_0=\left(\frac{|\lambda \overline{l}(E)|}{2\pi }\right)^{\frac{1}{2}}$$
they describe bounded motion while
$``$ for $`r>r_0`$ they describe the scattering.
In Sect. 3 we introduce the corresponding quantum mechanics, described by a nonstandard Schrรถdinger equation, and discuss the energy eigenvalues and wave eigenfunctions.
We find that the exact values of the energy levels for the confined motion ($`r<r_0`$) can be determined only numerically.
The final remarks, concerning possible physical interpretations of our solutions are presented in the last, fourth Section.
It should be added that the present model is a generalization of the $`(1+1)`$โdimensional model, presented by one of the present authors (PCS) in .
## 2 New Planar Particle Dynamics: Classical Theory
### 2.1 $`๐`$โParticle System
Let us discuss, in some detail, a set of coupled equations describing an interacting two-dimensional particle โ torsion field system, with the torsion fields being given by the Abelian ChernโSimons action (1.3). Introducing $`N`$ trajectories $`\stackrel{}{x}_\alpha (t)=\left\{x_\alpha ^i(t)\right\}`$ ($`i=1,2;\alpha =1,\mathrm{}N`$) for $`N`$ particles and the notation
$$E_{\mu ;\alpha }^{\underset{ยฏ}{a}}(t)E_\mu ^{\underset{ยฏ}{a}}(\stackrel{}{x}_\alpha (t),t),$$
(2.1)
we find that the action for $`N`$ particles in $`D=2`$ dimensions, in the first order formalism, can be written as
$$\text{S}_{\mathrm{part}}^{(N)}=\underset{\alpha =1}{\overset{N}{}}m_\alpha ๐t\left[\frac{1}{2}\xi _\alpha ^{\underset{ยฏ}{a}}\xi _\alpha ^{\underset{ยฏ}{a}}\xi _\alpha ^{\underset{ยฏ}{a}}\left(E_{j;\alpha }^{\underset{ยฏ}{a}}\dot{x}_\alpha ^j+E_{0,\alpha }^{\underset{ยฏ}{a}}\right)\right],$$
(2.2)
thus providing the constraint formula
$$\xi _\alpha ^{\underset{ยฏ}{a}}=E_{j;\alpha }^{\underset{ยฏ}{a}}\dot{x}^j+E_{0,\alpha }^{\underset{ยฏ}{a}}.$$
(2.3)
We can now derive the equations describing the $`D=2`$ coupled particleโtorsion field system, described by the action $`S=\text{S}_{\mathrm{part}}^{(N)}+\text{S}_\mathrm{T}^{(\mathrm{NR})}`$ (see (1.8) and (2.2)). They take the form:
$$\dot{\xi }_\alpha ^{\underset{ยฏ}{a}}E_{i,\alpha }^{\underset{ยฏ}{a}}+\xi _\alpha ^{\underset{ยฏ}{a}}T_{\mu i,\alpha }^{\underset{ยฏ}{a}}\dot{x}^{\mu ,\alpha }=\mathrm{\hspace{0.17em}0},$$
(2.4)
$$T_{\mu \nu }^{\underset{ยฏ}{a}}=_\mu E_\nu ^{\underset{ยฏ}{a}}_\nu E_\mu ^{\underset{ยฏ}{a}}=\frac{\lambda }{2}ฯต_{\mu \nu \rho }\underset{\beta }{}\xi _\beta ^{\underset{ยฏ}{a}}\dot{x}^{\rho ,\beta }\delta (\stackrel{}{x}\stackrel{}{x}^\beta ).$$
(2.5)
The second term in (2.4), due to the antisymmetry of the $`ฯต`$โtensor in (2.5), can be rewritten as
$$\xi _\alpha ^{\underset{ยฏ}{a}}T_{\mu i,\alpha }^{\underset{ยฏ}{a}}\dot{x}^{\mu ,\alpha }=\frac{\lambda }{2}ฯต_{\mu i\rho }\underset{\genfrac{}{}{0pt}{}{\beta }{\beta \alpha }}{}\xi _\beta ^{\underset{ยฏ}{a}}\xi _\alpha ^{\underset{ยฏ}{a}}\dot{x}^{\rho ,\beta }\dot{x}^{\mu ,\alpha }\delta (\stackrel{}{x}_\alpha \stackrel{}{x}_\beta ),$$
(2.6)
which is infinite for coinciding particle positions and vanishes otherwise. If we restrict our configuration space only to noncoinciding particle positions we have
$$\dot{\xi }_\alpha ^{\underset{ยฏ}{a}}E_{\alpha ;i}^{\underset{ยฏ}{a}}=\mathrm{\hspace{0.17em}0},$$
(2.7)
which leads, for points in the configuration space where the metric is non-degenerate, to
$$\dot{\xi }_\alpha ^{\underset{ยฏ}{a}}=\mathrm{\hspace{0.17em}0}.$$
(2.8)
Let us now consider the field equation (2.5). Its general solution can be written in the pure gauge form
$$E_\mu ^{\underset{ยฏ}{a}}(\stackrel{}{x},t)=\stackrel{~}{E}_\mu ^{\underset{ยฏ}{a}}(\stackrel{}{x},t)+_\mu \mathrm{\Lambda }^{\underset{ยฏ}{a}}(\stackrel{}{x},t)=_\mu \stackrel{~}{\mathrm{\Lambda }}^{\underset{ยฏ}{a}}(\stackrel{}{x},t),$$
(2.9)
where
$$\stackrel{~}{E}_\mu ^{\underset{ยฏ}{a}}(\stackrel{}{x},t):=\frac{\lambda }{4\pi }_\mu \underset{\alpha }{}\xi _\alpha ^{\underset{ยฏ}{a}}\mathrm{\Phi }(\stackrel{}{x}\stackrel{}{x}_\alpha )$$
(2.10)
and
\- $`\mathrm{\Lambda }^{\underset{ยฏ}{a}}`$ is an $`O(2)`$โvector valued pair of regular gauge functions,
\- $`\mathrm{\Phi }(\stackrel{}{x})`$ is a singular gauge function satisfying the following equation (see e.g. ):
$$ฯต^{ij}_i_j\mathrm{\Phi }(\stackrel{}{x})=2\pi \delta (\stackrel{}{x}).$$
(2.11)
As a solution of (2.11) we can take
$$\mathrm{\Phi }(\stackrel{}{x})=\text{arc}\text{tan}\frac{x_2}{x_1}$$
(2.12)
i.e.
$$_k\mathrm{\Phi }(x)=\epsilon _{kl}_l\mathrm{ln}|\stackrel{}{x}|.$$
(2.13)
where (2.13) has to be regularized in such a way that it is well defined everywhere and therefore vanishes for $`\stackrel{}{x}0`$ (see e.g. ).
Let us note that the solutions for the fields $`E_\mu ^{\underset{ยฏ}{a}}(\stackrel{}{x},t)`$ with asymptotically nonvanishing gauge function $`\mathrm{\Lambda }^a`$ are not solutions of the Hamiltonโs variational principle for the action $`S`$. The bad asymptotic (as $`r\mathrm{}`$) behaviour of $`E_\mu ^{\underset{ยฏ}{a}}(\stackrel{}{x},t)`$ leads to the appearance of nonvanishing surface integrals and, in consequence, our $`E_\mu ^{\underset{ยฏ}{a}}(\stackrel{}{x},t)`$ do not minimize the action. In order to take into consideration the particular asymptotic behaviour of $`\mathrm{\Lambda }^a`$ we fix our gauge as follows:
$$\mathrm{\Lambda }^i(\stackrel{}{x},t)=x^ia^i(t)$$
(2.14)
and modify the field Lagrangian by adding two surface integrals<sup>1</sup><sup>1</sup>1see for such a procedure.
$`I_1`$ $`=`$ $`{\displaystyle \frac{1}{\lambda }}{\displaystyle d^2xv^{\underset{ยฏ}{\text{a}}}ฯต^{ij}_i\stackrel{~}{E}_j^{\underset{ยฏ}{a}}(\stackrel{}{x},t)}`$ (2.15)
$`I_2`$ $`=`$ $`{\displaystyle \frac{1}{\lambda }}{\displaystyle d^2xฯต^{ij}_i\stackrel{~}{E}_0^j}`$ (2.17)
where we have defined
$$v^{\underset{ยฏ}{a}}(t):=\dot{a}^{\underset{ยฏ}{a}}(t).$$
(2.18)
Notice that with the gauge (2.14) the space part of our metric becomes asymptotically Euclidean. Due to the asymptotic behaviour
$$\stackrel{~}{E}_\mu ^{\underset{ยฏ}{a}}(\stackrel{}{x},t)O(r^1)\text{as}r\mathrm{}$$
(2.19)
the variations of the $`\stackrel{~}{E}_\mu ^{\underset{ยฏ}{a}}`$ and of the $`v^{\underset{ยฏ}{a}}`$ are independent of each other i.e. the fields $`\stackrel{~}{E}_\mu ^{\underset{ยฏ}{a}}`$ and $`v^{\underset{ยฏ}{a}}`$ can be introduced as new variables. The particle action (2.2) takes the form
$$\text{S}_{\mathrm{part}}^{(\mathrm{N})}[x,\xi ,E]=\text{S}_{\mathrm{part}}^{(\mathrm{N})}[x,\xi ,\stackrel{~}{E}]+๐t\left[\underset{\alpha }{}\xi _\alpha ^{\underset{ยฏ}{a}}\dot{x}_\alpha ^{\underset{ยฏ}{a}}v^{\underset{ยฏ}{a}}\underset{\alpha }{}\xi _\alpha ^{\underset{ยฏ}{a}}\right].$$
(2.19a)
where as the modified field action $`S_{\mathrm{field}}(\stackrel{~}{E})`$ is given by $`S_T^{(NR)}`$ (eq. (1.8)) but taken as a function of the $`\stackrel{~}{E}`$ now
$$\text{S}_{\mathrm{field}}(\stackrel{~}{E})=\text{S}_T^{(NR)}(\stackrel{~}{E})$$
(2.19b)
We see that in the action (2.19) we have separated the variables describing the โbulkโ $`(\stackrel{~}{E}_\mu ^a)`$ and asymptotic behaviour $`(v^a)`$. In technical terms, no surface integrals remain when we calculate the functional derivative of the action (2.19) with respect to the $`\stackrel{~}{E}`$. By varying $`S`$ with respect to $`v^{\underset{ยฏ}{a}}`$ we obtain the constraint
$$\underset{\alpha }{}\xi _\alpha ^{\underset{ยฏ}{a}}=0.$$
(2.20)
The choice (2.14) for $`\mathrm{\Lambda }^i`$ breaks asymptotically the invariance with respect to local space translations (1.4). However, as under general coordinate transformations the functions $`\stackrel{~}{\mathrm{\Lambda }}^i`$ are scalars, the changes $`\delta \stackrel{~}{\mathrm{\Lambda }}^i`$ under (1.4) of both the singular and regular parts of $`\stackrel{~}{\mathrm{\Lambda }}^i`$ must vanish and we obtain
$$\delta x^i=\delta a^i(t),$$
(2.21)
where $`x^i`$ as well as $`a^i`$ transform as vectors under rotations in tangent space. Therefore as a residual symmetry we get translations, local in time, and rigid rotations.
Let us calculate the corresponding generator of constant translations (2.21). We find by using Noetherโs theorem
$`P^i`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}p_\alpha ^i+{\displaystyle \frac{2}{\lambda }}{\displaystyle d^2xB^{\underset{ยฏ}{a}}\stackrel{~}{E}_i^{\underset{ยฏ}{a}}}`$ (2.22)
$`=`$ $`P_{\mathrm{part}}^i+P_{\mathrm{field}}^i,`$ (2.24)
where the canonical particle momenta are obtained from (2.2), (2.9) and (2.14) as
$$p_\alpha ^j=\xi _\alpha ^j\frac{\lambda }{4\pi }\underset{\genfrac{}{}{0pt}{}{\beta }{\beta \alpha }}{}\left(\xi _\alpha ^{\underset{ยฏ}{a}}\xi _\beta ^{\underset{ยฏ}{a}}\right)_j\mathrm{\Phi }(\stackrel{}{x}_{\alpha \beta }),$$
(2.25)
where $`\stackrel{}{x}_{\alpha \beta }=\stackrel{}{x}_\alpha \stackrel{}{x}_\beta `$.
Inserting the explicit expressions for $`B^{\underset{ยฏ}{a}}`$ and $`\stackrel{~}{E}_i^{\underset{ยฏ}{a}}`$ into (2.22) we obtain
$$P_{\mathrm{field}}^i=0.$$
(2.26)
Furthermore, from (2.23) we get
$$P_{\mathrm{part}}^i=\underset{\alpha }{}\xi _\alpha ^i.$$
(2.27)
The formula (2.25) describes the linear momentum of the CM-motion for the N-particle system which vanishes due to the constraint (2.20)<sup>2</sup><sup>2</sup>2Analogous situation has been observed in (2+ 1)-dimensional gravity .. The same result has been obtained by one of the present authors (PCS) in the one-dimensional case \[ 8\], \[ 9 \].
Next we pass to the Hamiltonian formulation. Applying the Legendre transformation to the Lagrangian (2.19) and using the constraint
$$B^{\underset{ยฏ}{a}}=ฯต_{ij}_jh_i^{\underset{ยฏ}{a}}=\frac{\lambda }{2}\underset{\alpha }{}\xi _\alpha ^{\underset{ยฏ}{a}}\delta ^{(2)}(xx_\alpha )$$
(2.28)
given by the components $`\mu =j`$, $`\nu =i`$ of (2.5) we find the N-particle Hamiltonian to be given by
$$H^{(\mathrm{N})}=\frac{1}{2}\underset{\alpha }{}\xi _\alpha ^{\underset{ยฏ}{a}}\xi _\alpha ^{\underset{ยฏ}{a}}+v^{\underset{ยฏ}{a}}\underset{\alpha }{}\xi _\alpha ^{\underset{ยฏ}{a}}.$$
(2.29)
### 2.2 Hamiltonian Formulation for the TwoโBody Problem
Let us consider now more explicitly the $`N=2`$ case.
We define:
$$\xi ^{\underset{ยฏ}{a}}:=\frac{1}{2}\left(\xi _1^{\underset{ยฏ}{a}}\xi _2^{\underset{ยฏ}{a}}\right),x^{\underset{ยฏ}{a}}:=x_1^{\underset{ยฏ}{a}}x_2^{\underset{ยฏ}{a}},p^{\underset{ยฏ}{a}}:=\frac{1}{2}(p_1^{\underset{ยฏ}{a}}p_2^{\underset{ยฏ}{a}}).$$
(2.30)
Then using the constraint (2.20) we get from (2.27)
$$H^{(2)}=\xi ^{\underset{ยฏ}{a}}\xi ^{\underset{ยฏ}{a}}$$
(2.31)
and
$$p_i=\xi _i+\frac{\lambda }{4\pi }(\xi ^{\underset{ยฏ}{a}}\xi ^{\underset{ยฏ}{a}})_i\mathrm{\Phi }(\stackrel{}{x}).$$
(2.32)
The Hamiltonian equations take the form
$$\dot{x}^i=\frac{H}{p^i}=\mathrm{\hspace{0.17em}2}\left(\xi ^{\underset{ยฏ}{a}}\frac{\xi ^{\underset{ยฏ}{a}}}{p^i}\right),$$
(2.31a)
$$\dot{p}_i=\frac{H}{x_i}=\mathrm{\hspace{0.17em}2}\left(\xi ^a\frac{\xi ^a}{x^i}\right).$$
(2.31b)
Using (2.32) we have
$$\xi ^{\underset{ยฏ}{a}}\frac{\xi ^{\underset{ยฏ}{a}}}{p^i}=\frac{\xi _i}{1+{\displaystyle \frac{\lambda }{2\pi }}(\xi ^{\underset{ยฏ}{a}}_a\mathrm{\Phi })},$$
(2.32a)
$$\xi ^a\frac{\xi ^a}{x^j}=\frac{\lambda }{4\pi }\frac{\left(\xi _a\xi ^a\right)\xi ^i_i_j\mathrm{\Phi }}{1+{\displaystyle \frac{\lambda }{2\pi }}(\xi ^a_a\mathrm{\Phi })}.$$
(2.32b)
Taking the time derivative of (2.32) and using (2.31โ32) we get
$$\dot{\xi }_i^a+\frac{\lambda }{2\pi }_i\mathrm{\Phi }(x)\xi _j\dot{\xi }_j=0.$$
(2.33)
However, instead of the canonical variables ($`x_i,p_i`$) we can use the variables ($`x_i,\xi _i`$). Then the Lagrangian obtained from the Hamiltonian (2.31) would have had the form:
$`L`$ $`=`$ $`p_l(\xi _i,x_l)\dot{x}_lH`$ (2.34)
$`=`$ $`\left(\xi _l+{\displaystyle \frac{\lambda }{4\pi }}\xi ^2_l\mathrm{\Phi }(x)\right)\dot{x}_l\xi ^2.`$ (2.35)
The variation with respect to $`\xi _i`$ is given by the expression
$$\xi _i=\frac{\frac{1}{2}\dot{x}_i}{1{\displaystyle \frac{\lambda }{4\pi }}\left(\dot{x}_j_j\mathrm{\Phi }\right)},$$
(2.36)
which is equivalent to the Hamiltonian eq. (2.31a) with the insertion of (2.32a). Inserting (2.35) in (2.34) we get
$$L=\frac{1}{4}\frac{\dot{x}_l^2}{1{\displaystyle \frac{\lambda }{4\pi }}\left(\dot{x}_l_l\mathrm{\Phi }\right)}.$$
(2.37)
In particular if we consider
$`det\left({\displaystyle \frac{^2L}{\dot{x}_i\dot{x}_j}}\right)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(1{\displaystyle \frac{\lambda }{4\pi }}\left(\dot{x}_l_l\mathrm{\Phi }\right)\right)^4`$ (2.38)
$`=`$ $`{\displaystyle \frac{1}{4}}\left(1+{\displaystyle \frac{\lambda }{2\pi }}\xi _l_l\mathrm{\Phi }\right)^{+4},`$ (2.39)
we see that when the velocities are expressible in terms of the canonical variables then (2.33) gives us
$$\dot{\xi }_l=0.$$
(2.40)
Using the Hamiltonian (2.31) we get for the pair of noncanonical variables the Hamilton equations (2.31a), as $`\dot{x}_l=\{x_l,H\}`$ and (2.38) as $`\dot{\xi }_l=\{\xi _l,H\}`$, if we assume the following nonstandard symplectic structure:
$$\{x^i,x^j\}=\{\xi _i,\xi _j\}=0,$$
(2.39a)
$$\{x^i,\xi _j\}=\delta _j^i\frac{\frac{\lambda }{2\pi }\xi ^i_j\mathrm{\Phi }}{1+{\displaystyle \frac{\lambda }{2\pi }}\left(\xi ^a_a\mathrm{\Phi }\right)}.$$
(2.39b)
It is easy to check that the Poisson brackets (2.39aโb) satisfy the Jacobi identity. For our system we have two conserved angular momenta. If we define (in $`D=2`$ $`\stackrel{}{a}\stackrel{}{b}=ฯต_{ij}a^ib^j`$)
$$l:=\stackrel{}{x}\stackrel{}{p}=\stackrel{}{x}\stackrel{}{\xi }+\frac{\lambda }{4\pi }H$$
(2.39c)
we find that $`l`$ is conserved as well as
$$\overline{l}:=\stackrel{}{x}\stackrel{}{\xi }$$
(2.39d)
because
$$\frac{d}{dt}\overline{l}=\dot{\stackrel{}{x}}\stackrel{}{\xi }=\mathrm{\hspace{0.17em}0}.$$
(2.39e)
However, as
$$\xi ^{\underset{ยฏ}{a}}_a\mathrm{\Phi }=\stackrel{}{\xi }\stackrel{}{}\mathrm{ln}|x^{\underset{ยฏ}{c}}|=\frac{\overline{l}}{r^2},$$
(2.39f)
where $`r:=|x^{\underset{ยฏ}{a}}|`$, we see that (2.35) can be rewritten as:
$$\dot{x}^{\underset{ยฏ}{a}}=\frac{2\xi ^{\underset{ยฏ}{a}}}{1+\frac{\lambda \overline{l}}{2\pi r^2}}.$$
(2.39g)
Note that if $`x^{\underset{ยฏ}{a}}(0)`$ is parallel to $`\xi ^{\underset{ยฏ}{a}}`$ we have $`\overline{l}=0`$ and, in consequence, a free motion on a line.
For $`\lambda \overline{l}<0`$ it is convenient to introduce the quantity
$$r_0^2:=\frac{|\lambda \overline{l}|}{2\pi }$$
(2.40)
as then, from the previous results, we see that the relative twoโbody motion separates in this case into a motion within the interior region given by
$$r<r_0$$
(2.41a)
and a motion within the exterior region given by
$$r>r_0.$$
(2.41b)
## 3 New Planar Dynamics: QuantumโMechanical TwoโBody Problem
The relation (2.32), after the substitution of (2.31), takes the form ($`HH^{(2)}`$):
$$\stackrel{}{\xi }=\stackrel{}{p}\frac{\lambda }{4\pi }H\stackrel{}{}\mathrm{\Phi }.$$
(3.1)
Squaring it and using again (2.29) we obtain
$$H=\stackrel{}{p}^2\frac{l^2}{r^2}+\frac{\overline{l}^2}{r^2},$$
(3.2a)
where we have used the definition (2.39d) of $`\overline{l}`$.
It should be stressed that
$$\overline{l}=l\frac{\lambda }{4\pi }H$$
(3.2b)
i.e. (3.2a) gives us a quadratic equation for $`H`$.
We quantize the problem by considering a Schrรถdinger-like equation
$$i\mathrm{}\frac{\psi (\stackrel{}{x},t)}{t}=\widehat{H}\psi (\stackrel{}{x},t)=\left[\widehat{\stackrel{}{p}}^2\frac{l^2}{r^2}+\frac{1}{r^2}\overline{l}^2\right]\mathrm{\Psi }(\stackrel{}{x},t)$$
(3.3)
in which the operators $`\widehat{H}`$ and $`\widehat{\stackrel{}{p}}`$ are defined by the usual quantization rules
$$\widehat{H}=i\mathrm{}\frac{}{t},\widehat{p}_l=\frac{\mathrm{}}{i}_l.$$
(3.4)
We see that the equation (3.3) describes a nonstandard form of a time dependent Schrรถdinger equation, with its right hand side containing both the first and second time derivatives (entering through $`\overline{l}`$).
For the stationary case, i.e. when $`\mathrm{\Psi }(\stackrel{}{x},t)=\mathrm{\Psi }_E(\stackrel{}{x})e^{\frac{iEt}{\mathrm{}}}`$ we can use the angular-momentum basis and put
$$\mathrm{\Psi }_{E,m}=f_{E,m}(r)e^{im\phi }$$
(3.5)
where $`m`$ is an integer, and find that $`f_{E,m}`$ satisfies a nonstandard time independent Schrรถdinger equation
$$\left[\mathrm{}^2\left(_r^2+\frac{1}{r}_r\frac{\overline{m}^2}{r^2}\right)E\right]f_{E,m}(r)=0,$$
(3.6)
where we have defined
$$\mathrm{}\overline{m}:=\mathrm{}m\frac{\lambda }{4\pi }E$$
(3.7)
i.e. $`\mathrm{}\overline{m}`$ is an eigenvalue of $`\overline{l}`$.
A characteristic feature of the Schrรถdinger equation (3.6) is the appearance of the noncanonical angular momentum $`\mathrm{}\overline{m}`$ with $`\overline{m}`$ not being an integer. i.e. our two-body system carries a fractional orbital angular momentum (see the discussion in or ). It can be shown that $`J=\overline{l}`$ is equal to the total angular momentum of the particle + field system with a nonvanishing angular momentum of the fields .
In the following we discuss only the most interesting case of $`\lambda \overline{l}<0`$.
Now the appropriate boundary conditions correspond to the requirement that $`f_{E,m}(r)`$ is nonzero in either the interior region ($`r<r_0`$) or in the exterior region ($`r>r_0`$). Thus our boundary condition is
$$f_{E,m}(r_0)=0$$
(3.8)
The general solution of (3.6) is given by
$$f_{E,m}(r)=Z_{\overline{m}}\left(\frac{\sqrt{E}}{\mathrm{}}r\right),$$
(3.9)
where $`Z_{\overline{m}}`$ is an appropriate Bessel function of order $`\overline{m}`$ (or a superposition of such functions).
### Interior solutions ($`r<r_0`$)
The only Bessel functions belonging to a self-adjoint differential operator in (3.6) are those of the first kind with $`\overline{m}0`$. Therefore a decomposition of space into interior/exterior regions appear only in the case of $`\lambda <0`$.
Then the possible eigenvalues $`E_n(m)`$ are determined by
$$J_{\overline{m}}\left[\frac{\sqrt{E}}{\mathrm{}}\left(\frac{\mathrm{}|\lambda |\overline{m}}{2\pi }\right)^{\frac{1}{2}}\right]=\mathrm{\hspace{0.17em}0}$$
(3.10)
with $`\overline{m}`$ given by (3.7).
To simplify (3.10) we define
$$ฯต=\frac{|\lambda |E}{2\pi \mathrm{}}.$$
(3.11)
Then (3.10) takes the form
$$J_{\overline{m}}(\overline{m}^{\frac{1}{2}}ฯต^{\frac{1}{2}})=\mathrm{\hspace{0.17em}0}.$$
(3.12)
As $`J_{\overline{m}}`$, for fixed $`\overline{m}>0`$, has no imaginary zeros (see e.g. ) we have to consider an infinite number of simple positive zeros, which we denote by $`y_n(\overline{m})`$, $`n=1,2..`$ arranged in ascending order of magnitude. Then we see that due to (3.7), the eigenvalues $`ฯต_n(m)`$ we are looking for, are the positive fixed points of the equation
$$ฯต=f_n(m+\frac{1}{2}ฯต),$$
(3.13)
subject to the conditions $`m+\frac{1}{2}ฯต_n(m)>0`$ and $`y_n(\overline{m})>0`$, where we have defined
$$f_n(\overline{m}):=\frac{1}{\overline{m}}y_n^2(\overline{m}).$$
(3.14)
The solutions of (3.13) can be obtained only numerically.
We insert (3.14) into (3.13) and after the use of approximate formulae for the zeros of Bessel functions of first kind we conclude that the low-lying energy levels are given by small values of $`n`$. In this case all values of $`mZZ`$ contribute to the energy spectrum ($`E>0`$) but the positivity of $`\overline{m}=m+\frac{1}{2}ฯต_n(m)>0`$ for higher $`n`$ implies that some finite range of values of $`m`$ has to be excluded. Qualitatively one can say that for every $`n`$ there appears however an infinite tower of energy values $`ฯต_n(m)`$.
Only in the case of large $`m`$ the analytical asymptotic results are available. Let us consider the choice $`n=1`$. Using the asymptotic behaviour ()
$$y_1(\overline{m})=\overline{m}+1.855757\overline{m}^{\frac{1}{3}}+O(\overline{m}^{\frac{1}{3}}),$$
(3.15)
valid for large $`\overline{m}`$, we obtain
$$ฯต_1(m)=2m+9.35243m^{\frac{1}{3}}+O(m^{\frac{1}{3}}).$$
(3.16)
valid for large positive $`m`$.
### Exterior solutions $`(r>r_0)`$
First of all, let us note that there are no bound states solutions of (3.6) for $`r>r_0`$ as the only square integrable Bessel functions in $`[r_0,\mathrm{})`$ are the modified Bessel functions of the third kind, which have neither positive nor pure imaginary zeros (see e.g. ).
Scattering solutions are given by a superposition of Bessel functions of the first and second kind
$$f_{E,m}(r)=A_m(E)J_{\overline{m}}\left(\frac{\sqrt{E}}{\mathrm{}}r\right)+B_m(E)Y_{\overline{m}}\left(\frac{\sqrt{E}}{\mathrm{}}r\right)$$
(3.17)
with the ratio $`\frac{A_m}{B_m}`$ determined by the boundary condition (3.8). These solutions describe scattering on an obstruction of radius $`r_0`$, which is dynamically determined.
## 4 Final Remarks
Gauging of translations for nonrelativistic point particles in a $`D=2`$ dimensional space coupled to a translational ChernโSimons action leads to a new and interesting $`D=(2+1)`$ particle dynamics. In particular, we have shown that for negative values of the product $`\lambda \overline{l}`$ of the coupling strength $`\lambda `$ and the fractional orbital angular momentum $`\overline{l}`$ the twoโbody motions separate into interior and exterior motions. In the quantum mechanical case we have obtained a new type of nonlinear Schrรถdinger equation, with a second order time derivative, which leads, in the stationary case, to the nonlinear energy eigenvalue problem. The interior solutions are described by a infinite tower of bound states for each fixed $`n`$. They describe confinement within a dynamically determined compact space region (see (2.41a)) (geometric bag). We should mention that analogous results were obtained by one of the present authors (PCS) in the oneโdimensional case . It is clear that a possible physical relevance of these results, which recall the features of confinement in one and two space dimensions, should be investigated further.
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# Entangled Simultaneous Measurement and Elementary Particle Representations
## 1 Introduction
The measurement of non-commuting operators in a quantum system is at the root of the well-known paradoxes of quantum mechanics. For example, the Kochen-Specker paradox involves measurement of a set of non-commuting operators with mutually commuting subsets . The Bell inequalities are similarly formulated by incompatible space-like separated spin measurements . The connection of these problems to quantum non-locality motivates the consideration of an alternative definition for the measurement of incompatible quantum observables. Because measured quantities are inherently statistical, quantum mechanics requires an ensemble of identical systems to establish an expectation value. The types of measurements on the ensemble in the situation of non-commuting observables is critical to the interpretation of the result. For example, the usual simultaneous measurement of position and momentum would require separate measurements of each operator on half the ensemble at the same time. This procedure for the assignment of position and momentum to a system, coupled with wavefunction collapse, leads to a built-in non-locality in system observables. One approach to this problem is to restrict valid measurement of non-commuting observables to be with special-purpose ancillary quantum systems (referred to as vacuum meters in this paper) that are entangled with the original system. In the alternative entangled simultaneous measurement scheme, all systems in the ensemble have the same experimental set-up with coupling to vacuum meters for a joint position/momentum measurement. The premise of this paper is that this procedure, while avoiding quantum non-locality, yields results that are properly interpreted as a joint measurement of non-commuting observables. Furthermore, it is assumed that, probably due to a principle involving quantum non-locality, entangled states are fundamental to particle representations. The principle mandates that particle states allow the simultaneous determination of all particle observables. These assumptions impose the structure of a Naimark extension on particle multiplets.
Allowed quantum numbers and statistics for relativistic particles are strictly constrained by the dual principles of Lorentz invariance and quantum mechanics . The latter requires states that represent symmetries of the lagrangian. A larger structure, the Naimark extension of the Poincare algebra , results if states are required for the measurement of incompatible observables with entangled vacuum meters. New commuting operators are defined in a Naimark extension that project to the original set upon meter measurement. A key component of the measurement scheme is the entanglement of the system with a vacuum state containing independent meters. The assumption of Fermi-Dirac statistics for the meters forces a different particle identity (flavor) onto the ancillary particles that make up the meters. It is shown in this paper that the minimum Hilbert space for the Naimark extension of the Poincare algebra contains six independent fermions, which are identified with left- and right-handed lepton/quark generations. The particle set is the minimum required for realization of the Poincare algebra on commuting operators such that particles can entangle with vacuum meters. The universal nature of fermionic multiplets, existing for both leptons and quarks, motivates the suggestion of a single underlying principle rooted in relativistic quantum mechanics. An elementary particle is interpreted as a complex, entangled quantum system in which the entire space of Poincare observables is realized on commuting operators.
The Naimark extension or embedding is a mathematical description of measurement with a quantum apparatus . In the original construction non-orthogonal projection operators, such as generated by optical coherent states , are extended to orthogonal projection operators in a combined system/meter Hilbert space . As first proposed by von Neumann , and developed by Arthurs and Kelly , a realization of the Naimark extension for position and momentum is obtained by the entangling of a harmonic oscillator with meter harmonic oscillators in ground states. In the Arthurs-Kelly model, measurement on two meters results in the collapse of the system wavefunction to a coherent state corresponding to the measured position and momentum. As explained by Levine and Tucci , entangled simultaneous measurement of position and momentum with a single meter results in the collapse of the system to the eigenstate of the operator that was measured on the system. However, the system/meter expectation values are still proportional to the desired system expectation values. The theory of entangled simultaneous quantum measurement was extended to non-relativistic spin by coupling to spin-1/2 meters by Levine and Tucci . In this case measurements project the system to Bloch states corresponding to the measured spin components. Analogous simultaneous spin measurement schemes are found in Refs. -. Modified Stern-Gerlach experiments, based on hamiltonian models for simultaneous spin measurement, are discussed in Ref. . The relativistic generalization of position, momentum, and angular momentum measurements leads to the consideration of the Poincare algebra of observables and elementary particles. In an attempt to describe elementary particle multiplets, Levine suggested the generalization of spin $`\left(SU\left(2\right)\right)`$ measurement to the measurement of $`SO\left(2n\right)`$ Clifford algebra operators using $`2n`$ spin meters. A more fundamental application to elementary particle multiplets, originating in the structure of the Poincare algebra, is proposed in this paper. Finally, a general discussion of entangled simultaneous measurement with second quantized relativistic fields is found in Ref. .
The mechanism for entangled simultaneous measurement is particularly transparent for second quantized systems. The vacuum, which is defined as the state projected to zero by annihilation operators, is critical to isolate incompatible system observables on independent meters. Another key property of the Naimark extension is that the minimal extension is determined from the pattern of commutativity in the algebra of observables. These properties are demonstrated for non-relativistic position, momentum, and angular momentum measurements in Section 2. In the latter case the three components of angular momentum are measured through the entanglement of the system with two vacuum meters. The formalism for first quantized harmonic oscillators, reviewed in the section, is directly generalized to relativistic fields. Section 3 contains the relativistic generalization of simultaneous measurement for Dirac fermions. It is shown that the Naimark extension of the Poincare algebra consists of six independent fermionic fields. The structure results from the embedding of the three 3-vector operators of momentum, angular momentum, and boost generators. It is suggested that the three entangled lepton/quark generations, an entanglement that is described for left-handed quarks by the Cabibbo-Kobayashi-Maskawa matrix , is a result of the Naimark extension of the Poincare algebra for massive fermions. Appendices A and B contain the formal details of Naimark extensions for second quantized angular momentum and Dirac fields, respectively.
As mentioned above, non-commuting operators are the basis for observed non-locality in quantum systems. The quantum logical consequence of incompatible operators is the non-distributivity of the von Neumann subspace lattice of quantum measurement outcomes . The possibility of a deeper reality, involving a new principle of quantum measurement, motivates the suggestion that measurement of non-commuting observables must be confined to Naimark-extended Hilbert spaces. In Appendix C, it is shown that this principle of entangled simultaneity can be applied to the subspace lattice to avoid paradoxical (non-distributive) statements. As an example, the distributive Naimark-extended lattice is constructed for spin-1/2 measurement. A conclusion follows in Section 4.
## 2 Non-relativistic Systems
In this section the simultaneous measurement of position and momentum in a first quantized harmonic oscillator, and of angular momentum in a second quantized system, are considered. Both cases involve the construction of non-relativistic Naimark extensions, and have properties that generalize to the relativistic case.
An example of a non-relativistic Naimark extension without second quantization uses a pair of one-dimensional harmonic oscillators with the hamiltonian $`\left(\mathrm{}=c=1\right)`$
$$H=\underset{j=1}{\overset{2}{}}\left(\frac{p_j^2}{2m_j}+\frac{1}{2}m_j\omega _j^2q_j^2\right).$$
(1)
In terms of annihilation operators, $`a_j=\sqrt{m_j\omega _j/2}\left(q_j+ip_j/m_j\omega _j\right)`$, the expression in Eq. (1) is written,
$$H=\underset{j}{}\omega _ja_j^{}a_j=\stackrel{}{a}^{}D\stackrel{}{a},$$
(2)
where $`\stackrel{}{a}^{}=(a_1^{},a_2^{})`$,$`D=diag(\omega _1,\omega _2)`$, and the constant zero point energy is dropped. A rotation by angle $`\theta `$, $`\stackrel{}{A}=R\left(\theta \right)\stackrel{}{a}`$, is applied to these operators to define another Hilbert space with combined hamiltonian, $`H=H_0+H_{int}`$, given by
$$H_0=\left(c^2\omega _1+s^2\omega _2\right)A_1^{}A_1+\left(s^2\omega _1+c^2\omega _2\right)A_2^{}A_2,$$
(3)
and
$$H_{int}=sc\left(\omega _1\omega _2\right)\left(A_1^{}A_2+A_2^{}A_1\right),$$
(4)
with $`s=\mathrm{sin}\theta `$ and $`c=\mathrm{cos}\theta `$. Note that, from the commutativity of $`a_1`$ and $`a_2`$, the operators $`A_1`$ and $`A_2`$ commute. Through this simple rotation, the operators $`a_j\left(t\right)=e^{i\omega _jt}a_j\left(0\right)`$, $`j=1,2`$, combine incompatible information about the $`A_1`$ system into operators that are simultaneously measured. For example, consider the commuting operators $`q_1\left(t\right)=\left(a_1\left(t\right)+a_1^{}\left(t\right)\right)/\sqrt{2m_1\omega _1}`$, and $`p_2\left(t\right)=i\sqrt{m_2\omega _2/2}\left(a_2\left(t\right)a_2^{}\left(t\right)\right)`$ for system 1 position and system 2 momentum, respectively. In terms of the operators $`A_1\left(0\right)`$ and $`A_2\left(0\right)`$ at time zero, these operators are given by
$$q_1\left(t\right)=\frac{1}{\sqrt{2m_1\omega _1}}\left(e^{i\omega _1t}\left(cA_1\left(0\right)+sA_2\left(0\right)\right)+e^{i\omega _1t}\left(cA_1\left(0\right)+sA_2\left(0\right)\right)^{}\right),$$
(5)
and
$$p_2\left(t\right)=i\sqrt{m_2\omega _2/2}\left(e^{i\omega _2t}\left(sA_1\left(0\right)cA_2\left(0\right)\right)e^{i\omega _2t}\left(sA_1\left(0\right)cA_2\left(0\right)\right)^{}\right).$$
(6)
The expectation of the operators $`q_1\left(t\right)`$ and $`p_2\left(t\right)`$ in the state $`|\psi >_{A_1}|0>_{A_2}`$, in which the $`A_2`$ system is in the vacuum, is given by
$`<q_1\left(t\right)>`$ $`=`$ $`{\displaystyle \frac{c}{\sqrt{2m_1\omega _1}}}[\mathrm{cos}\left(\omega _1t\right)<\psi |(A_1\left(0\right)+A_1^{}\left(0\right))|\psi >`$ (7)
$`i\mathrm{sin}\left(\omega _1t\right)<\psi \left|(A_1\left(0\right)A_1^{}\left(0\right))\right|\psi >]`$
and
$`<p_2\left(t\right)>`$ $`=`$ $`is\sqrt{{\displaystyle \frac{m_2\omega _2}{2}}}[i\mathrm{sin}\left(\omega _2t\right)<\psi |(A_1\left(0\right)+A_1^{}\left(0\right))|\psi >`$ (8)
$`\mathrm{cos}\left(\omega _2t\right)<\psi \left|(A_1\left(0\right)A_1^{}\left(0\right))\right|\psi >].`$
Inversion of the expressions in Eqs.(7) and (8) for the measurement of $`\left(A_1\left(0\right)+A_1^{}\left(0\right)\right)`$ and $`\left(A_1\left(0\right)A_1^{}\left(0\right)\right)`$ on commuting operators $`q_1\left(t\right)`$ and $`p_2\left(t\right)`$ is given by
$`\left[\begin{array}{c}<A_1\left(0\right)+A_1^{}\left(0\right)>\\ <A_1\left(0\right)A_1^{}\left(0\right)>\end{array}\right]`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{cos}\left(\left(\omega _1\omega _2\right)t\right)}}\left|\begin{array}{cc}\mathrm{cos}\omega _2t& i\mathrm{sin}\omega _1t\\ i\mathrm{sin}\omega _2t& \mathrm{cos}\omega _1t\end{array}\right|\times `$ (16)
$`\left[\begin{array}{c}\left(\sqrt{2m_1\omega _1}/c\right)<q_1\left(t\right)>\\ \left(i\sqrt{2}/\left(\sqrt{m_2\omega _2}s\right)\right)<p_2\left(t\right)>\end{array}\right].`$
The rotation $`\stackrel{}{A}=R\left(\theta \right)\stackrel{}{a}`$ corresponds to the system/meter entanglement that defines a Naimark extension . In this construction, a meter Hilbert space $`_{}`$ is combined with the system space $`_๐ฎ`$ to define the space $`_{}_๐ฎ`$ in which relevant operators commute. For the case of position and momentum, commuting operators can be defined as $`Q_T=Q_1Q_2`$ and $`P_T=P_1+P_2`$, where the subscripts $`1`$ and $`2`$ correspond to system and meter, respectively. Information of the system position and momentum is contained in the orthonormal eigenstates $`\left\{|\xi ,\eta >\right\}`$ of $`(Q_T,P_T)`$ in $`_{}_๐ฎ`$. Mathematically, the projection property of a Naimark extension is defined in terms of density operators $`\rho _{S+M}=|\xi ,\eta ><\xi ,\eta |`$ and $`\rho _M=|0><0|`$, which correspond to Naimark and meter states, respectively. Projection to a system coherent state is then given by
$$|\alpha ><\alpha |=Trace_M\left(\rho _{S+M}\rho _M\right),$$
(17)
where $`\alpha =\sqrt{m\omega /2}\left(q+ip/\left(m\omega \right)\right)`$ for a system of mass $`m`$, position $`q`$, and momentum $`p`$; and where $`Trace_M`$ is the trace over meter states. The projection property in Eq.(17), a defining characteristic of Naimark extensions, is the basis for entangled simultaneous measurement of position and momentum. In a harmonic oscillator model of the measurement, Arthurs and Kelly observed that the entanglement of the system with two vacuum meters resulted in the collapse of the system to a coherent state upon position and momentum measurement on meters. However, the single system/meter entanglement above, $`\stackrel{}{A}=R\left(\theta \right)\stackrel{}{a}`$, is sufficient to simultaneously measure position and momentum. Of course in this case the system is projected onto an eigenstate of the measured operator rather than a coherent state. This realization of the Naimark extension, which is the minimum possible for entangled simultaneous measurement, is used in this paper.
In Appendix A the above construction is extended to a second quantized system described by an angular momentum basis set $`\left\{|jm>;j=0,1,\mathrm{};m=j,\mathrm{},j\right\}`$. It is shown that, by entangling the system to two ancillary meters in vacuum states $`|0>`$, where the vacuum is defined with zero occupation of the states $`\left\{|jm>\right\}`$, the original system angular momentum components are simultaneously measured. Extended operators $`\stackrel{~}{J}_x^{\left(1\right)}`$, $`\stackrel{~}{J}_y^{\left(2\right)}`$, and $`\stackrel{~}{J}_z^{\left(3\right)}`$, corresponding to the original system $`\left(1\right)`$ and meters $`\left(2\right)`$ and $`\left(3\right)`$, have a role similar to $`Q_T`$ and $`P_T`$ defined above. These operators commute and project onto the system operators $`J_x`$, $`J_y`$, and $`J_z`$ upon meter measurement. An analogy to the Naimark projection property in Eq. (17) is also described.
The examples in this section and Appendix A demonstrate the difference between the formulation of entangled simultaneous measurement in first and second quantized systems. In the former case, which has found applications in quantum optics -, the entangling hamiltonian in Eqs. (3) and (4) is dependent on the measured position and momentum operators. In a second quantized system the entangling interaction is independent of the measured operators, which are represented as bilinear functions of creation and annihilation operators (see Eqs.(33)-(35)). Because of this property, the structure of the Naimark extension depends only on the algebra represented by the observables. The minimal extension requires a sufficient number of independent meters for the measurement of non-commuting operator sets. For example, the Naimark extension of the combined galilean algebra of position $`Q_i`$, momentum $`P_i`$, and angular momentum $`J_i`$, $`i=x,y,z`$, appropriate for a non-relativistic particle , is derived from the commutation relations $`j,k,m=x,y,z`$,
$$\begin{array}{ccc}[Q_k,P_j]=i\delta _{kj}\hfill & [J_k,Q_j]=iฯต_{kjm}Q_m\hfill & [J_k,P_j]=iฯต_{kjm}P_m\hfill \\ [J_k,J_j]=iฯต_{kjm}J_m\hfill & [Q_k,Q_j]=0\hfill & [P_k,P_j]=0.\hfill \end{array}$$
A second quantized Naimark extension for the above algebra consists of three distinguishable particles upon which mutually commuting operator pairs $`(J_i,P_i)`$,$`i=x,y,z`$, are measured. All three commuting components of the position vector $`Q_i`$,$`i=x,y,z`$, are measured on a fourth particle. The four entangled particles of the minimum Naimark extension result from the mutually commuting subsets of the algebra in Eq.(2). This construction is generalized to the Poincare algebra in the next section.
## 3 Relativistic Systems
The previous section contains constructions for entangled simultaneous measurement of position, momentum, and angular momentum with quantum meters. The relativistic generalization of these results is the entangled measurement of operators in the Poincare algebra, from which particle representations satisfying Lorentz invariance are obtained . It is suggested that the Naimark extension of this algebra is the basis for a larger structure than the usual particle spin representations, and is the source of the quark/lepton multiplets. Appendix B contains the formulation of entangled simultaneity for fermions described by second quantized Dirac fields. It is shown that $`n`$ non-commuting operators $`\left\{\mathrm{\Theta }_k,k=1,\mathrm{},n\right\}`$ are simultaneously measurable by the system entanglement with $`\left(n1\right)`$ fermions. The Naimark-extended state for this measurement is $`|\varphi >_{\left(1\right)}|0>_{\left(2\right)}\mathrm{}|0>_{\left(n\right)}`$, in which the ancillary particles $`\left(j\right),j=2,\mathrm{},n`$ are in the vacuum state.
The relativistic generalization of position and momentum is the Poincare algebra of operators involving the generators of Lorentz transformations. The nine operators, not including the hamiltonian, can be grouped into three 3-vectors; momentum,
$$P_j=i๐\stackrel{}{x}\psi ^{}(\stackrel{}{x},t)_j\psi (\stackrel{}{x},j),$$
(18)
angular momentum, $`J_k=iฯต_{kjm}J^{jm}`$, $`k,j,m=1,2,3`$, with
$$J^{jm}=i๐\stackrel{}{x}\psi ^{}(\stackrel{}{x},t)\left(x^j^mx^m^ji๐ฅ^{jm}\right)\psi (\stackrel{}{x},t),$$
(19)
and boost generators,
$$K_j=๐\stackrel{}{x}\psi ^{}(\stackrel{}{x},t)\left(it_j+ix_j_0+๐ฆ_j\right)\psi (\stackrel{}{x},t),$$
(20)
with (assuming Pauli matrices $`\sigma _k`$, $`k=1,2,3`$)
$$๐ฅ^{jm}=iฯต^{jmk}\left[\begin{array}{cc}\sigma _k& 0\\ 0& \sigma _k\end{array}\right],$$
(21)
and
$$๐ฆ_j=i/2\left[\begin{array}{cc}\sigma _j& 0\\ 0& \sigma _j\end{array}\right].$$
(22)
The boost operators in Eq.(20) contain the average position of particle energy through the integral $`๐\stackrel{}{x}\stackrel{}{x}\psi ^{}H\psi `$. The commutation relation $`[P_k,K_j]=i\delta _{jk}H`$ further suggests the identification of $`\stackrel{}{K}`$ with position. However the boost operators are not hermitian. Consequently, in order to describe the Poincare algebra with measurable observables, the real and imaginary parts are separately defined as $`K_j=K_j^1+iK_j^2`$ with $`K_j^i`$, $`i=1,2`$, hermitian. The operator $`K_j^1`$ is the physical location of the particle energy, and $`K_j^2`$ reflects spin state effects in the boost operator.
The operators in Eqs.(18)-(21) fit the model in Eq.(45) for which a second quantized Naimark extension is discussed in Appendix B. The complete Poincare algebra in terms of boost operators $`K_j^i`$, $`i=1,2`$, is given by
$`[J_j,J_k]=iฯต_{jkm}J_m,`$ $`[J_j,K_k^1]=iฯต_{jkm}K_m^1,`$ $`[K_i^2,P_j]=0,`$
$`[J_j,P_k]=iฯต_{jkm}P_m,`$ $`[J_j,K_k^2]=iฯต_{jkm}K_m^2,`$ $`[K_i^2,H]=0,`$
$`[K_j^1,H]=iP_j,`$ $`[J_i,H]=0,`$ $`[P_i,H]=0,`$
$`[K_j^1,P_k]=i\delta _{jk}H,`$ $`[K_j^1,K_k^2]+[K_j^2,K_k^1]=0,`$ (23)
$`[K_j^1,K_k^1][K_j^2,K_k^2]=iฯต_{jkm}J_m.`$
A minimal embedding of the Poincare algebra for massive fermions uses the commutations $`[J_j,K_j^1]=[P_j,K_j^2]=0`$,$`j=1,2,3`$, in Eq. (23). Three independent fields $`\psi _j`$ are sufficient for $`(J_j,K_j^1)`$, $`j=1,2,3`$, measurements; and on an additional three fields $`\psi _j^{}`$ the operators $`(P_j,K_j^2)`$, $`j=1,2,3`$, are measured. The hamiltonian operator, satisfying $`[H,P_j]=[H,K_j^2]=0`$, could be measured on any of the $`\psi _j^{}`$ fields. Because of the time derivative in the boost operator in Eq. (20), the energy density position corresponds to the original mass (before entanglement, $`m_j`$ in Eq.(40)) rather than the system mass ($`M_1`$ in Eq.(41)). This suggests that the three fields $`\psi _j`$ used to measure boost operators should have the same mass as the original system, which implies constraints on the entanglement in Eqs. (40)-(42).
The Naimark embedding of the Poincare algebra suggests an underlying explanation for the generation structure of quarks and leptons. The fermions $`(u,c,t)`$, $`(d,s,b)`$, $`(e,\mu ,\tau )`$, and $`(\nu _e,\nu _\mu ,\nu _\tau )`$ are triplets within which particles differ only by mass. A realization of the Naimark extension is obtained by identifying massive quarks and leptons with the fields $`\psi _j\left(\psi _j^{}\right)`$, $`j=1,2,3`$, to obtain the mapping
$$\left[\begin{array}{ccc}f_L^1& f_L^2& f_L^3\\ f_R^1& f_R^2& f_R^3\end{array}\right]\left[\begin{array}{ccc}(J_1,K_1^1)& (J_2,K_2^1)& (J_3,K_3^1)\\ (P_1,K_1^2)& (P_2,K_2^2)& (P_3,K_3^2)\end{array}\right]$$
(24)
where $`f^i=(u,c,t),(d,s,b)`$, and $`(e,\mu ,\tau )`$. The entanglement between different generations of left and right-handed fermions is the result of spontaneous symmetry breaking of the vacuum by Higgs particles. As described in Ref. , the entanglement for quarks results from the diagonalization of the coupling to Higgs particles expressed as rotations $`\stackrel{}{f}_R^{}=W_{u\left(d\right)}\stackrel{}{f}_R`$ and $`\stackrel{}{f}_L^{}=U_{u\left(d\right)}\stackrel{}{f}_L`$, where $`u\left(d\right)`$ corresponds to up (down) quarks, and $`\stackrel{}{f}=(f^1,f^2,f^3)^T`$ corresponds to the triplets $`(u,c,t)`$ and $`(d,s,b)`$. The Cabibbo-Kobayashi-Maskawa matrix, $`V=U_u^{}U_d`$, is the only observable (other than mass) arising from fermionic mixing in the standard model . Mass generation from a non-zero Higgs vacuum expectation provides the entanglement connecting left- and right-handed particles that completes the fermionic Naimark extension.
Massless fermions form representations of a reduced Poincare algebra . For left-handed particles (like massless neutrinos) with fixed $`+1/2`$ helicity, the only measurements required in an inertial frame are boosts $`K_j^1`$ and momentum $`P_j`$ operators. From the commutation $`[P_j,K_i^1]=0`$ for $`ij`$, the operator correspondence of left-handed fermions
$$(\nu _e,\nu _\mu ,\nu _\tau )((P_1,K_2^1),(P_2,K_3^1),(P_3,K_1^1)),$$
(25)
is a sufficient Naimark extension.
## 4 Conclusions
In this paper it is suggested that quantum mechanics is incomplete, and the complete formulation is relevant at relativistic energies. A complete quantum system is in an entangled state with the property that the entire observable phase space of non-commuting operators is measurable.
The definition of a relativistic fermion as a minimally entangled system for the representation of the Poincare algebra is examined. It is shown that the generation structure of leptons and quarks fits this definition with entanglement provided by Higgs couplings as observed in the Cabibbo-Kobayashi-Maskawa matrix. In the high energy limit, the primary structure is not a particle, but rather a system of six entangled particles upon which the complete phase space is represented. This is a generalization of the usual group theoretical particle representation arising from Lorentz invariance of the lagrangian.
## Appendix A Entanglement for Angular Momentum
In this appendix we construct the Naimark extension for angular momentum measurement in a second quantized system described by angular momentum states $`\left\{|jm>;j=0,\mathrm{};m=l,\mathrm{},l\right\}`$ with a non-interacting hamiltonian given by
$$H=\underset{j=0}{\overset{\mathrm{}}{}}\underset{m=j}{\overset{j}{}}\underset{k=1}{\overset{3}{}}E_k\left(j\right)a_{jm}^{\left(k\right)}a_{jm}^{\left(k\right)},$$
(26)
where $`a_{jm}^{\left(k\right)}`$, $`k=1,2,3`$, are the annihilation operators for three independent systems satisfying commutation relations
$`[a_{jm}^{\left(k\right)},a_{j^{}m^{}}^{\left(k^{}\right)}]=\delta _{jj^{}}\delta _{mm^{}}\delta _{kk^{}},`$ $`[a_{jm}^{\left(k\right)},a_{j^{}m^{}}^{\left(k^{}\right)}]=0.`$ (27)
The system/meter operators are defined by a rotation of the original system as
$$\left[\begin{array}{c}A_{jm}^{\left(1\right)}\\ A_{jm}^{\left(2\right)}\\ A_{jm}^{\left(3\right)}\end{array}\right]=R\left[\begin{array}{c}a_{jm}^{\left(1\right)}\\ a_{jm}^{\left(2\right)}\\ a_{jm}^{\left(3\right)}\end{array}\right],$$
(28)
where $`R`$ is a $`|jm>`$-independent rotation matrix. Substitution of Eq. (28) into Eq. (26), with $`E=diag(E_1\left(j\right),E_2\left(j\right),E_3\left(j\right))`$, yields a system/meter hamiltonian given by
$$H=\underset{j=0}{\overset{\mathrm{}}{}}\underset{m=j}{\overset{j}{}}\underset{k=1}{\overset{3}{}}\underset{k^{}=1}{\overset{3}{}}A_{jm}^{\left(k\right)}D_{kk^{}}\left(j\right)A_{jm}^{\left(k^{}\right)},$$
(29)
where
$$D\left(j\right)=RE\left(j\right)R^{}.$$
(30)
The expression in Eq. (29) can be written as a non-interacting hamiltonian (terms with $`k=k^{}`$),
$$H_0=\underset{j=0}{\overset{\mathrm{}}{}}\underset{m=j}{\overset{j}{}}\underset{k=1}{\overset{3}{}}A_{jm}^{\left(k\right)}D_{kk}\left(j\right)A_{jm}^{\left(k\right)},$$
(31)
and an interaction term
$`H_{int}`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=j}{\overset{j}{}}}[D_{12}\left(j\right)A_{jm}^{\left(1\right)}A_{jm}^{\left(2\right)}+D_{13}\left(j\right)A_{jm}^{\left(1\right)}A_{jm}^{\left(3\right)}`$ (32)
$`+`$ $`D_{23}\left(j\right)A_{jm}^{\left(2\right)}A_{jm}^{\left(3\right)}]+h.c.`$
Assume that the operators $`A_{jm}^{\left(1\right)}`$ and $`A_{jm}^{\left(k\right)}`$, $`k=2,3`$, correspond to the system and meters, respectively, and consider the commuting angular momentum operators for the original independent systems in Eq.(26),
$$\stackrel{~}{J}_x^{\left(1\right)}=\underset{j}{}\underset{mm^{}}{}a_{jm}^{\left(1\right)}\left(t\right)a_{jm^{}}^{\left(1\right)}\left(t\right)<jm\left|๐ฅ_x\right|jm^{}>,$$
(33)
$$\stackrel{~}{J}_y^{\left(2\right)}=\underset{j}{}\underset{mm^{}}{}a_{jm}^{\left(2\right)}\left(t\right)a_{jm^{}}^{\left(2\right)}\left(t\right)<jm\left|๐ฅ_y\right|jm^{}>,$$
(34)
and
$$\stackrel{~}{J}_z^{\left(3\right)}=\underset{j}{}\underset{mm^{}}{}a_{jm}^{\left(3\right)}\left(t\right)a_{jm^{}}^{\left(3\right)}\left(t\right)<jm\left|๐ฅ_z\right|jm^{}>,$$
(35)
where $`๐ฅ_\alpha `$, $`\alpha =x,y,z`$, are matrix representations of the angular momentum components. The tilde notation in Eqs. (33)-(35) is used to emphasize the distinction between the operators in the entangled space and the system. The time dependence of $`a_{jm}^{\left(k\right)}`$, given by $`a_{jm}^{\left(k\right)}\left(t\right)=e^{iE_k\left(j\right)t}a_{jm}^{\left(k\right)}\left(0\right)`$, indicates that the operators in Eqs. (33)-(35) are time independent. The substitution of Eq. (28) into Eqs. (33)-(35), and evaluation in the state $`|\psi >_{\left(1\right)}|0>_{\left(2\right)}|0>_{\left(3\right)}`$ in which the meters $`\left(2\right)`$ and $`\left(3\right)`$ are in vacuum states and the system $`\left(1\right)`$ is in the state $`|\psi >`$, results in the expectation values,
$`<\stackrel{~}{J}_x^{\left(1\right)}>`$ $`=`$ $`{\displaystyle \underset{j}{}}{\displaystyle \underset{mm^{}}{}}<\psi \left|A_{jm}^{\left(1\right)}\left(t\right)A_{jm^{}}^{\left(1\right)}\left(t\right)\right|\psi ><jm\left|๐ฅ_x\right|jm^{}>\left|R_{11}\right|^2`$ (36)
$`=`$ $`\left|R_{11}\right|^2<\psi \left|J_x^{\left(1\right)}\right|\psi >,`$
$`<\stackrel{~}{J}_y^{\left(2\right)}>`$ $`=`$ $`{\displaystyle \underset{j}{}}{\displaystyle \underset{mm^{}}{}}<\psi \left|A_{jm}^{\left(1\right)}\left(t\right)A_{jm^{}}^{\left(1\right)}\left(t\right)\right|\psi ><jm\left|๐ฅ_y\right|jm^{}>\left|R_{12}\right|^2`$ (37)
$`=`$ $`\left|R_{12}\right|^2<\psi \left|J_y^{\left(1\right)}\right|\psi >,`$
and
$`<\stackrel{~}{J}_z^{\left(3\right)}>`$ $`=`$ $`{\displaystyle \underset{j}{}}{\displaystyle \underset{mm^{}}{}}<\psi \left|A_{jm}^{\left(1\right)}\left(t\right)A_{jm^{}}^{\left(1\right)}\left(t\right)\right|\psi ><jm\left|๐ฅ_z\right|jm^{}>\left|R_{13}\right|^2`$ (38)
$`=`$ $`\left|R_{13}\right|^2<\psi \left|J_z^{\left(1\right)}\right|\psi >.`$
The expressions in Eqs. (36)-(38) demonstrate the simultaneous measurement of system angular momentum components on the commuting operators $`\stackrel{~}{J}_x^{\left(1\right)}`$, $`\stackrel{~}{J}_y^{\left(2\right)}`$, and $`\stackrel{~}{J}_z^{\left(3\right)}`$. The Naimark projection property corresponding to Eq. (17) is given by
$$J_i^{\left(1\right)}=Trace_{\left(2\right)\left(3\right)}\left[\frac{\stackrel{~}{J}_i^{\left(i\right)}\rho _M}{\left|R_{1i}\right|^2}\right],i=1,2,3,$$
(39)
where $`J_i^{\left(i\right)}`$ is the $`i^{th}`$ component of angular momentum for the system, $`\rho _M=|0><0|`$ is the density matrix for the product vacuum state $`|0>=|0>_{\left(2\right)}|0>_{\left(3\right)}`$, and the trace is over the Hilbert space for meters $`\left(2\right)`$ and $`\left(3\right)`$.
## Appendix B Entanglement for Dirac Fields
Consider the Dirac hamiltonian for $`n`$ independent fermionic fields,
$$H=\underset{j=1}{\overset{n}{}}๐\stackrel{}{x}\left(\pi _j\gamma ^0\stackrel{}{\gamma }\stackrel{}{}\psi _j+m_j\pi _j\gamma ^0\psi _j\right),$$
(40)
where $`\psi _j(\stackrel{}{x},t)`$, $`j=1,\mathrm{},n`$ are second quantized Dirac spinors, $`\gamma ^\mu `$ are Dirac matrices (Ref. notation), and $`\pi _j=i\psi _j^{}`$ is the conjugate momentum to the field $`\psi _j`$. The rotation of the vector of fermionic fields, $`\stackrel{}{\psi }^{}=(\psi _1^{},\mathrm{},\psi _n^{})`$, given by $`\stackrel{}{\mathrm{\Psi }}=R\stackrel{}{\psi }`$, defines the system/meter Hilbert spaces. Substitution into Eq. (40) results in a hamiltonian $`H=H_0+H_{int}`$ with
$$H_0=\underset{j=1}{\overset{n}{}}๐\stackrel{}{x}\left(\mathrm{\Pi }_j\gamma ^0\stackrel{}{\gamma }\stackrel{}{}\mathrm{\Psi }_j+M_j\mathrm{\Pi }_j\gamma ^0\mathrm{\Psi }_j\right),$$
(41)
and
$$H_{int}=๐\stackrel{}{x}\left(\stackrel{}{\mathrm{\Pi }}\gamma ^0\left(diag(M_1,\mathrm{},M_n)\right)\stackrel{}{\mathrm{\Psi }}\right),$$
(42)
with $`\mathrm{\Pi }_j=i\mathrm{\Psi }_j^{}`$, $`m=diag(m_1,\mathrm{},m_n)`$, $`=RmR^{}`$, and $`M_j=_{jj}`$.
The free fields $`\psi _j`$, $`j=1,\mathrm{},n`$, satisfy the equation,
$$\psi _j(\stackrel{}{x},t)=e^{iH_0\left(j\right)t}\psi _j(\stackrel{}{x},0)e^{iH_0\left(j\right)t},$$
(43)
where $`H_0\left(j\right)`$ is the $`j^{th}`$ term in Eq. (40). The second quantized field in Eq. (43) is written in terms of the vector $`\stackrel{}{\mathrm{\Psi }}(\stackrel{}{x},0)`$, at time zero, as
$$\psi _j(\stackrel{}{x},t)=e^{iH_0\left(j\right)t}\left(R^{}\stackrel{}{\mathrm{\Psi }}(\stackrel{}{x},0)\right)_je^{iH_0\left(j\right)t}.$$
(44)
Consider the entangled simultaneous measurement of field operators,
$$\mathrm{\Theta }_k=๐\stackrel{}{x}\psi _1^{}(\stackrel{}{x},t)\theta _k\left(\stackrel{}{x}\right)\psi _1(\stackrel{}{x},t),k=1,\mathrm{},n,$$
(45)
where $`\theta _k\left(\stackrel{}{x}\right)`$ are $`4\times 4`$ $`\stackrel{}{x}`$-dependent operators. The substitution of Eq. (44) into Eq. (45), with $`\psi _1`$ replaced with $`\psi _k`$, results in the expression,
$$\stackrel{~}{\mathrm{\Theta }}_k^{\left(k\right)}\left(t\right)=e^{iH_0\left(k\right)t}๐\stackrel{}{x}\left(\stackrel{}{\mathrm{\Psi }}^{}(\stackrel{}{x},0)R^{}\right)_k\theta _k\left(\stackrel{}{x}\right)\left(R\stackrel{}{\mathrm{\Psi }}(\stackrel{}{x},0)\right)_ke^{iH_0\left(k\right)t}.$$
(46)
Note that the operators $`\stackrel{~}{\mathrm{\Theta }}_k^{\left(k\right)}`$,$`k=1,\mathrm{},n`$, are mutually commuting from the independence of the fields $`\psi _k`$. The commuting set of operators, $`\stackrel{~}{\mathrm{\Theta }}_k^{\left(k\right)}\left(0\right)=e^{iH_0\left(k\right)t}\stackrel{~}{\mathrm{\Theta }}_k^{\left(k\right)}\left(t\right)e^{iH_0\left(k\right)t}`$, evaluated in the state $`|\varphi >_{\left(1\right)}|0>_{\left(2\right)}\mathrm{}|0>_{\left(n\right)}`$, where $`\left(k\right)`$ corresponds to the system with field $`\mathrm{\Psi }_k`$, results in the expression,
$$<\stackrel{~}{\mathrm{\Theta }}_k^{\left(k\right)}\left(0\right)>=<\varphi \left|\mathrm{\Psi }_1^{}(\stackrel{}{x},0)\theta _k\left(\stackrel{}{x}\right)\mathrm{\Psi }_1(\stackrel{}{x},0)\right|\varphi >\left|R_{1k}\right|^2.$$
(47)
The expression in Eq. (47) represents the entangled simultaneous measurement of $`\mathrm{\Theta }_k`$, $`k=1,\mathrm{},n`$, in the state $`|\varphi >_{\left(1\right)}`$ through the Naimark extension to $`n`$ fermions. The equation is a relativistic generalization to an arbitrary set of operators of Eqs. (36)-(38) for non-relativistic angular momentum. The Naimark projection property in Eqs. (17) and (39) is generalized to the relativistic case by the condition
$$\mathrm{\Theta }_k=Trace_{\left(2\right)\left(3\right)\mathrm{}\left(n\right)}\left[\frac{\stackrel{~}{\mathrm{\Theta }}_k^{\left(k\right)}\left(0\right)\rho _M}{\left|R_{1k}\right|^2}\right],k=1,\mathrm{},n,$$
(48)
where $`\rho _M=|0><0|`$ for the meter state $`|0>=|0>_{\left(2\right)}\mathrm{}|0>_{\left(n\right)}`$.
## Appendix C Quantum Logical Implications
Measurement of non-commuting operators in a quantum system is the cause of the non-distributivity in the von Neumann subspace lattice. One approach to the inconsistencies and paradoxes arising from this property is to restrict the interpretation of these measurements as simultaneous only if joint measurements occur on an entangled system at a fixed space-time location. In this appendix, this procedure for the measurement of incompatible observables is included in the subspace lattice of logical propositions. It is shown that the Naimark-extended von Neumann lattice is distributive for a simple example of spin-1/2 component measurement.
### C.1 Entangled Simultaneous Spin Measurement
In this section the entangled simultaneous measurement of spin-1/2 operators $`S_x`$ and $`S_z`$ is demonstrated by the coupling to an independent spin-1/2 meter. This derivation is a simplified version of the mechanism proposed by Levine and Tucci .
Consider two independent spin-1/2 systems $`S`$ and $`R`$ defined by commuting spin-1/2 operators $`S_j`$ and $`R_j`$, $`j=1,2,3`$, which are both expressed as Pauli matrices $`\sigma _j/2`$, $`\left(\mathrm{}=1\right)`$. Assume an entangling hamiltonian given by
$$H=kS_xR_z,$$
(49)
where $`k`$ is an arbitrary coupling constant. The evolution operator $`U=e^{iHt}`$, corresponding to the hamiltonian $`H`$ in Eq. (49), is given by
$$U=\left(c4isS_xR_z\right)$$
(50)
where $`c=\mathrm{cos}\left(kt/4\right)`$ and $`s=\mathrm{sin}\left(kt/4\right)`$. The evolution of operators $`S_z\left(t\right)`$ and $`R_x\left(t\right)`$ in the Heisenberg picture determines the entangled spin component measurement at time $`t`$ on the $`S`$ and $`R`$ systems. The evolved, entangled, and commuting operators are given from Eq. (50) by
$$R_x\left(t\right)=\left[\left(c^2s^2\right)R_x+4csS_xR_y\right],$$
(51)
and
$$S_z\left(t\right)=\left[\left(c^2s^2\right)S_z4csS_yR_z\right].$$
(52)
Assume that the $`R`$-system is aligned along the $`+y`$ direction in the state $`|r>=|+1/2>_y^R`$, and note that the expectations $`<r\left|R_z\right|r>`$ and $`<r\left|R_x\right|r>`$ vanish, to obtain the projection of Eqs. (51) and (52) to $`S`$ system components,
$$<r\left|R_x\left(t\right)\right|r>=2scS_x,$$
(53)
and
$$<r\left|S_z\left(t\right)\right|r>=\left(c^2s^2\right)S_z.$$
(54)
The operators $`R_x\left(t\right)`$ and $`S_z\left(t\right)`$ provide an entangled simultaneous measurement of $`S_x`$ and $`S_z`$. The above projections are typically expressed as a partial trace over the meter Hilbert space of the product of the meter density operator $`\rho _M=|r><r|`$ and the relevant extended-space operator . For example, Eqs. (53) and (54) are expressed as
$$S_x=\frac{1}{2sc}Trace_R\left[R_x\left(t\right)\rho _M\right],$$
(55)
and
$$S_z=\frac{1}{\left(c^2s^2\right)}Trace_R\left[S_z\left(t\right)\rho _M\right].$$
(56)
Mathematically, the $`(S,R)`$-Hilbert space $`_S_R`$ is the Naimark extension of the $`S`$-Hilbert space $`_S`$ -. Among the proposals in this paper is that the Naimark extension is a logically consistent mechanism for joint quantum measurements of non-commuting operators.
### C.2 Extended Subspace Lattice
The essential property of an extended quantum system is seen in the lattice of spin-1/2 measurements, which is defined in Refs. and . Consider a spin-1/2 system $`S`$ and an apparatus that measures spin in the $`x`$ or $`z`$ directions with the result up $`u`$ or down $`d`$. As discussed in the previous section, a spin-1/2 meter $`R`$ is entangled with $`S`$ such that the $`R`$-vacuum expectation value directly reveals the corresponding spin operator of $`S`$. Denote by $`\alpha \beta \gamma `$, with $`\alpha \{r,s\}`$, $`\beta \{x,z\}`$, and $`\gamma \{u,d\}`$, the projected subspace for the measurement on the system $`\alpha `$ of spin value $`\gamma `$ along direction $`\beta `$. For example, $`sxu`$ defines the subspace generated by the $`S_x`$ eigenvector with eigenvalue $`u`$ in system $`S`$. A subspace lattice of these outcomes on the system $`S`$ is shown in Fig. 1, where $`()`$ represents or (and) connectives on the measurement outcomes . The bold lines in Fig. 1 correspond to measurement outcomes, and the connecting dotted lines denote inclusion into a higher dimensional subspace. The horizontal/vertical and $`45^{}`$-rotated coordinate systems correspond to $`S_x`$ and $`S_z`$ measurements (either $`u`$ or $`d`$), respectively.
A defining feature of quantum lattices is a failure of the distributive property of meet $`\left(m\right)`$ over join $`\left(j\right)`$, which is seen in Fig. 1 by the different outcomes in
$$sxum\left(szujszd\right)=sxu,$$
(57)
and
$$\left(sxumszu\right)j\left(sxumszd\right)=\left\{s\right\}j\left\{s\right\}=\left\{s\right\},$$
(58)
where $`m`$ and $`j`$ denote meet and join operations on the lattice, and the notation $`\left\{\alpha \right\}=\left(\alpha xu\alpha xd\right)=\left(\alpha zu\alpha zd\right)`$ and $`\left[\alpha \right]=\left(\alpha xu\alpha xd\right)=\left(\alpha zu\alpha zd\right)`$ is used. The failure of the distributive property is due to the fact that the observables $`S_x`$ and $`S_z`$ are incompatible; suggesting that the distributive property will hold in an extended lattice with measurements on commuting operators.
The extended Naimark subspace lattice, which combines the $`S`$ and $`R`$ Hilbert spaces, is shown in Fig. 2. The vertically-placed coordinates in each row correspond to $`S`$ (top) and $`R`$ (bottom) Hilbert spaces that have been entangled to provide simultaneous measurement of system $`S`$ components. A mechanism for the entangled simultaneous measurement, with $`R`$ in a special initial state as a meter, was discussed in Section C.1. Consider a section of the extended lattice referring to the simultaneous measurement of the $`S`$ system spin along the $`x`$ and $`z`$ directions. The relevant lattice operations
$`\left(sxu\left[r\right]\right)m\left[\left(szu\left[r\right]\right)j\left(rzd\left[s\right]\right)\right]`$ $`=`$ $`\left(sxu\left[r\right]\right)m\left(\left[r\right]\left[s\right]\right)`$ (59)
$`=`$ $`sxu\left[r\right],`$
and
$$\begin{array}{ccc}\left[\left(sxu\left[r\right]\right)m\left(szu\left[r\right]\right)\right]j\left[\left(sxu\left[r\right]\right)m\left(rzd\left[s\right]\right)\right]\hfill & =\hfill & \\ \left(\left\{s\right\}\left[r\right]\right)j\left(sxurzd\right)=sxu\left[r\right],\hfill & & \end{array}$$
(60)
are distributive. All other complex statements mixing $`S_x`$ on $`S`$ and $`S_z`$ on $`R`$ are distributive in the lattice. For example, another distributive expression is given by
$$\begin{array}{ccc}\left(sxu\left[r\right]\right)m\left[\left(\left\{s\right\}rzu\right)j\left(\left\{s\right\}rzd\right)\right]=\hfill & & \\ \left[\left(sxu\left[r\right]\right)m\left(\left\{s\right\}rzu\right)\right]j\left[\left(sxu\left[r\right]\right)m\left(\left\{s\right\}rzd\right)\right]\hfill & =\hfill & \\ \left(\left\{s\right\}\left[r\right]\right).\hfill & & \end{array}$$
(61)
In this appendix a distributive quantum theory is defined using special measurement-dependent statements in cases of incompatible observables. The appropriate measurements involve meters in special-purpose initial states that are entangled with the system. In the body of the paper it was suggested that such measurements are fundamental in the definition of relativistic particle observables. The special meter initial state for particles is the vacuum, the Naimark extension is the repeated fermionic generations, and entanglement is observed in the Cabibbo-Kobayashi-Maskawa matrix and spontaneous breaking of the electro-weak gauge symmetry.
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# 1 Introduction
## 1 Introduction
The synthesis of light nuclei in the early universe, the Big Bang Nucleosynthesis (BBN), represents one of the most striking evidence in favour of standard cosmology, and since its proposal , it has been extensively used as one of the best laboratories where to test cosmological models and/or elementary particle physics. The appealing feature of BBN is that, in its standard version, it relies on quite solid theoretical grounds, which makes the predictions for $`D`$, $`{}_{}{}^{3}He`$, $`{}_{}{}^{4}He`$ and $`{}_{}{}^{7}Li`$ abundances, whose primordial values are only partially modified by the subsequent stellar activity, quite robust.
Recently, the experimental accuracy in measurements of the light primordial nuclide abundances, mainly $`{}_{}{}^{4}He`$, has been highly improved, reaching a precision of the order of $`1\%`$. Similar improvements have been also obtained for both Deuterium ($`D`$) and $`{}_{}{}^{7}Li`$ relative abundances, $`Y_DD/H`$ and $`Y_{{}_{}{}^{7}Li}{}_{}{}^{7}Li/H`$, but, unfortunately, the refinement of the experimental techniques does not yet correspond to a clear picture of the primordial nuclide densities. This is mainly due to an uncomplete understanding of systematic errors. In particular, measuring the primordial $`{}_{}{}^{4}He`$ mass fraction, $`Y_p`$, from regression to zero metallicity in Blue Compact Galaxies, two independent surveys obtained two results, a low value ,
$$Y_p^{(l)}=0.234\pm 0.003,$$
(1.1)
and a sensibly larger one ,
$$Y_p^{(h)}=0.244\pm 0.002,$$
(1.2)
which are compatible at $`2\sigma `$ level only.
As in a recent analysis , we here adopt a more conservative value, with a larger error (hereafter we always use 1$`\sigma `$ errors)
$$Y_p=0.238\pm 0.005.$$
(1.3)
A similar dichotomy holds in $`D`$ measurements as well, where the study of distant Quasars Absorption line Systems (QAS), is thought to represent a reliable way to estimate the primordial Deuterium. In this case, observations in different QAS leads to the incompatible results
$`Y_D^{(l)}`$ $`=`$ $`\left(3.4\pm 0.3\right)10^5,`$ (1.4)
$`Y_D^{(h)}`$ $`=`$ $`\left(2.0\pm 0.5\right)10^4.`$ (1.5)
Finally, a reliable estimate for $`{}_{}{}^{7}Li`$ primordial abundance is provided by the Spite plateau, observed in the halo of POP II stars . The observations give the primordial abundance ,
$$Y_{{}_{}{}^{7}Li}=\left(1.73\pm 0.21\right)10^{10}.$$
(1.6)
From the theoretical point of view, the BBN predictions are obtained by numerically solving a set of coupled Boltzmann equations, which trace the abundances of the different nuclides in the framework of standard Big Bang cosmology . The collisional integrals of the above equations contain all $`np`$ weak reaction rates and a large nuclear reaction network .
The increasing precision in measuring the primordial abundance has recently pushed the theoretical community to make an effort to develop new generation BBN codes , with a comparable level of accuracy. In a recent series of papers, in the framework of the standard cosmological model, the present authors and other groups have performed a comprehensive and accurate analysis of all the physical effects which influence $`{}_{}{}^{4}He`$ mass fraction up to $`0.1\%`$.
Since almost all neutrons present at the onset of nucleosynthesis are fixed into $`{}_{}{}^{4}He`$, its abundance is mainly function of the neutron versus proton abundances, at the time of $`np`$ weak interactions freeze out, which takes place for $`T1MeV`$. To improve the accuracy on the $`{}_{}{}^{4}He`$ prediction it is demanding to reach an accuracy level of the order of $`1\%`$ in the estimate of the $`np`$ rates, well beyond the simple Born approximation. This has been performed by considering a number of additional contributions, which, ordered according to their relative weight, are the following:
i) electromagnetic radiative and Coulomb corrections ;
ii) finite nucleon mass corrections ;
iii) thermal radiative effects induced by the presence of a surrounding plasma of $`e^\pm `$ and $`\gamma `$ ;
iv) corrections to the equation of state of the $`e^\pm ,\gamma `$ plasma due to thermal mass renormalization ;
v) the residual neutrino coupling to the plasma during $`e^+e^{}`$ annihilation, affecting the neutrino to photon temperature ratio .
Unfortunately, a similar systematic analysis of all corrections up to the desired level of accuracy cannot be performed for the other input parameters of the theory. The theoretical estimates do depend in fact on the values of the neutron lifetime, as well as on several nuclear reaction cross sections, which are poorly known in the energy range relevant for BBN. This introduces a certain level of uncertainty on the light element abundances. This aspect has been studied in two different approaches, either using Monte Carlo methods to sample the error distributions of the relevant reaction cross sections, or, alternatively, using a linear error propagation , with comparable results.
In this paper, we further refine our previous predictions by using a new version of our numerical code, where the full dependence of the BBN equations on the electron chemical potential is accurately implemented. Furthermore, for the standard scenario (vanishing neutrino chemical potentials) we perform an accurate likelihood analysis in the space of the two free parameters of the model, the effective number of neutrinos, $`N_\nu `$, and the baryon to photon ratio $`\eta `$.
While the value of electron chemical potential is bounded, by neutrality, by the value of $`\eta `$, this is not the case for neutrinoโantineutrino asymmetries which, in principle, can be quite large. The influence of the neutrino chemical potentials on BBN predictions (degenerate BBN) has been considered in the past . One relevant aspect of this analysis, which is worth stressing, is that degenerate BBN allows for a better agreement with observations at values of $`\eta `$ larger than, say, $`610^{10}`$, while standard BBN prefers smaller values, $`\eta (2รท6)10^{10}`$. We will discuss this in detail in the paper. This feature is particularly relevant in view of the recent analysis of the BOOMERANG and MAXIMA-1 results on the acoustic peak of the Cosmic Microwave Background Radiation (CMBR), which seems to favour a value for the baryonic asymmetry $`\eta `$ of the order of $`10^9`$, larger, as we said, than what expected in the framework of standard BBN. This discrepancy may be looked as a signal in favour of a large neutrinoโantineutrino asymmetry. It seems therefore demanding to reโanalyze the degenerate scenario, making profit of the now available more precise BBN codes. We have performed this study, and the main result is that degenerate BBN and CMBR data seems to be compatible for large values of the effective number of neutrinos, $`N_\nu 10`$ and $`\eta 10^9`$.
The paper is organized as follows. In section 2 we briefly describe the BBN set of differential equations, recast in a suitable form for a numerical solution. The light element abundances obtained from our code, for standard BBN, are then presented in section 3 as a functions of $`\tau _n`$, $`N_\nu `$ and $`\eta `$. Section 4 is devoted to degenerate BBN and to the bounds on neutrino chemical potentials coming from nucleosynthesis and CMBR data. Finally, in section 4, we give our conclusions.
## 2 The BBN set of equations
Consider $`N_{nuc}`$ species of nuclides, whose number densities, $`n_i`$, are normalized with respect to the total number density of baryons, $`n_B`$,<sup>1</sup><sup>1</sup>1We will use the same notations of our previous paper , which we refer to for further details.
$$X_i=\frac{n_i}{n_B}i=1,..,N_{nuc}.$$
(2.1)
Alternatively, we will also make use in the following of the notation:
$`X_1=X_n,`$ $`X_2=X_p=X_H,`$ $`X_3=X_D,`$
$`X_5=X_{{}_{}{}^{3}He},`$ $`X_6=X_{{}_{}{}^{4}He},`$ $`X_8=X_{{}_{}{}^{7}Li}.`$ (2.2)
The set of differential equations ruling primordial nucleosynthesis is given by :
$`{\displaystyle \frac{\dot{R}}{R}}`$ $`=`$ $`H=\sqrt{{\displaystyle \frac{8\pi }{3M_P^2}}\rho _T},`$ (2.3)
$`{\displaystyle \frac{\dot{n}_B}{n_B}}`$ $`=`$ $`\mathrm{\hspace{0.17em}3}H,`$ (2.4)
$`\dot{\rho _T}`$ $`=`$ $`\mathrm{\hspace{0.17em}3}H(\rho _T+\mathrm{p}_T),`$ (2.5)
$`\dot{X}_i`$ $`=`$ $`{\displaystyle \underset{j,k,l}{}}N_i\left(\mathrm{\Gamma }_{klij}{\displaystyle \frac{X_l^{N_l}X_k^{N_k}}{N_l!N_k!}}\mathrm{\Gamma }_{ijkl}{\displaystyle \frac{X_i^{N_i}X_j^{N_j}}{N_i!N_j!}}\right)\mathrm{\Gamma }_i(X_j),`$ (2.6)
$`L({\displaystyle \frac{m_e}{T}},\varphi _e)`$ $`=`$ $`{\displaystyle \frac{n_B}{T^3}}{\displaystyle \underset{j}{}}Z_jX_j,`$ (2.7)
where $`\rho _T`$ and $`\mathrm{p}_T`$ denote the total energy density and pressure, respectively,
$`\rho _T`$ $`=`$ $`\rho _\gamma +\rho _e+\rho _\nu +\rho _B\rho _{NB}+\rho _B,`$ (2.8)
$`\mathrm{p}_T`$ $`=`$ $`\mathrm{p}_\gamma +\mathrm{p}_e+\mathrm{p}_\nu +\mathrm{p}_B,`$ (2.9)
$`i,j,k,l=(1,..,N_{nuc})`$, $`Z_i`$ is the charge number of the $`i`$th nuclide, and the function $`L(z,y)`$ is defined as
$$L(z,y)\frac{1}{\pi ^2}_z^{\mathrm{}}๐xx\sqrt{x^2z^2}\left(\frac{1}{e^{xy}+1}\frac{1}{e^{x+y}+1}\right).$$
(2.10)
Eq.(2.3) is the definition of the Hubble parameter, $`H`$, whereas Eq.s (2.4) and (2.5) state the total baryon number and entropy conservation in the comoving volume, respectively. The set of $`N_{nuc}`$ Boltzmann equations (2.6) describe the density evolution of each nuclide specie, and finally Eq.(2.7) states the universe charge neutrality in terms of the electron chemical potential, $`\varphi _e\mu _e/T`$, with $`T`$ the temperature of $`e^\pm ,\gamma `$ plasma. Note that the neutrino energy density and pressure are included in Eq. (2.5) only for $`TT_D`$<sup>2</sup><sup>2</sup>2We assume that all neutrinos decouple at the same temperature, $`T_D=2.3MeV`$ ..
In a previous analysis, , we neglected the contribution of $`\varphi _e`$ in the BBN equations, since its effect results to be very small. In this way we obtained a substantial simplification of the set of equations (2.3)-(2.7), since the unknown functions can be reduced in this case to the $`N_{nuc}+1`$ ($`\widehat{h}n_B/T^3,X_j`$). We do here release this assumption and report the results of an improved code where we take into complete account the electron chemical potential evolution. We will see, however, that all changes on the final abundances are of minor impact. In order to obtain the new equations, we note that it is more convenient to follow the evolution of the $`N_{nuc}+1`$ unknown functions $`(\varphi _e,X_j)`$ in terms of the dimensionless variable $`z=m_e/T`$, and to use Eq. (2.7) to get $`n_B`$ as a function of $`\varphi _e`$. The new set of differential equations may be cast in the form
$$\frac{d\varphi _e}{dz}=\frac{1}{z}\frac{LEF+(zL_z3L)G}{LE\frac{\delta \widehat{\rho }_e}{\delta \varphi _e}L_{\varphi _e}G},$$
(2.11)
$$\frac{dX_i}{dz}=\frac{\widehat{\mathrm{\Gamma }}_i}{z}\frac{L_{\varphi _e}F+(zL_z3L)\frac{\delta \widehat{\rho }_e}{\delta \varphi _e}}{LE\frac{\delta \widehat{\rho }_e}{\delta \varphi _e}L_{\varphi _e}G},$$
(2.12)
where the functions $`E`$, $`F`$ and $`G`$ are given by
$`E(z,\varphi _e,X_j)`$ $`=`$ $`3\widehat{H}{\displaystyle \frac{{\displaystyle \underset{i}{}}Z_i\widehat{\mathrm{\Gamma }}_i}{{\displaystyle \underset{j}{}}Z_jX_j}},`$ (2.13)
$`F(z,\varphi _e,X_j)`$ $`=`$ $`4\widehat{\rho }_{NB}+{\displaystyle \frac{3}{2}}\widehat{\mathrm{p}}_Bz{\displaystyle \frac{\delta \widehat{\rho }_e}{\delta z}},`$ (2.14)
$`G(z,\varphi _e,X_j)`$ $`=`$ $`3\widehat{H}(\widehat{\rho }_{NB}+\widehat{\mathrm{p}})+{\displaystyle \frac{zL}{_jZ_jX_j}}{\displaystyle \underset{i}{}}\left(\mathrm{\Delta }\widehat{M}_i+{\displaystyle \frac{3}{2z}}\right)\widehat{\mathrm{\Gamma }}_i,`$ (2.15)
and $`Hm_e\widehat{H}`$, $`n_Bm_e^3\widehat{n}_B`$, $`\mathrm{\Gamma }_im_e\widehat{\mathrm{\Gamma }}_i`$, $`\rho _TT^4\widehat{\rho }_T`$, $`\mathrm{p}_TT^4\widehat{\mathrm{p}}_T`$. Note that by using Eq. (2.7) and the previous definitions, it is possible to express $`\widehat{\rho }_B`$ and $`\widehat{\mathrm{p}}_B`$ as functions of $`z`$, $`\varphi _e`$, and $`X_i`$ only
$`\widehat{\rho }_B`$ $`=`$ $`{\displaystyle \frac{zL(z,\varphi _e)}{_jZ_jX_j}}\left[\widehat{M}_u+{\displaystyle \underset{j}{}}\left(\mathrm{\Delta }\widehat{M}_j+{\displaystyle \frac{3}{2z}}\right)X_j\right],`$ (2.16)
$`\widehat{\mathrm{p}}_B`$ $`=`$ $`{\displaystyle \frac{L(z,\varphi _e)}{_jZ_jX_j}}{\displaystyle \underset{j}{}}X_j.`$ (2.17)
With $`\mathrm{\Delta }\widehat{M}_i`$ and $`\widehat{M}_u`$ we denote the i-th nuclide mass excess and the atomic mass unit, respectively, normalized to $`m_e`$. In Appendix A we report the partial derivative of $`L`$ with respect to $`z`$ and $`\varphi _e`$, denoted with $`L_z`$ and $`L_{\varphi _e}`$, and the quantities $`\widehat{\rho }_e`$, $`\delta \widehat{\rho }_e/\delta z`$ and $`\delta \widehat{\rho }_e/\delta \varphi _e`$ in a form which is suitable for a BBN code implementation.
Eq.s (2.11)-(2.12) are solved by imposing the following initial conditions at $`z_{in}=m_e/(10MeV)`$:
$`\varphi _e(z_{in})`$ $`=`$ $`\varphi _{e}^{}{}_{}{}^{0},`$ (2.18)
$`X_n(z_{in})`$ $`=`$ $`\left(\mathrm{exp}\{\widehat{q}z_{in}\}+1\right)^1,X_p(z_{in})=\left(\mathrm{exp}\{\widehat{q}z_{in}\}+1\right)^1,`$ (2.19)
$`X_i(z_{in})`$ $`=`$ $`{\displaystyle \frac{g_i}{2}}\left(\zeta (3)\sqrt{{\displaystyle \frac{8}{\pi }}}\right)^{A_i1}A_i^{\frac{3}{2}}\left({\displaystyle \frac{m_e}{M_Nz_{in}}}\right)^{\frac{3}{2}(A_i1)}\eta ^{A_i1}X_p^{Z_i}X_n^{A_iZ_i}`$ (2.20)
$`\times `$ $`\mathrm{exp}\left\{\widehat{B}_iz_{in}\right\}\text{with}i=3,..,N_{nuc}.`$
In the previous equations $`\widehat{q}=(M_nM_p)/m_e`$, and the quantities $`A_i`$ and $`\widehat{B}_i`$ denote the atomic number and the binding energy of the $`i`$th nuclide normalized to electron mass, respectively. Finally $`\eta `$ is, as usual, the baryon to photon number density ratio, and $`\varphi _{e}^{}{}_{}{}^{0}`$ the solution of the implicit equation
$$L(z_{in},\varphi _{e}^{}{}_{}{}^{0})=\frac{11}{4}\frac{2\zeta (3)}{\pi ^2}\eta \underset{i}{}Z_iX_i(z_{in}).$$
(2.21)
The method of resolution of the BBN equations (2.11)-(2.12) is the same applied in . It is the Backward Differentiation Formulas with Newtonโs method, implemented in a NAG routine with adaptive step-size (see for more details). We used the reduced network of nuclear reactions, made of 25 reactions involving 9 nuclides, since the use of the complete network affects the abundances for no more than $`0.01\%`$.
## 3 Primordial abundances for standard BBN
The new BBN code has been used to produce the primordial abundances, in the standard scenario, for different values of the input parameters, namely the neutron lifetime $`\tau _n`$, the effective number of neutrinos $`N_\nu `$, defined as
$$\rho _\nu =N_\nu \frac{7}{4}\frac{\pi ^2}{30}T_\nu ^4,$$
(3.1)
with $`\rho _\nu `$ the total neutrino energy densities, and the final baryon to photon number density ratio $`\eta `$.
We consider the abundances relative to hydrogen for $`D`$, $`{}_{}{}^{3}He`$, and $`{}_{}{}^{7}Li`$,
$$Y_D=\frac{X_D}{X_H},Y_{{}_{}{}^{3}He}=\frac{X_{{}_{}{}^{3}He}}{X_H},Y_{{}_{}{}^{7}Li}=\frac{X_{{}_{}{}^{7}Li}}{X_H}.$$
(3.2)
In the case of $`{}_{}{}^{4}He`$, the quantity usually defined as mass fraction,
$$Y_p=\frac{A_{{}_{}{}^{4}He}X_{{}_{}{}^{4}He}}{_jA_jX_j},$$
(3.3)
is rather the baryon number fraction, since $`A_{{}_{}{}^{4}He}=4`$. Note that expression (3.3) does not correspond to the true mass fraction, obviously defined as
$$Y_p^m=\frac{M_{{}_{}{}^{4}He}X_{{}_{}{}^{4}He}}{_jM_jX_j}.$$
(3.4)
The difference between $`Y_p`$ and $`Y_p^m`$ are of the order of $`1\%`$ and thus relevant for an accurate analysis like the one presented here. Since it is customary to express the experimental value in terms of $`Y_p`$, see (1.3), we will consider this quantity for a comparison with experimental data.
Table 1 shows the results obtained for $`N_\nu =3`$, $`\tau _n=886.7s`$, and $`\eta =510^{10}`$, compared with our previous results . As one can see, the inclusion of the complete evolution of $`\varphi _e`$ modifies the final $`{}_{}{}^{4}He`$ abundances for less than $`0.1\%`$ and for few percent the other nuclides.
In Figures 1, 3, 5, 7 we show the theoretical predictions for the abundances (3.2)-(3.3), compared with the experimental values. Figures 2, 4, 6, 8 show the comparison between our results and the ones of Ref.s . The relative differences among our predictions and the results of Ref.s for $`{}_{}{}^{4}He`$, in the relevant range for $`\eta _{10}10^{10}\eta `$, are less than $`0.25\%`$, and thus probably due to different ways of taking into account subdominant effects (thermal radiative corrections). The differences are not much larger for Deuterium, but reach few percents for $`{}_{}{}^{3}He`$ and $`{}_{}{}^{7}Li`$, which however have very large theoretical errors due to the uncertainties on nuclear reaction rates. Note that for $`D`$, $`{}_{}{}^{3}He`$ and $`{}_{}{}^{7}Li`$ only the results of are available.
We have performed a fit of the previous abundances as functions of $`N_\nu `$, $`\tau _n`$, and $`x\mathrm{log}_{10}\eta _{10}`$, with an accuracy which is better than $`1\%`$ in the all ranges $`0.5N_\nu 6`$, $`882.9s\tau _n890.5s`$, $`0x1`$. In particular, for $`N_\nu =3`$ and $`{}_{}{}^{4}He`$, our fit is accurate at $`0.01\%`$. The fitting functions have been chosen as
$`k_iY_i(N_\nu ,\tau _n,x)`$ $`=`$ $`[{\displaystyle \underset{j=0}{\overset{8}{}}}a_jx^j+\left({\displaystyle \underset{j=0}{\overset{8}{}}}b_jx^j\right)(\tau _n\tau _n^{ex})+\left({\displaystyle \underset{j=0}{\overset{8}{}}}c_jx^j\right)(N_\nu 3)+`$ (3.5)
$`\left({\displaystyle \underset{j=0}{\overset{8}{}}}d_jx^j\right)(N_\nu 3)^2\left]\mathrm{exp}\right\{{\displaystyle \underset{j=1}{\overset{6}{}}}e_jx^j\},`$
with all coefficients given in Table 2.
By using these expressions (3.5), as the theoretical predictions for the light element abundances, and their experimental measurements (1.3), (1.4), (1.5), (1.6) it is possible to test the compatibility of standard BBN scenario in the $`N_\nu `$$`\eta _{10}`$ plane. To this end, we define a total likelihood function as
$$(N_\nu ,\eta _{10})=L_D(N_\nu ,\eta _{10})L_{{}_{}{}^{4}He}(N_\nu ,\eta _{10})L_{{}_{}{}^{7}Li}(N_\nu ,\eta _{10}),$$
(3.6)
where the likelihood function for each abundance, assuming Gaussian distribution for the errors, is given by the overlap
$`L_i(N_\nu ,\eta _{10})=`$ (3.7)
$`={\displaystyle \frac{1}{2\pi \sigma _i^{th}(N_\nu ,\eta _{10})\sigma _i^{ex}}}`$ $`{\displaystyle ๐Y\mathrm{exp}\left\{\frac{(YY_i^{th}(N_\nu ,\eta _{10}))^2}{2\sigma _i^{th\mathrm{\hspace{0.17em}2}}(N_\nu ,\eta _{10})}\right\}\mathrm{exp}\left\{\frac{(YY_i^{ex})^2}{2\sigma _i^{ex\mathrm{\hspace{0.17em}2}}}\right\}}`$
In order to evaluate $`L_i(N_\nu ,\eta _{10})`$ we need the theoretical uncertainties $`\sigma _i^{th}(N_\nu ,\eta _{10})`$. In Ref. a new, alternative approach to the standard Monte Carlo technique has been proposed. This method is based on linear error propagation for the estimate of theoretical uncertainties on primordial abundances due to the poor knowledge of nuclear reaction rates. The light element abundances, $`Y_i`$, and the logarithmic derivatives of $`Y_i`$ with respect to the nuclear rates, $`\lambda _{ik}`$, are given as polynomial fits, while the variation of the abundances for a change $`\delta R_k`$ in the rate $`R_k`$ is obtained as
$$\delta Y_i=Y_i\underset{k}{}\lambda _{ik}\frac{\delta R_k}{R_k}.$$
(3.8)
Taking into account the error correlation, the error matrix results
$$\sigma _{ij}^2=Y_iY_j\underset{k}{}\lambda _{ik}\lambda _{jk}\left(\frac{\mathrm{\Delta }R_k}{R_k}\right)^2,$$
(3.9)
where $`\mathrm{\Delta }R_k`$ are the $`1\sigma `$ uncertainties. In particular the theoretical $`1\sigma `$ uncertainties are given by the square root of the diagonal elements,
$$\sigma _i^{th}=\sqrt{\sigma _{ii}^2}.$$
(3.10)
We used the Fortran code provided by the authors to calculate the theoretical uncertainties $`\sigma _i^{th}(N_\nu ,\eta _{10})`$, which we used in Eq. (3.7).
The total likelihood function (3.6) and the corresponding contour plots for $`50\%`$, $`68\%`$ and $`95\%`$ CL, for low and high $`D`$ are shown in Figures 9, 10 and 11, 12, respectively. As already clear from Figure 1, the two different experimental estimates of $`D`$ single out different regions for $`\eta _{10}`$. From the $`95\%`$ CL contour of Figure 10, for low $`D`$ we have $`4.0\eta _{10}5.7`$, which is in fair agreement with the similar results of Ref. , whereas for high $`D`$ from Figure 12 we get $`1.4\eta _{10}3.7`$.
As far as $`N_\nu `$ is concerned, for low $`D`$ and for $`95\%`$ CL one has $`1.7N_\nu 3.3`$, whereas for high $`D`$ one gets $`2.3N_\nu 4.4`$. In both cases we obtain comparable ranges for $`N_\nu `$ with respect to Ref. . The total likelihood function is peaked around the points shown as the crosses in Figures 10 and 12, which correspond to $`N_\nu =2.44`$ and $`\eta _{10}=4.69`$ for low $`D`$, and to $`N_\nu =3.29`$ and $`\eta _{10}=1.81`$ for high $`D`$. The position of the two maxima is easily understood. The Deuterium abundance is a decreasing function of $`\eta _{10}`$. Since lowering $`\eta _{10}`$ results in a smaller $`{}_{}{}^{4}He`$ mass fraction $`Y_p`$, it is necessary to compensate this effect by increasing the universe expansion rate via a larger $`N_\nu `$. This in fact leads to a larger value for the freezeโout temperature for nucleon weak interactions and thus gives a larger amount of the initial neutron to proton density ratio. The single contributions (3.7) to the total likelihood (3.6) can be easily recognized by looking at Figures 13-16, where we show the $`D`$, $`{}_{}{}^{4}He`$ and $`{}_{}{}^{7}Li`$ likelihood functions, for both high and low Deuterium results, for the two preferred values $`N_\nu =2.44`$ and $`N_\nu =3.29`$. For the low $`D`$ case, a better overlap between the maxima of the single likelihoods of $`D`$ and $`{}_{}{}^{4}He`$ is realized for $`N_\nu =2.44`$. The opposite situation occurs for high $`D`$ where to $`N_\nu =3.29`$ corresponds a better overlap of the single likelihoods.
## 4 BBN predictions for degenerate neutrinos
The assumption of vanishing neutrino chemical potentials can be only justified by the sake of simplicity or by the theoretical prejudice that their order of magnitude should be set by the ones of the corresponding charged leptonic partner. However, physical scenarios in which large lepton asymmetries are produced, which do not lead to a large baryon asymmetry, have been proposed in literature . They are based on the Affleck-Dine mechanism or on activeโsterile neutrino oscillations . In particular, the expected asymmetry may be different for each neutrino family. For these reasons, at least in principle, it is worth-while considering the primordial nucleosynthesis in presence of large neutrino asymmetries, i.e. for non-vanishing neutrino chemical potentials. It seems to us that this topic, which has been extensively studied in past , receives a renewed interest in view of the recent BOOMERANG and MAXIMA-1 results on the acoustic peaks of the CMBR, which suggests a larger value of the baryonic matter contribution to the total energy density $`\mathrm{\Omega }`$. In their analysis of the data they find a baryon density $`\mathrm{\Omega }_bh^20.02รท0.03`$, or $`\eta (6รท10)10^{10}`$, which, as is clear from the considerations of the previous section, is incompatible at $`95\%`$ CL with the standard BBN result (see Figures 10 and 12). Actually even larger values for $`\mathrm{\Omega }_bh^2`$ are obtained if no constraints are imposed in the likelihood analysis (see ). It is the aim of this section to discuss whether a finite neutrino chemical potential may reconcile BBN theoretical predictions, the observed nuclide abundances and a higher value for $`\eta `$. In particular we report the results of a new analysis of the degenerate BBN scenario we have performed with our code.
The effect of neutrino chemical potentials on BBN predictions is twofold. Due to the definition of $`N_\nu `$ (3.1), non-vanishing $`\xi _\alpha =\mu _{\nu _\alpha }/T_\nu `$, with $`\alpha `$ denoting the neutrino specie, change its value from the non-degenerate case ($`\xi _\alpha =0`$). In fact for three massless neutrinos with degeneracy parameter $`\xi _\alpha `$, $`N_\nu `$ becomes
$$N_\nu =3+\mathrm{\Sigma }_{\alpha =e,\mu ,\tau }\left[\frac{30}{7}\left(\frac{\xi _\alpha }{\pi }\right)^2+\frac{15}{7}\left(\frac{\xi _\alpha }{\pi }\right)^4\right],$$
(4.11)
implying a larger expansion rate of the universe with respect to the non-degenerate scenario. In this case, nucleons freeze out at a larger temperature, with a higher value for the neutron to proton density ratio, which implies a larger value of $`Y_p`$. This effect does not depend on the particular neutrino chemical potentials but rather on the whole neutrinoโantineutrino asymmetry, via the sum in the r.h.s. of (4.11). In addition, we also note that the neutrino decoupling temperature and the ratio $`T_\nu /T`$, entering in the BBN equations, are affected by a change of $`N_\nu `$ as well. However, it has been checked that this effect is quite negligible on the predictions for the element abundances.
Electron neutrinos entering in the $`np`$ processes, can modify the corresponding rates if their distribution has a non-vanishing $`\xi _e`$. In particular, a positive value for $`\xi _e`$ means a larger number of $`\nu _e`$ with respect to $`\overline{\nu }_e`$ and thus enhances $`np`$ processes with respect to the inverse processes. This, of course, reduces the number of neutrons available at the onset of BBN. Moreover, since the initial condition for BBN at $`T10MeV`$ is fixed by the nuclear thermal equilibrium, which is also kept by reactions like $`\nu _e+np+e^{}`$, the chemical equilibrium fixes $`(\mu _p\mu _n)/T\xi _e`$. This implies an initial $`n/p`$ ratio which is lowered by the factor $`\mathrm{exp}(\xi _e)`$, which again reduces the number of neutrons available at the onset of BBN.
These effects strongly influence the nuclide production, so with no $`\nu _\mu `$ and $`\nu _\tau `$ degeneracy, the value of $`\xi _e`$ is strongly constrained. In this case, the authors of found the limits $`0.06\xi _e0.14`$. For a fully degenerate BBN, at least for $`Y_p`$, the effect of a positive $`\xi _e`$ can be compensated by the contribution to $`N_\nu `$ coming from $`\xi _{\mu ,\tau }`$. In the neutrino degeneracy parameters result to be in the ranges $`0.06\xi _e1.1`$ and $`|\xi _{\mu ,\tau }|6.9`$ for $`\xi _\mu 0`$, $`\xi _\tau =0`$ or viceversa, and $`|\xi _{\mu ,\tau }|5.6`$ for $`\xi _\mu =\xi _\tau 0`$.
We do consider in our analysis as input parameters $`\xi _e`$, $`N_\nu `$ and $`\eta _{10}`$. Likelihood analysis of compatibility between theoretical predictions and experimental data yields contour levels which are surfaces in the three-dimensional space of these parameters. Because of the partial cancellation of the effects due to $`\xi _e`$ and $`\xi _{\mu ,\tau }`$ we have discussed, if we define, in analogy to the nonโdegenerate case, a total likelihood function, $`(\xi _e,N_\nu ,\eta _{10})`$, it is reasonable to expect that it may sensibly differ from zero in a quite wide parameter region, so some bound should be chosen to the possible range of variation for the parameters. We have chosen to constrain $`N_\nu `$ to be smaller than 13. This bound has been obtained, at $`2\sigma `$ level, in , by a likelihood analysis of the BOOMERANG data, as function of $`N_\nu `$, maximizing for each $`N_\nu `$ the likelihood function over all other parameters, including $`\eta `$. As we will see our conclusion on the degenerate BBN scenario versus BOOMERANG and MAXIMA-1 data is completely different than what has been argued in . Nevertheless we think that the upper bound on $`N_\nu `$ is quite robust. We also consider $`\xi _e`$ and $`\eta _{10}`$ in the wide range $`1รท1`$ and $`1รท30`$, respectively.
By using the results of our BBN code for the degenerate neutrino case and performing a study similar to the one presented for the standard scenario, we obtained likelihood functions for both low and high $`D`$ experimental values. We first report the maxima for these functions in the low $`D`$ case
$$\xi _e=0.06,N_\nu =3.43,\eta _{10}=5,$$
(4.12)
and in the high $`D`$ scenario
$$\xi _e=0.35,N_\nu =13,\eta _{10}=4.20,$$
(4.13)
We notice that in the low Deuterium case the solution prefers an almost non degenerate scenario, with values for $`\eta `$ compatible with the one obtained in the previous section and a slightly larger result for $`N_\nu `$.
More interesting is to follows the allowed ranges for $`\xi _e`$ and $`N_\nu `$ as functions of $`\eta `$. In Figures 17 and 18 we show $`(\xi _e,N_\nu ,\eta _{10})`$, evaluated for $`\eta _{10}`$ taking the values in (4.12) and (4.13). Increasing $`\eta _{10}`$ in the considered range, the likelihood functions move towards higher values of $`\xi _e`$ and $`N_\nu `$, along, approximatively, a linear path. This is due to the fact that increasing $`\xi _e`$, which results in a lower value for the $`n/p`$ ratio, and then for $`Y_p`$, must be compensated by a faster expansion produced by a higher value for $`N_\nu `$. This behaviour is highlighted in the $`95\%`$ exclusion plots for the $`\xi _e`$ and $`N_\nu `$ parameters for different values of $`\eta _{10}`$, which are reported in Figures 19 and 20 with the bound $`N_\nu 13`$. These two plots summarize our analysis of the degenerate BBN scenario, providing the combined bounds on $`\xi _e`$ and $`N_\nu `$. For low $`D`$ the allowed range for $`\eta _{10}`$ is $`3.3รท9.9`$, while for high $`D`$ we have $`1.1รท5.8`$. In both cases the degenerate BBN scenario, at $`2\sigma `$ level is compatible with observational results on nuclide abundances even for quite large values for both $`N_\nu `$ and $`\xi _e`$. As expected, however, these two parameters are strongly correlated.
From our analysis we see that a large value for $`\eta 10^9`$ is compatible, at $`2\sigma `$ level, with degenerate BBN for quite large $`N_\nu `$ only, $`N_\nu 10`$, and large positive $`\xi _e0.3`$. The BOOMERANG and MAXIMA-1 results seem to indicate a value for $`\eta 10^9`$, but at the same time suggest $`N_\nu 13`$, again at $`2\sigma `$ level . We may conclude that the degenerate BBN scenario and BOOMERANG and MAXIMA-1 results are compatible only in the very top area in the exclusion plot of Figure (19) for the low $`D`$ scenario. In the high $`D`$ case, the upper limit on $`N_\nu `$ is already reached for $`\eta 6.510^{10}`$. As a conclusion we can say that, if the result $`\eta 10^9`$ will be confirmed, the degenerate BBN scenario is not consistent with CMBR results at $`2\sigma `$ level, while there is agreement for the low $`D`$ case. For slightly lower values for $`\eta (0.5รท0.6)10^{10}`$, which corresponds to the central value measured by MAXIMA-1 experiment, the agreement at $`2\sigma `$ level definitely improves, and it would represent a signal in favour of a large neutrinoโantineutrino asymmetry in the universe, $`0.2\xi _e0.5`$ and $`N_\nu 10`$ . We note that a similar result has been obtained in .
## 5 Conclusions
In this paper we have reported the results of a likelihood analysis of Big Bang Nucleosynthesis theoretical predictions versus the experimental data, both for vanishing neutrino chemical potentials (standard BBN) and in presence of neutrino-antineutrino asymmetries (degenerate BBN). The theoretical estimates have been obtained by using a new updated code we have developed in the recent years to increase the accuracy at the $`1\%`$ level.
In the standard scenario, BBN predictions are in good agreement with experimental data for both low and high $`D`$ results. In the first case, since $`Y_D`$ is a decreasing function of baryonic asymmetry $`\eta `$, the low Deuterium abundance fixes the maximum of total likelihood to larger values of $`\eta `$, and to compensate the effect on $`Y_p`$ yields smaller values of $`N_\nu `$. In this way we have obtained as preferred values $`N_\nu =2.44`$ and $`\eta _{10}=4.69`$. The situation is reversed for high $`D`$ where the maximum lies at $`N_\nu =3.29`$ and $`\eta _{10}=1.81`$. From the plots of $`95\%`$ CL for both the above likelihood functions one gets the corresponding compatibility regions
$`\begin{array}{ccccc}\text{low}D& 1.7N_\nu 3.3& & 4.0\eta _{10}5.7& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0.015}\mathrm{\Omega }_Bh^20.021\\ & & & & \\ \text{high}D& 2.3N_\nu 4.4& & 1.4\eta _{10}3.7& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0.005}\mathrm{\Omega }_Bh^20.014\end{array}.`$ (5.17)
Both ranges on $`N_\nu `$ are comparable with the results of Ref. . For the baryon asymmetry $`\eta `$, the result for low $`D`$ fairly coincides with , whereas the distortion of the contour for $`95\%`$ CL in case of high $`D`$ makes our upper limit on $`\eta _{10}`$ larger than the value of . This effect is mainly due to the non-trivial dependence on the parameters $`N_\nu \eta _{10}`$ of the nuclear theoretical uncertainties.
We have also analyzed the degenerate BBN scenario, motivated by the BOOMERANG and MAXIMA-1 results on the CMBR spectrum. For both high and low Deuterium measurements we have obtained a $`2\sigma `$ exclusion plots for the $`\xi _e`$ and $`N_\nu `$ parameters, and we have shown that the theoretical estimates are in good agreement with the experimental measurements of light nuclei abundances even for large values of these parameters, provided they lie in the regions shown in Fig.s (19) and (20). In particular large values for $`\eta `$ implies high values for the neutrino degeneracy parameters.
A first analysis of the recent CMBR data, , suggests quite large values for $`\eta 10^9`$. From our results we see that this baryon to photon density is actually too high to be in good agreement with degenerate BBN, the compatibility being only at $`2\sigma `$ level, with the low Deuterium scenario. This conclusion holds if we constraint the effective number of neutrinos to be bounded by $`N_\nu 13`$, which, as pointed out in , is again suggested by BOOMERANG data. The agreement improves for slightly smaller values for $`\eta `$, as the central value obtained by MAXIMA-1 results, and would be an evidence in favour of large neutrino degeneracy, $`0.2\xi _e0.5`$, $`N_\nu 10`$, an intriguing feature of the neutrino cosmic background. New experimental results, as the one expected from the MAP and PLANCK experiments on CMBR, as well as new analysis of the BOOMERANG and MAXIMA-1 data, will be crucial in clarifying this issue.
Acknowledgements
The authors are pleased to thank Dr.s Sergio Pastor, Marco Peloso and Francesco Villante for useful discussions and comments. We thank A. Bottino for pointing us the MAXIMA-1 Collaboration results, and M. Kaplinghat for valuable remarks.
## Appendix A Series expansion for integral quantities
The partial derivatives of $`L`$ with respect to $`z`$ and $`\varphi _e`$, $`L_z`$ and $`L_{\varphi _e}`$, can be expressed as the following series<sup>3</sup><sup>3</sup>3We truncate all the series at $`n=7`$ while in the standard code the truncation is at $`n=5`$.
$`L_z(z,\varphi _e)={\displaystyle \frac{dz_R}{dz}}\left\{{\displaystyle \frac{2L}{z_R}}\left({\displaystyle \frac{z_R}{\pi }}\right)^2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^{n+1}\left[K_1(nz_R)+K_3(nz_R)\right]\mathrm{sinh}(n\varphi _e)\right\},`$ (A.1)
$`L_{\varphi _e}(z,\varphi _e)={\displaystyle \frac{2z_R^2}{\pi ^2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^{n+1}K_2(nz_R)\mathrm{cosh}(n\varphi _e),`$ (A.2)
where the $`K_i`$ are the modified Bessel functions of the second kind and in $`z_R(z)=m_e^R/T`$ we consider the renormalized electron mass,
$$m_e^R(T)m_e\left[1+\frac{\pi }{3}\alpha \left(\frac{T}{m_e}\right)^2\right].$$
(A.3)
The series expansion for the remaining electron quantities are
$`\widehat{\rho }_e(z,\varphi _e)`$ $`=`$ $`{\displaystyle \frac{z_R^3}{2\pi ^2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^{n+1}}{n}}\left[3K_3(nz_R)+K_1(nz_R)\right]\mathrm{cosh}(n\varphi _e),`$ (A.4)
$`{\displaystyle \frac{\delta \widehat{\rho }_e}{\delta z}}(z,\varphi _e)`$ $`=`$ $`{\displaystyle \frac{dz_R}{dz}}\{{\displaystyle \frac{3\widehat{\rho }_e}{z_R}}+{\displaystyle \frac{z_R^3}{2\pi ^2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^{n+1}[{\displaystyle \frac{15K_3(nz_R)+K_1(nz_R)}{nz_R}}`$ (A.5)
$`4K_4(nz_R)]\mathrm{cosh}(n\varphi _e)\},`$
$`{\displaystyle \frac{\delta \widehat{\rho }_e}{\delta \varphi _e}}(z,\varphi _e)`$ $`=`$ $`{\displaystyle \frac{z_R^3}{2\pi ^2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^{n+1}\left[3K_3(nz_R)+K_1(nz_R)\right]\mathrm{sinh}(n\varphi _e).`$ (A.6)
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# A new model of a tidally disrupted star
## 1 Introduction
Starting from the seminal paper by Roche the problem of the tidal influence of a gravitating source on a satellite has been addressed by numerous researchers. More recently, interest in this problem has been raised by a paper of Hills (Hills 1975), who proposed tidal disruption processes as the main processes of fueling of QSOโs and AGNโs. From the point of view of the astrophysics of QSOโs and AGNโs there are several approaches to that problem. Firstly one can consider the tidal interaction event as an elementary process in the complicated astrophysical environment of a supermassive black hole, presumably situated in the cores of QSOโs and AGNโs. Then one could find the main evolutionary characteristics of such a system and its average luminosity, taking into account additional gas dynamical and stellar dynamical processes occurring in the cores (e.g. Hills, 1975, Frank $`\&`$ Rees 1976, Young et al 1977; Young 1977; Hills 1978; Frank 1979; Gurzadian $`\&`$ Ozernoi 1980; Lacy et al 1982; Illarionov $`\&`$ Romanova 1986a; Illarionov $`\&`$ Romanova 1986b; Dokuchaev 1991; Beloborodov et al 1992; Roos 1992; Syer $`\&`$ Ulmer 1999; Magorrian $`\&`$ Tremaine 1999). One can also consider the evolution of remnants of a single tidal stripping or tidal disruption event and find the characteristic luminosity change of the object due to the accretion of the remnants through an accretion disk or quasi-spherical configuration onto the central black hole (e.g. Lacy et all 1982; Rees 1988; Evans $`\&`$ Kochanek 1989; Cannizzo et all 1990; Roos 1992; Kochanek 1994; Ulmer et al 1998; Kim et al 1999; Syer $`\&`$ Ulmer 1999; Ulmer 1999).
On the other hand it is very important to understand quantitatively the main characteristics of the tidal encounter itself, and a lot of of work has been devoted to the physical processes occurring in a star during its fly-by around a black hole. The papers on that subject could be classified by the different stellar models used in the calculations. The simplest possible approach to the problem uses an incompressible model of the star. Thus one can reduce the complicated hydrodynamical nonlinear partial differential equation governing the evolution of the stellar gas to a set of ordinary differential equations, which are easy to analyze by analytical and numerical means. The study of incompressible models has been performed for Newtonian and relativistic tidal fields and different kinds of orbits of the star (e.g. Nduka, 1971; Fishbone, 1973; Mashhoon, 1975; Luminet $`\&`$ Carter 1986; Kosovichev $`\&`$ Novikov, 1992). However, this approach is highly unrealistic, since effects determined by the compressibility of the star can play a major role during the tidal disruption event (e.g. Carter $`\&`$ Luminet 1982).
A significant step forward was made by Lattimer and Schramm (Lattimer $`\&`$ Schramm 1976) and by Carter and Luminet (Carter $`\&`$ Luminet 1983, 1985) who proposed the so-called affine model of the tidally disrupted star, which allows for the compressibility of the stellar gas. In this model the law of time evolution of different elements of the star is defined in terms of some spatially uniform $`3\times 3`$ matrix $`๐(t)`$:
$$x^i=Q_j^i(t)x_0^j,$$
where $`x^i`$ are the components of the position vector of a gas element, $`x_0^j`$ are the components of the position vector in some reference state (say, before the tidal field โis switched onโ), and summation over repeated indices is assumed. Then one can find the evolution equations for the matrix elements from the so-called virial relations written for the whole star. The affine model has successfully been applied to the problem of tidal interaction and tidal disruption of a star by a supermassive black hole during close encounters (e.g. Carter $`\&`$ Luminet 1983, 1985; Luminet $`\&`$ Mark 1985; Luminet $`\&`$ Carter 1986; Luminet $`\&`$ Pichon 1986; Novikov et al 1992; Diener et al 1995). Lai, Rasio and Shapiro used the same model for an approximate treatment of an isolated rotating star, as well as for a star in a binary system (e.g. Lai et al 1994; Lai $`\&`$ Shapiro 1995).
Recent progress in numerical simulations has allowed researchers to perform direct 3D simulations of the tidal interaction and tidal disruption events. The first SPH simulations were run in the beginning of eighties by Nolthenius and Katz, and by Bicknell and Gingold, although the number of particles in these simulations was too small to be representative (Nolthenius $`\&`$ Katz 1982, 1983; Bicknell $`\&`$ Gingold 1983). In the following decades, SPH simulations were improved both by increasing of the number of particles, and by using of more complicated stellar models (e.g. Evans $`\&`$ Kochanek 1989; Laguna et al 1993; Laguna 1994; Fulbright et al 1995; Ayal et al 2000). Three dimensional finite difference simulations were done by Khokhlov, Novikov and Pethick for a polytropic star in a Newtonian tidal field (Khokhlov et al 1993a,b, hereafter Kh a,b), by Frolov, Khokhlov, Novikov and Pethick for a white dwarf (Frolov et al 1994), and by Diener, Frolov, Khokhlov, Novikov and Pethick for a polytropic star in the tidal field of a Kerr black hole (Diener et al 1997). An interesting attempt to combine the affine model and a simple version of 3D finite difference hydrodynamics has been made by Mark, Lioure and Bonazzola (Mark et al 1996).
Although the 3D simulations promise the most direct and thoughtful approach to the problem, they are still very time consuming. All in all, less than one hundred different sets of values of the problem parameters have been tested with numerical experiments, and due to very poor statistics these experiments cannot be used to characterize the general properties of the tidal encounters for a broad range of available parameters. There is another, more fundamental difficulty connected with the 3D simulations. The complexity of 3D hydrodynamical flows makes the interpretation of the results of numerical work increasingly difficult. The situation is reminiscent of a real physical experiment, and a simple โreferenceโ model of the tidally disrupted star would be very welcome in order to interpret the results of the numerical simulations. On the other hand, the astrophysics of AGNโs and QSOโs requires a rather rough description of a single tidal encounter, and only a few โaveragedโ quantities such as e.g. the amount of mass lost by the star during the tidal interaction, or the amount of energy deposited in the star by the tidal forces are of interest from the astrophysical viewpoint.
In this paper we propose a new, semi-analytical model of the tidally interacting or tidally disrupted star which could be used for intensive calculations covering the whole parameter space of the problem, and also as a โreferenceโ model for 3D simulations. Our model is a straightforward generalization of the affine model. However, in contrast to the affine model, the different layers of the star evolve differently in our model, and are connected to each other by a force determined by pressure. This allows us to employ our model for calculation of quantities such as the loss of mass from the star after a fly-by over a black hole without complete disruption, which cannot be calculated in the affine approximation. Instead of the position matrix $`๐(t)`$ of the affine model, we use the position matrix $`๐(t,r_0)`$, which depends not only on time, but in addition on the value $`r_0`$ of the โreferenceโ vector $`x_0^i`$ (obviously the radius $`r_0`$ plays the role of a Lagrangian coordinate, so we will later call it the Lagrangian radius). Thus, in our model the star consists of elliptical shells which are composed of all elements of the star with a given Lagrangian radius $`r_0`$. The evolution of the shell depends on the Lagrangian radius, and therefore the shells have different ratios between their major axes and different rotation angles with respect to a (locally inertial) coordinate frame centered on the starโs center of mass, for the different values of the Lagrangian radius. The evolution equations of our model follow from the virial relations written for each shell (see e.g. Chandrasekhar 1969, hereafter Ch). Unlike the affine model the virial relations written for a shell inside the star must contain surface terms, and these surface terms lead to interactions between shells with different Lagrangian radii, and therefore to the propagation of a disturbance through the star. In fact, the evolution equations are of hyperbolic type, and the disturbance induced by a tidal field propagates over the star as a non-linear sound wave. We derive the evolution equations of our model in the next Section using certain approximations for the pressure terms, and the terms describing the self-gravity of the star. In the simplest formulation of our model the interaction between the shells depends only on their relative volumes, and therefore the shells are allowed to intersect each other. Therefore the position matrix $`๐`$ has no direct physical meaning in such a case, and one should only use quantities averaged over many shells in order to infer physical information (such as e.g. the energy and the angular momentum contained inside some part of the star, the components of the quadrupole moment tensor for that part, the amount of mass lost by the star, and so on). We show that the energy and the angular momentum are well defined in our model, and derive the law of evolution of these quantities due to the presence of a tidal field. We also show that the circulation of velocity of a gas element over the shell is exactly conserved even in the presence of the tidal field. Then we apply our model to the simplest problem, the parabolic fly-by of a polytropic star around a source of Newtonian gravity, and numerically calculate the evolution of the quantities characterizing the star during the fly-by. We compare our results with 3D finite difference simulations of Kh a,b for the same problem and the same parameters, and find very good agreement. Additionally we compare our model with results of SPH simulations, calculations based on the affine model and results from the linear theory of tidal perturbations (Press $`\&`$ Teukolsky 1977; Lee $`\&`$ Ostriker 1986). Then we calculate the energy deposited in the star, its angular momentum and the amount of mass lost by the star as a function of the pericentric separation between the star and the center of gravity. As it will be clear for the results in Section 3, our model gives a better agreement with the results of the 3D simulations than the affine model.
We use a rather unusual summation convention assuming that summation is performed over all indices appearing in our expressions more than once, but summation is not performed if indices are enclosed in brackets. Bold letters represent matrices in abstract form. All indices can be raised or lowered with help of the Kronecker delta symbol, but nevertheless we distinguish between the upper and lower indices in order to enumerate the rows and columns of matrices, respectively. Therefore, the expression $`A_i^kB_i^l=C^{kl}`$ means $`\mathrm{๐๐}^T=๐`$, and $`A_k^iB_l^i=C_{kl}`$ means $`๐^T๐=๐`$ (here $`T`$ stands for the transpose of a matrix).
Finally, we would like to list the main approximations made in the derivation of the dynamical equations of our model.
1) We assume that the star is composed of elliptical shells, and the shells are not deformed during the evolution of the star in a tidal field.
2) We calculate the self-gravity of the star in a simplified manner. Namely, in order to calculate the force of gravity acting on some particular shell, we neglect the contribution of the starโs mass concentrated in the outer (with respect to that shell) layers of the star. It is also assumed that the gravitational force determined by the inner layers is equivalent to the gravitational force of a uniform density ellipsoid inserted in this shell.
3) We assume a polytropic equation of state of the stellar gas.
4) We use an โaveragedโ density and an โaveragedโ pressure instead of the exact quantities. These averaged quantities depend only on time and the Lagrangian radius of the star.
The approximations 1,2 are essential for our model. The approximations 3,4 can be relaxed in a more advanced variant of the model.
## 2 Building up of the model
As was mentioned in the Introduction, we divide the star into a set of elliptical shells. Each shell consists of the gas elements which had the same distance from the center of the star in the unperturbed spherical state. The initial Cartesian coordinates of the gas elements of the star $`x_0^i`$ play the role of Lagrangian coordinates in the course of the starโs evolution under the influence of a tidal field, and hereafter we simply call them the Lagrangian coordinates. We assume that the star layers corresponding to the same Lagrangian radius $`r_0=\sqrt{x_{0i}x_0^i}`$ always keep the elliptical form. The parameters of the shells are different for the different values of $`r_0`$, and evolve with time according to some dynamical equations, which are derived below. Let us consider the Eulerian coordinates of the gas elements $`x^i`$ with respect to some inertial reference system centered at the starโs geometrical center. From our discussion it follows that the law of transformation between the Lagrangian and the Eulerian coordinates can be written as
$$x^i=T_j^i(t,r_0)e_0^j,$$
$`1`$
where $`e_0^i=x_0^i/r_0`$, and $`e_{0i}e_0^i=1`$. We also introduce the matrix $`๐`$ which is the inverse of the matrix $`๐`$.
The position matrix $`๐`$ and its inverse $`๐`$ can be represented as a product of two rotational matrices $`๐`$ and $`๐`$, and a diagonal matrix $`๐`$:
$$T_j^i=A_l^iB_m^lE_j^m=a_lA_l^iE_j^l,S_j^i=a_l^1A_l^jE_i^l,$$
$`2`$
where $`B_m^l=a_{(l)}\delta _m^{(l)}`$, and $`a_l`$ are the principal axes of the elliptical shell. Clearly, the matrices $`๐`$ and $`๐`$ describe the rotation of the principal axes of the shell with respect to the reference frames in the Eulerian and Lagrangian spaces, respectively. For our purpose it is useful to define several quantities connected with the matrices $`๐`$ and $`๐`$, namely the determinant $`g`$ of the position matrix:
$$g=|๐|=a_1a_2a_3,$$
$`3`$
the Jacobian $`D=|\frac{x^i}{x_0^j}|`$ of the mapping between the Lagrangian and Eulerian spaces:
$$D(x_0^i)=\frac{ge_{0l}e_{n0}R^{ln}}{r_0^2},$$
$`4`$
where the symmetric matrix $`R^{ln}`$ determines the local shear and change of volume of the neighboring shells:
$$R^{ln}=\frac{1}{2}(S_m^l(T_n^m)+S_m^n(T_l^m)),$$
and the prime stands for the differentiation with respect to $`r_0`$. We also use the โaveragedโ Jacobian
$$\overline{D}=\frac{1}{4\pi }๐\mathrm{\Omega }D=\frac{dg}{dr_0^3},$$
$`5`$
where the integration is performed over a unit sphere in the Lagrangian space and $`d\mathrm{\Omega }`$ is the elementary solid angle. Eq. (4) immediately gives the law of evolution of the gas density $`\rho `$ in our model. Taking into account the law of mass conservation we have
$$\rho (t,x^i)=\rho _0(r_0)/D,$$
$`6`$
where $`\rho _0(r_0)`$ is the density distribution in the unperturbed star. The โaveragedโ density $`\overline{\rho }`$ is defined analogously to (6), but with the โaveragedโ Jacobian (5)
$$\overline{\rho }(t,r_0)=\rho _0/\overline{D}=\frac{3}{4\pi }\frac{dM}{dg},$$
$`7`$
where the mass differential $`dM=4\pi \rho _0r_0^2dr_0`$.
Of course, the transformation law (1) is incompatible with the exact hydrodynamical equations of motion of a perturbed star, and we must introduce some reasonable approximations which allow us to reduce the equations of motion to a dynamical equation for the matrix components $`T_i^j`$. As a starting point of our analysis we use the integral consequences of the exact equations of motion, namely the equation of energy conservation and the so-called virial relations (see e. g. Ch, p. 20). In the adiabatic approximation the energy equation has the form
$$\frac{d}{dt}\{d^3x(\rho v^2/2+ฯต)+๐ซ\}=๐S_i(pv^i)+d^3x(\rho C_{ij}v^ix^j),$$
$`8`$
and the virial relations are
$$\frac{d}{dt}d^3x(\rho x^kv^i)=d^3x(\rho v^kv^i)+\delta ^{ki}d^3xp$$
$$๐S_i(x^kp)+๐ซ^{ki}+d^3x(\rho C_j^ix^kx^j).$$
$`9`$
Here $`v^i`$ is the velocity of the gas element, $`p`$ is the pressure and $`ฯต`$ is the energy density per unit volume. The matrix $`C_j^i`$ represents the tidal tensor, and therefore it is symmetric and traceless. The potential energy $`๐ซ`$ and the potential-energy tensor $`๐ซ^{ki}`$ are
$$๐ซ=\frac{1}{2}d^3xd^3x_1\rho (x^i)\rho (x_1^i)\frac{1}{|\stackrel{}{x}\stackrel{}{x}_1|},$$
$`10`$
and
$$๐ซ^{ki}=\frac{1}{2}d^3xd^3x_1\rho (x^i)\rho (x_1^i)\frac{(x^kx_1^k)(x^ix_1^i)}{|\stackrel{}{x}\stackrel{}{x}_1|}.$$
$`11`$
Obviously, $`๐ซ^{ki}`$ is a symmetric matrix, and the relation
$$๐ซ=๐ซ^{kk}$$
$`12`$
holds. The volume integration in eqs. (8,9) and (10,11) is performed over the volume surrounded by a surface $`r_0=`$const, and the surface integration is performed over that surface <sup>1</sup><sup>1</sup>1Strictly speaking we must extend the integration in the inner integral in eqs. (10,11) over the whole volume of the star. However in our approximation the outer part of the star does not influence gravitationally the inner part of the star, and that part of the integrals can be omitted..
One can try to calculate the integrals containing the pressure and the energy density (the thermal terms) in eqs (8,9), and the potential energy tensor and the potential energy (the gravitational terms) directly, using the transformation law (1), the density distribution (6), and the adiabatic condition. However this approach leads to rather complicated expressions for the thermal terms. We want to construct our model in the simplest possible way, and therefore we make several additional approximations with respect to these terms. For the thermal terms we use an โaveragedโ pressure $`\overline{p}(t,r_0)`$, and an โaveragedโ energy density $`\overline{ฯต}(t,r_0)`$ instead of the exact quantities $`p(t,x^i)`$ and $`ฯต(t,x^i)`$. That approximation allows us to represent the thermal terms in eqs (8,9) in a very simple form
$$\delta ^{ki}d^3xp๐S_i(x^kp)=\frac{4\pi }{3}\delta ^{ki}_0^{\overline{p}(r_0)}g๐\overline{p},$$
$`13`$
$$๐S_i(pv^i)=4\pi H(t,r_0)\overline{p}(t,r_0),$$
$`14`$
where the expansion rate is
$$H=\frac{1}{3}S_i^l\dot{T}_l^i=\frac{1}{3}(\frac{\dot{a}_1}{a_1}+\frac{\dot{a}_2}{a_2}+\frac{\dot{a}_3}{a_3}),$$
$`15`$
and the dot stands for time differentiation. In our approximation the pressure force acting on the shell from the side of the neighboring shells depends on the relative values of the shell volumes, and does not depend on the orientation of the shells with respect to each other. Therefore, the shells can intersect each other, and the interpretation of the position matrix $`๐`$ as describing the Eulerian positions of the gas elements is rather ambiguous. As we have already mentioned in the Introduction quantities are only meaningful in this approximation when averaged over many shells. In order to obtain the physical interpretation of the position matrix $`๐`$ one should use a more complicated model, where the pressure force depends on the relative orientation of the shells (see also Discussion).
In order to calculate an โaveragedโ potential energy tensor $`\overline{๐ซ}^{ik}`$ we assume that the gravitational force acting on the gas near the shell with some Lagrangian radius $`r_0`$ is equivalent to the gravitational force of a uniform density ellipsoid with a mass equal to the part of the starโs mass within that shell. The principal axes of that ellipsoid coincide with the principal axes of the shell, and the density is averaged over the volume enclosed in the shell. Under this assumption the โaveragedโ potential energy tensor $`\overline{๐ซ}^{ik}`$ has the form
$$\overline{๐ซ}^{ik}=\frac{1}{2}๐MGMA_j^iA_j^k\frac{a_j^2D_j}{g},$$
$`16`$
and the averaged potential energy $`\overline{๐ซ}=\overline{๐ซ}^{kk}`$ is
$$\overline{๐ซ}=\frac{1}{2}๐MGM\frac{a_j^2D_j}{g}.$$
$`17`$
Here we use the mass $`M=4\pi _0^{r_0}\rho _0(r_1)๐r_1`$ of the gas inside the shell of the radius $`r_0`$ as a new Lagrangian coordinate instead of $`r_0`$. The dimensionless quantities $`D_j`$ have been described by e. g. Ch, p. 41. They have the form:
$$D_j=g_0^{\mathrm{}}\frac{du}{\mathrm{\Delta }(a_j^2+u)},$$
$`18`$
where
$$\mathrm{\Delta }=\sqrt{(a_1^2+u)(a_2^2+u)(a_3^2+u)}.$$
These quantities obey very useful relations
$$\frac{a_j^2D_j}{g}=_0^{\mathrm{}}\frac{du}{\mathrm{\Delta }},$$
$`19`$
and
$$\underset{i=1}{\overset{3}{}}D_i=2.$$
Now we can substitute eqs (13,14) and (16,17) into eqs (8,9), and perform the integration in the other terms with the help of the law of mass conservation: $`d^3x\rho =d^3x_0\rho _0(r_0)`$. From the energy equation we obtain
$$\frac{d}{dt}\{๐M\frac{\dot{T}_n^i\dot{T}_n^i}{2}+4\pi \overline{ฯต}๐g+3\overline{๐ซ}\}=12\pi H(r_0)\overline{p}(r_0)g(r_0)+๐MC_j^i\dot{T}_l^iT_l^j,$$
$`20`$
and from the virial relations we obtain
$$๐MT_l^k\ddot{T}_l^i=4\pi \delta ^{ik}g๐\overline{p}+3๐ซ^{ik}+๐MC_j^iT_l^jT_l^k.$$
$`21`$
Differentiating these equations over the mass coordinate we have
$$\frac{d}{dt}\{\frac{\dot{T}_n^i\dot{T}_n^i}{2}+4\pi \overline{ฯต}\frac{dg}{dM}\frac{3}{2}a_j^2D_j\frac{GM}{g}\}=12\pi \frac{d(H\overline{p}g)}{dM}+C_j^i\dot{T}_l^iT_l^j,$$
$`22`$
and
$$\ddot{T}_n^i=4\pi S_i^ng\frac{d\overline{p}}{dM}\frac{3}{2}A_j^ia_jD_jE_n^j\frac{GM}{g}+C_j^iT_n^j,$$
$`23`$
The dynamical equations (23) are the main result of this Section. Note that if the tidal term is absent and the position matrix is proportional to $`\delta _n^i`$: $`T_n^i=a\delta _n^i`$ the equations (23) are reduced to a single equation, which describes the radial adiabatic oscillations of a star. In that case this equation follows from the hydrodynamical equations without any approximation. For the position matrix of a special type $`T_n^i=\widehat{T}_n^i(t)r_0`$ our equations are reduced to the dynamical equations of the affine model (Appendix A). Also, in the case of an incompressible fluid $`\rho (t,x^i)=\rho _0`$=const the equations (23) are exact (see Appendix A).
Eq. (22) must follow from eq. (23). To prove it we contract both sides of (23) with the velocity matrix $`\dot{T}_n^i`$ over all indices, and subtract the result from eq. (22). The remainder can be separated into thermal and gravitational parts. The thermal part is easily reduced to the relation:
$$d(\frac{\overline{ฯต}}{\overline{\rho }})+\overline{p}d(\frac{1}{\overline{\rho }})=0,$$
$`24`$
which obviously reflects the application of the first law of thermodynamics to our case. Neglecting the possible presence of shocks, for an ideal gas with polytropic index $`\gamma `$ this equation is integrated to give
$$\overline{p}=C(r_0)\overline{\rho }^\gamma ,$$
$`25`$
where the entropy constant $`C(r_0)`$ is determined from an unperturbed model of the star. The gravitational part is reduced to the equality
$$\frac{d}{dt}(\frac{a_j^2D_j}{g})+\frac{D_ja_j\dot{a}_j}{g}=0.$$
$`26`$
Differentiating eq. (19) and using the definition of $`D_j`$, one can see that in fact, the equality (26) is an identity.
The law of evolution of angular momentum can be easily obtained from (23). Contracting both parts of (23) with $`\dot{T}_n^k`$ and taking the antisymmetric part of the result, we have:
$$\frac{d}{dt}(T_n^k\dot{T}_n^iT_n^i\dot{T}_n^k)=C_j^iT_n^jT_n^kC_j^kT_n^jT_n^i=a_n^2A_n^j(C_j^iA_n^kC_j^kA_n^i),$$
$`27`$
where we use eq. (2) to obtain the last equality. Obviously, the eq. (27) describes the rate of change of angular momentum due to a tidal torque.
Similar to the dynamical system describing the motion of an incompressible ellipsoid in a tidal field (e. g. Ch, p. 74) and the affine model (Carter and Luminet, 1985), our system has three additional quantities which are exactly conserved if the system is evolving in a tidal field <sup>2</sup><sup>2</sup>2Unlike these models, our quantities are not numbers, but functions of the Lagrangian variable $`M`$.. Contracting (23) with $`\dot{T}_m^i`$, taking the antisymmetric part of the result, and using the symmetric properties of the tidal tensor, we see that
$$\chi _{mn}(M)=T_n^i\dot{T}_m^iT_m^i\dot{T}_n^i,$$
$`28`$
($`\chi _{mn}=\chi _{nm}`$) do not depend on time. It is easy to find the physical meaning of $`\chi _{mn}`$. For that let us introduce the dual vector $`\chi _l=\frac{1}{2}ฯต_{lmn}\chi _{mn}`$, and consider the circulation $`C=_Lv^i๐x^i`$, where the integration is performed over a closed path $`L`$ on the surface of an ellipsoid of a given Lagrangian radius. Along this path the vector $`e_0^i`$ describes a closed curve on a sphere of unit radius. Let us consider a surface on that sphere enclosed in that curve, and denote projections of that surface on the coordinate planes of a Cartesian coordinate system where the components $`e_0^i`$ are defined, as $`S^i`$. Then it is easy to see that the circulation can be expressed as
$$C=\chi _iS^i,$$
$`29`$
and therefore the conservation of the quantities $`\chi _{mn}`$ corresponds to the conservation of the circulation of the fluid over our elliptic shells. It is also interesting to note that the angular momentum tensor and the quantities $`\chi _{mn}`$ are adjoint in a certain sense provided the tidal interaction is switched off. For that it is sufficient to note that the motion determined by the transpose $`๐^T`$ of $`๐`$ also provides a solution of the system (23) (for a motion of an incompressible ellipsoid the similar statement is known as Dedekindโs theorem). Clearly, the quantities $`\chi _{mn}`$ play the role of angular momentum for the motion determined by $`๐^T`$, and the angular momentum tensor plays the role of $`\chi _{mn}`$. The presence of tidal interactions breaks this symmetry. Since the star is usually assumed to be non-rotating before the tidal field is โswitched onโ, we set $`\chi _{mn}=0`$.
At the end of this Section let us note that similar to the affine model, the dynamical equations of our model can be deduced from a variational principle. Consider a Lagrangian of the form:
$$=๐M\{\frac{\dot{T}_n^i\dot{T}_n^i}{2}+\frac{3}{2}a_j^2D_j\frac{GM}{g}+\frac{C_k^iT_l^iT_l^k}{2}\}4\pi \overline{ฯต}๐g,$$
$`30`$
The first variance of that Lagrangian can be written as <sup>3</sup><sup>3</sup>3 The variance of the thermal part of (30) is transformed with help of (5),(24): $`\delta \overline{ฯต}๐g=\delta g๐\overline{p}`$.
$$\delta =๐M\{\dot{T}_n^i\delta \dot{T}_n^i\frac{3}{2}A_j^ia_jD_jE_n^j\frac{GM}{g}\delta T_n^i+C_k^iT_n^k\delta T_n^i4\pi g\frac{d\overline{p}}{dM}S_i^n\delta T_n^i\}$$
$`31`$
Substituting the variance (31) into the Lagrange equations
$$\frac{d}{dt}\frac{\delta }{\delta \dot{T}_n^i}=\frac{\delta }{\delta T_n^i},$$
$`32`$
we arrive again at eq. (23).
## 3 Results of numerical calculations
As we mentioned in the Introduction, for our numerical work we choose a simple problem, the tidal encounter of a polytropic star moving around a source of Newtonian gravity (referred to as a black hole) on a parabolic orbit. We assume that the star consists of an ideal gas with constant specific heat ratio $`\gamma =5/3`$. In this case our problem can be described by two parameters: a) the polytropic index $`n`$, and b) the parameter $`\eta `$ reflecting the strength of a tidal encounter:
$$\eta =\sqrt{\frac{M_{}}{M_h}\frac{R_p^3}{R_{}^3}}=(\frac{R_p}{R_T})^{\frac{3}{2}},$$
$`33`$
where $`M_h`$ is mass of the black hole, $`M_{}`$ and $`R_{}`$ are the mass and radius of the star, respectively, and $`R_p`$ is the value of the pericentric separation distance between the star and the black hole <sup>4</sup><sup>4</sup>4From a physical point of view, $`\eta `$ is approximately the ratio of the time which star spends near the pericentric distance to the characteristic stellar time $`t_{}=\sqrt{\frac{R_{}^3}{GM_{}}}`$. Therefore the strong tidal encounters ($`\eta 0`$) are short..
$$R_T=\sqrt[3]{\frac{M_h}{M_{}}}R_{}0.91\sqrt[3]{\frac{M_h}{\pi \rho _{}}},$$
$`34`$
is a characteristic โtidalโ radius, and the average density $`\rho _{}=\frac{3M_{}}{4\pi R_{}^3}`$. In a very approximate sense it can be said that stars moving on orbits with $`R_p`$ smaller than $`R_T`$ (and $`\eta <1`$) experience a strong tidal influence and could be disrupted, and stars moving on orbits with $`R_p>R_T`$ ($`\eta >1`$) are rather weakly perturbed by the tidal field. In fact, as we see later this conclusion depends on the polytropic index of the star (see also Kh b) <sup>5</sup><sup>5</sup>5Let us remind that in the classic Roche problem, the infinitesimal incompressible satellite of density $`\rho `$ rotating about an object of mass $`M`$ in a circular Keplerian orbit, loses its stability if the radius of the orbit is smaller than $`R_{Rh}2.23\sqrt[3]{\frac{M}{\pi \rho }}`$, e.g. Ch, p. 12..
We calculate the main characteristics of the tidal disruption event with a simple explicit conservative Lagrangian numerical scheme (see Appendix B for the details), and compare the results of our calculations with the results of 3D simulations reported by Kh a,b. The values of the polytropic index $`n`$ are $`n=1.5`$, $`2`$, $`3`$, and the parameter $`\eta `$ is changed over a rather broad range. The Eulerian coordinate system coincides with the Lagrangian coordinate system in the beginning of calculations. The plane $`(XOY)`$ of the Cartesian frame corresponding to the Eulerian coordinate system coincides with the orbital plane, and the axis $`OX`$ is directed opposite to the black hole when the star passes the point of minimal separation. All results of the calculations are expressed in terms of natural units. Our spatial unit is the starโs radius $`R_{}`$, and the time unit is $`t_{}=\sqrt{\frac{R_{}^3}{GM_{}}}`$. We use dimensionless time $`\tau =t/t_{}`$, and the moment of passage of the minimal separation distance by the star corresponds to $`\tau =0`$. The energy and the angular momentum gained by the star are expressed in units of $`GM_{}^2/R_{}`$ and $`M_{}\sqrt{GM_{}R_{}}`$ respectively. As a Lagrangian coordinate we use the ratio $`x`$ of mass within a particular shell to the total mass of the star: $`x=M(r_0)/M_{}`$.
The results of our calculations are presented in Figures 1-17.
In Figure 1 (a-d) we show projections of our elliptical shells on the plane $`(XOY)`$ for four different moments of time $`\tau =0`$; $`1`$; $`2`$; $`3`$. The shells shown correspond to four different Lagrangian coordinates $`x=0.2`$; $`0.4`$; $`0.6`$; $`0.8`$. Although the positions of the shells have no direct physical meaning since the shells do not describe directly the mass distribution over the star and can intersect each other, these plots give useful qualitative information about the evolution of our model. The model parameters for these plots $`n=1.5`$ and $`\eta =1.5`$ correspond to a tidal encounter of moderate strength. Let us recall that one of the principal axes of the tidal tensor is always oriented toward the black hole, and two others are perpendicular to this direction. Therefore, as we see from Fig.1 (a-d), the principal axes of the shells do not coincide with the principal axes of the tidal tensor, and it can be said that our shells lag behind the changing tidal field. A similar effect has been observed in 3D numerical simulations of the tidal encounter. At the moment $`\tau =0`$ all shells are almost aligned with respect to each other, but during the subsequent evolution this alignment disappears. The outer layers of the star lag more than the inner layers, and at the time $`\tau =3`$ intersections between the shells are observed. Approximately at the same time in the 3D finite difference and SPH models the star takes a typical $`S`$-shaped form.
Figures 2-5 show the evolution of the central density $`\rho _c`$ (expressed in units of the central density $`\rho _{c0}`$ of the unperturbed star) with time for the different parameters of the model. The dashed lines in Figures 2-4 correspond to the results taken from the 3D finite difference simulations of Kh a,b, and the dashed line in Figure 5 corresponds to the result of SPH simulations by Fulbright et al 1995. The model parameters for these Figures are: $`n=1.5`$, $`\eta =3`$ for Fig. 2a; $`n=2`$, $`\eta =3`$ for Fig. 2b; $`n=2`$, $`\eta =2.5`$ for Fig. 2c; $`n=2`$, $`\eta =2`$ for Fig. 2d; $`n=3`$, $`\eta =0.5`$ for Fig. 3; $`n=2`$, $`\eta =0.1`$ for Fig. 4; $`n=1.5`$, $`\eta =5^{3/2}0.0894`$ for Fig. 5. Figs 2(a-d) correspond to rather weak tidal forces; Fig. 3 presents the case of moderately strong tidal disruption and Figs 4,5 present the case of a very strong tidal disruption event. For the cases of weak tidal influence we observe oscillations of the star after the moment $`\tau =0`$. The period of these oscillations is close to the period $`T_f`$ of the fundamental mode of radial stellar pulsations: $`T_f=3.81`$ for $`n=1.5`$, and $`T_f=\pi `$ for $`n=2`$ (e.g. Cox 1980). It seems that in our model the amplitude of these oscillations is always smaller than in the 3D simulations. We checked that this decrease of the amplitude does not depend on the value of artificial viscosity used in the computations. Fig. 3 gives an example of a moderately strong tidal disruption event (only about 10 percent of the mass remains gravitationally bounded after this encounter, see Kh b and Figure 17). After $`\tau =0`$ the central density monotonically decreases with time toward a very small number (we obtained $`\rho _c/\rho _{c0}1.410^2`$ at the end of our calculations $`\tau _{end}=15`$). In general our curve is close to the curve obtained in the 3D simulations, but there is a significant deviation near the time $`\tau 1`$. Unfortunately, the curve obtained in the 3D case is not continued after $`\tau =1`$, and we cannot say whether our asymptotic value for $`\rho _c/\rho _{c0}`$ is close to the 3D result or not. During the tidal disruption events corresponding to the results presented in Figs 4,5 a very strong tidal disruption of the star occurs, and a typical โspikeโ in the curves for the evolution of the central density is observed. The presence of this โspikeโ has been predicted analytically by Carter $`\&`$ Luminet $`1982`$. The amplitude of this โspikeโ is 0.64 times smaller than was obtained by Kh b for the $`n=2`$, $`\eta =0.1`$ case, and about $`1.6`$ times larger than was obtained by Fulbright et al 1995 in their SPH simulations of $`n=1.5`$, $`\eta =0.0894`$ case. The asymptotic values of $`\rho _c/\rho _{c0}`$ at the limit of large time almost coincide between our curves and the curves obtained by other methods.
In Figures 6-9 the dependence of total, potential, thermal and kinetic energies on time is shown. In all cases we plot the difference between the energy of a particular kind and its corresponding equilibrium value; $`E_{tot}`$, $`E_g`$, $`E_{th}`$ and $`E_K`$ stand for the differences of total, gravitational, thermal and kinetic energies, respectively. We found that too much energy is stored in the outer layers of the star in our model, and calculated the energy differences using $`97.5`$ percent of the starโs mass <sup>6</sup><sup>6</sup>6Typically we obtained $`1020`$ percent of the total energy stored in the outer layers of the star at the end of the computation.. The cases shown are: $`n=1.5`$, $`\eta =3`$ (Fig. 6); $`n=3`$, $`\eta =1.5`$ (Fig. 7); $`n=3`$, $`\eta =1`$ (Fig. 8) and $`n=3`$, $`\eta =0.5`$ (Fig. 9). Fig. 6 and Fig. 7 correspond to the case of a small tidal influence on the star, and Figure 8 presents the case of the tidal encounter of a moderate strength. As in the density plots we see the oscillations of the potential and thermal energies, and these oscillations are more pronounced in the 3D simulations (dashed lines). The oscillations of the potential and the thermal energies seem to compensate each other, and there is no oscillation of the total and kinetic energies. The relative difference between the asymptotic (at large time) values of $`E_{tot}`$ in our calculation and in the 3D calculations is about $`10`$ percent for the $`n=1.5`$, $`\eta =3`$ case, $`25`$ percent for the $`n=3`$, $`\eta =1.5`$ case. The same difference looks small for the case $`n=3`$, $`\eta =1`$ (about $`7`$ percent), but the 3D simulations end at the time $`\tau 3.5`$ which is too short to make conclusions about the asymptotic value of the total energy. Figure 9 presents the case of a strong tidal encounter. The relative difference of the total energy $`E_{tot}`$ is about $`20`$ percent for that case, but again the 3D simulations end at time $`\tau 1`$ and the asymptotic values of $`E_{tot}`$ cannot be compared. We see that $`E_{tot}`$ and $`E_K`$ continue to grow all the time in that case, but this growth takes place only for the gas stripped from the star. Let us define the gravitationally bound debris of the star as a combination of all stars elements where the sum of kinetic and potential energies is less than zero. The dotted line in Figure 9 presents the difference between the total energy of the gravitationally bounded debris and the total energy of the unperturbed star. This difference should tend to an asymptotic value at the large time limit, which is equal to $`0.718`$ for that case.
Figures 10-12 show the value of angular momentum in the direction perpendicular to the starโs orbit, gained by the star during the tidal encounters. The shown cases are: $`n=1.5`$, $`\eta =3`$ (Fig. 10); $`n=3`$, $`\eta =1.5`$ (Fig. 11); $`n=3`$, $`\eta =0.5`$ (Fig 12). For the cases of weak tidal influence the relative deviation between our model and the 3D results is about $`1020`$ percent. For the case of strong tidal disruption the same difference is approximately $`10`$ percent. In that case the total angular momentum grows with time, but the angular momentum of the gravitationally bound debris tends to a small asymptotic value (the dotted line in Figure 12).
Now we would like to discuss the dependence of the main characteristics of a star after a tidal encounter on the parameter $`\eta `$ and the polytropic index $`n`$. For that we should introduce some quantities describing the perturbed star which are conserved when the star leaves the region of effective action of the tidal forces ($`\tau \mathrm{}`$). One such quantity is the difference between the total energy of the gravitationally bounded debris of the star and the total energy of the unperturbed star: $`E_{tot}^{bn}`$. This quantity sharply decreases with increase of the parameter $`\eta `$, and therefore it is convenient to introduce another quantity $`T(\eta )=\eta ^4E_{tot}^{bn}`$. $`T(\eta )`$ can be easily compared with results of linear theory of tidal perturbations of a star (Press $`\&`$ Teukolsky 1977, Lee $`\&`$ Ostriker 1986), and with results of 3D simulations. We show the result of calculations of this quantity in Figures 13 for a polytropic star with $`n=3`$, in Figure 14 for the case $`n=2`$, and in Figure 15 for the case $`n=1.5`$, respectively <sup>7</sup><sup>7</sup>7 For the calculations reported hereafter we used the time interval $`10<\tau <10`$.. One can see that the difference between the linear theory and our model is small if $`\eta `$ is rather large ($`\eta 2`$ for the $`n=3`$ case, and $`3.54`$ for the $`n=2`$, $`1.5`$ cases), but the linear theory underestimates significantly the energy gain at intermediate values of $`\eta `$ ($`\eta 1,1.5,2`$ for the $`n=3`$, $`2`$, $`1.5`$ cases respectively). This fact is confirmed by the 3D simulations (open circles in Figures 13-15), and is very transparent from a qualitative point of view. Indeed, the action of tidal forces on the star causes an increase of a characteristic size of the star. In turn, this increase leads to an increase of the characteristic amplitude of the tidal forces. Since the $`n=3`$ polytrope is more concentrated toward the stellar center of mass than the $`n=2`$ case, and the $`n=2`$ polytrope is more concentrated than the $`n=1.5`$ case, for a given $`\eta `$ the energy gain of the $`n=3`$ polytrope is always smaller than the energy gain of the $`n=2`$, and the energy gain of the $`n=2`$ polytrope is respectively smaller than the energy gain of the $`n=1.5`$ case. Note, that the results corresponding to large values of $`\eta `$ should be taken with care. At first our model may not give exactly the same results as the linear theory in the limit of small tidal perturbations. Second in our numerical model the difference between the numerical value of the total energy for an unperturbed star and its theoretical value is of the order of $`10^2`$, and this difference is larger than the energy gain corresponding to large values of $`\eta `$ (e.g $`\eta 2`$ for the $`n=3`$ case). Also the circle corresponding to $`\eta =0.5`$ is shown for illustrative purposes only. For those values of $`\eta `$ a significant stripping of the star occurs (see Figure 17), and a deviation of the total energy of the star (shown by the circle) and the total energy of the gravitationally bound debris can be significant. In Figure 14 the open circles almost coincide with our curve except for the circle corresponding to $`\eta =1.5`$, where the difference is rather large (about 40 percent). Perhaps this difference is related to effect of the partial mass stripping of the star, which is rather pronounced for that case. This effect makes the determination of energy of the gravitationally bound debris rather ambiguous, since the gas leaves the star with almost parabolic velocities for this tidal encounter of a moderate strength.
In Figures 16 we show the values of the angular momentum perpendicular to the starโs orbit (taken at the end of calculations $`\tau _{end}=10`$) and for the gravitationally bound debris (the dotted line) versus $`\eta `$ for the $`n=3`$ polytrope. The open circles correspond to the 3D simulations and the dashed curve represents the same quantity calculated by Diener et al 1995 in the framework of the affine model. It seems that in our model the star gets more angular momentum than in the affine model.
In general Figures 13-16 show very good agreement between our results and the results of the 3D simulations. Unfortunately the number of the 3D experiments is too small to make a robust quantitative comparison.
In Figure 17 we show the amount of mass lost by the star after a fly-by around the black hole as a function of $`\eta `$. Contrary to a โnaiveโ Roche-like criterion of tidal disruption (the discussion after eq. 34) we see that there is a partial stripping of mass from the star for a particular range of $`\eta `$. For the $`n=3`$ polytrope (the solid line) the star loses its mass if $`\eta <\eta _{strip}1.5`$, and the star is completely disrupted if $`\eta <\eta _{crit}0.4`$. For the $`n=2`$ polytrope (the dashed line) we have $`\eta _{strip}2`$ and $`\eta _{crit}0.9`$ and for the $`n=1.5`$ polytrope (the dotted line) we have $`\eta _{strip}2.5`$ and $`\eta _{crit}1.14`$. The values of $`\eta _{strip}`$ and $`\eta _{crit}`$ are in excellent agreement with estimates based on the 3D simulations (see Kh b). It is instructive to compare these results with a criterion of tidal disruption obtained in the affine model. In this model Diener et al found that the star is disrupted if $`\eta <\eta _{crit}^a=0.844`$ for $`n=3`$, $`\eta <\eta _{crit}^a=1.482`$ for $`n=2`$, $`\eta <\eta _{crit}^a=1.839`$ for $`n=1.5`$. Since in the affine model the velocity field is self-similar, in this model a partial stripping of outer layers of the star cannot occur. It can be observed that the inequalities $`\eta _{crit}<\eta _{crit}^a<\eta _{strip}`$ always hold, and the difference between the values of $`\eta _{crit}^a`$ and values $`\eta _{crit}`$ is about 40-50 percent. This difference is much larger than the difference between our values and the results of the 3D calculations. This demonstrates the advantage of our model.
## 4 Conclusions and Discussion
We have constructed a new model of a star perturbed by a tidal field. In this model the star consists of a set of elliptical shells which in general have different principal axes and different orientations with respect to a fixed locally inertial reference frame comoving with the starโs center of mass. The model obeys certain evolution equations. The results of calculations of a simple problem of tidal encounter of a polytropic star moving on a parabolic orbit around a source of Newtonian gravity have been compared with the results of three dimensional finite difference simulations. We found that the main characteristics of the tidal encounter agree with the results of the finite difference approach with a typical accuracy $`1020`$ percent. Taking into account that the astrophysical applications do not demand very high accuracy in description of a single tidal encounter, we think that our model could be used in order to investigate all possible variants of the problem of tidal interaction between a supermassive black hole and a star interesting from an astrophysical point of view. The main advantage of our model is its effectively one dimensional character, which allows us to calculate all interesting variables much faster than the 3D approach, and over a much longer time of evolution. Also, all characteristics of the tidally perturbed star could be inferred from our model in a straightforward and unambiguous way. Our model could also be used in a study of a rotating single star, or a binary star.
In principal, the agreement between our model and 3D computations could be improved if one uses more advanced variants of our model (see below). However we would like to note that in order to make a comparison between an advanced variant of our model and 3D computations, one should also increase both the number of the numerical experiments and their resolution. For example, recently it has been claimed (Ayal et al, 2000) that the difference between the different numerical experiments in SPH models is about 10 percent depending on number of particles used in the calculations. We are not aware of similar convergence studies for the 3D finite difference models. We think that such studies should be undertaken parallel to work on improvement of our model.
Now we would like to discuss possible extensions of our formalism. One obvious way for such extension consists in using a more complicated stellar model for the unperturbed state. Then the realistic stellar models could be generalized to our problem by using the energy evolution equation and the virial relations similar to the eq. (8,9), but written for a realistic stellar gas, with possible inclusion of e.g. non-adiabatic effects, viscosity, effects of radiative transfer, and generation of heat due to nuclear reactions in the stellar core. The evolution equations for such a model could be derived from the energy equation and the virial relations in the way described above. One could also use a more refined numerical scheme, say an implicit scheme with a nonuniform grid. We suppose that certain powerful methods developed in numerical investigations of pulsating stars could be directly applied to our problem. Another interesting extension consists in using the real distribution of the pressure and the density over the volume of the star in our approximation. Thus it is possible to construct a model with no intersection between the shells, which could provide more information about displacements of particular elements of the star during the tidal encounter. Assuming that the star consists of an ideal gas with constant ratio of specific heats $`\gamma `$, and neglecting the possible presence of shocks in the system, this problem is reduced to the evaluation of the following integrals:
$$I_1=_SR^{2(\gamma 1)}๐\mathrm{\Omega },$$
$`35`$
and
$$I_{2i}=_Sl_{(i)}^2R^{2\gamma }๐\mathrm{\Omega },$$
$`36`$
where the integration is performed over the surface of an ellipsoid $`\lambda _{(i)}y_{(i)}^2=1`$. $`\lambda _i>0`$ are the eigenvalues of the matrix $`R^{ln}`$ (see. eq. 4), $`R`$ is the value of the radius vector joining the center of the ellipsoid to a point on its surface, $`l_i`$ are the direction cosines of the radius vector, and $`d\mathrm{\Omega }`$ is the elementary solid angle. The Cartesian coordinates $`y_i`$ are associated with the frame of eigenvectors of the matrix $`R^{ln}`$. If $`\gamma =1`$ the integral (35) is trivial, and the integrals (36) are related to the integrals $`D_i`$ defined above (eq. 18), but with $`\lambda _j^1`$ playing the role of $`a_j^2`$. One can see that the pressure tensor (defined as the left hand side of the eq. 13) can be evaluated with help of the integrals (35, 36). The volume part of the pressure tensor (the first term on l.h.s. of the eq. (13)) is obtained by integration over the mass of a quantity proportional to the integral (35) with the upper limit of integration determined by some given Lagrangian coordinate $`x_{}`$. The surface term ( the second term on l.h.s. of the eq. (13)) is expressed with help of the integrals (36) with $`\lambda _i=\lambda _i(x_{})`$. Note that now the surface term is not symmetric, and therefore transfer of angular momentum between the neighboring shells due to pressure is allowed. If some shells are close to intersection, the density at some particular value of $`x`$ tends to infinity. That means that some eigenvalues $`\lambda _i`$ go to zero, and as a consequence the integrals (35), (36) tend to infinity causing an increase of pressure. In turn the increase of pressure could prevent the shells from intersecting. The integrals (35), (36) can be evaluated e.g. by numerical means, and serve as main building blocks of a model without intersections between the shells corresponding to different Lagrangian radii.
Finally one can generalize our model to the case of a relativistic tidal field. In a separate paper we apply our model to the problem of tidal interaction of a star with a supermassive Kerr black hole.
We are grateful to P. Diener, I. Igumenshchev, A. Khokhlov, O. Rognvaldsson, and V. Turchaninov for useful remarks concerning the numerical methods, to J. Schmalzing for very useful comments, to E. Kotok for the help in preparation of this paper, and also to the referee for very helpful suggestions. It is our sincere pleasure to thank A. Doroshkevich and V. Frolov for many stimulating and fruitful discussions. This work was supported in part by RFBR grant N 00-02-161335 and in part by the Danish Research Foundation through its establishment of the Theoretical Astrophysics Center.
## Appendix A Reduction of the dynamical equations to the equations of the affine model and the case of an incompressible fluid
Let us consider the position matrix of a special form:
$$๐(t,r_0)=\widehat{๐}(t)r_0,$$
$`A1`$
and introduce the matrices $`\widehat{๐}(t)`$, $`\widehat{๐}(t)`$, $`\widehat{๐}(t)`$, and the quantities $`\widehat{a}_i(t)`$ defined with the help of the matrix $`\widehat{๐}(t)`$ by analogy with the matrices $`๐`$, $`๐`$, $`๐`$, and the quantities $`a_i`$, respectively. For the position matrix of the form (A1) the matrix $`R^{ln}`$ is:
$$R^{ln}=\frac{\delta ^{ln}}{r_0},$$
$`A2`$
and the Jacobian $`D=\widehat{g}(t)\widehat{a}_1\widehat{a}_2\widehat{a}_3`$ does not depend on the spatial coordinate. Therefore in this case there is no difference between the exact and the averaged values of the density and the pressure. Substituting the expression (A1) into the equations (23), multiplying the result on $`r_0`$ and integrating over $`M`$, we obtain
$$\ddot{\widehat{T}_n^i}=\frac{\mathrm{\Pi }}{M_{}}\widehat{S}_i^n+\frac{\mathrm{\Omega }_{}}{2M_{}}\widehat{g}^1\widehat{A}_j^i\widehat{a}_j\widehat{D}_j\widehat{E}_n^j+C_j^i\widehat{T}_n^j,$$
$`A3`$
where $`\mathrm{\Pi }=๐M\frac{P}{\rho }`$, the constant factor $`M_{}=\frac{1}{3}๐Mr_0^2`$ is the scalar quadrupole moment of the star in the unperturbed state and the constant factor $`\mathrm{\Omega }_{}=G\frac{MdM}{r_0}`$ is the self-gravitational energy value in the unperturbed state. The quantities $`\stackrel{~}{D}_j`$ are defined with the help of the quantities $`\stackrel{~}{a}_j`$ by analogy with the quantities $`D_j`$. The equations (A2) are just the dynamical equations of the affine model (e.g. Carter $`\&`$ Luminet 1983, 1985 and Luminet $`\&`$ Carter, 1986).
Now let us consider the case of an incompressible fluid. In this case the density $`\rho `$ is not changing durind the motion and therefore the constraint $`\widehat{g}=1`$ must be imposed, see eq (6). Assuming that the density is also constant over the whole volume of the star $`\rho =\rho _0`$=const, the equations (A3) can be reduced to a very simple form. Taking into account the constraint $`\widehat{g}=1`$, and performing the integration over $`M`$ in the expressions for $`\mathrm{\Pi }`$, $`M_{}`$, and $`\mathrm{\Omega }_{}`$, we have
$$\ddot{\widehat{T}_n^i}=\frac{\widehat{S}_i^nP_c(t)}{2\rho _0R_{}^2}2\pi G\rho _0\widehat{A}_j^i\widehat{a}_j\widehat{D}_j\widehat{E}_n^j+C_j^i\widehat{T}_n^j,$$
$`A4`$
where $`R_{}`$ is the radius of the star and $`P_c(t)`$ is the pressure in the center of the star. Together with the constraint $`\widehat{g}=1`$ the equations (A4) form a complete set of equations. These equations have already been obtained by Luminet and Carter (see eq. (3.1a) in Luminet $`\&`$ Carter, 1986) <sup>8</sup><sup>8</sup>8Note that there is a misprint in the equation (3.1a). The factor $`\frac{1}{2\rho _0R_{}^2}`$ in the last term on the right hand side is missed.. It has been proved that the equations (A4) follow from the hydrodynamical equations without any approximation.
## Appendix B The numerical scheme
For our numerical work it is convenient to introduce the potential energy tensor density per unit of mass:
$$F^{ik}=\frac{d}{dM}\overline{๐ซ}^{ik}=\frac{1}{2}GMA_l^iA_l^k\frac{a_l^2D_l}{g},$$
$`B1`$
and use dimensionless representations of all variables. We use the dimensionless mass coordinate $`x=M/M_{}`$, and the dimensionless time $`\tau =\sqrt{\frac{GM_{}}{R_{}^3}}t`$. $`\mathrm{\Delta }x`$ stands for the mass coordinate step, and $`\mathrm{\Delta }\tau `$ stands for the time step. The dimensionless dynamical variables are defined as follows:
$$\stackrel{~}{๐}=๐/R_{},\stackrel{~}{๐}=\frac{}{\tau }\stackrel{~}{๐},$$
$`B2`$
$$\stackrel{~}{๐
}=\frac{R_{}}{GM_{}}๐
,$$
$`B3`$
$$\stackrel{~}{g}=g/R_{}^3,\stackrel{~}{\rho }=\frac{4\pi R_{}^3}{3M_{}}\rho ,$$
$`B4`$
$$\stackrel{~}{P}=\frac{4\pi R_{}^4}{GM_{}^2}P,\stackrel{~}{ฯต}=\frac{\stackrel{~}{P}}{(\gamma 1)},$$
$`B5`$
the dimensionless tidal tensor is
$$\stackrel{~}{๐}=\frac{R_{}^3}{GM_{}}๐.$$
$`\mathrm{๐๐}`$
We use a uniform grid along the mass coordinate with number of the grid points $`J=201`$ (the first and the last numbers correspond respectively to the center of the star, and to the starโs boundary). The pressure, the density and the energy density are defined at the centers of zones between the grid points, and are denoted by the indices $`j+1/2`$, $`j1/2`$, and the position matrix and its time derivative are defined at the grid points. Hereafter the index $`n`$ denotes a time level, and the index $`j`$ denotes a particular grid point, and no summation over those indices is assumed.
We use several intermediate steps to advance our model from $`\tau ^n`$ to $`\tau ^{n+1}`$. First we calculate $`\stackrel{~}{๐}`$, and $`\stackrel{~}{g}_j^n`$ and $`\stackrel{~}{๐
}_j^n`$ with help of special subroutines. To calculate the rotational matrix $`๐`$ and the principal axes $`a_l`$ we use the standard Jacobi method. Then we calculate the dimensionless expansion rate
$$\stackrel{~}{H}_j^n=\frac{1}{3}(\stackrel{~}{๐}\stackrel{~}{๐})_j^n,$$
$`B7`$
where $``$ stands for the contraction over all indices, and the density
$$\stackrel{~}{\rho }_{j1/2}^n=\frac{\mathrm{\Delta }x}{(\stackrel{~}{g}_j^n\stackrel{~}{g}_{j1}^n)},$$
$`B8`$
Next the artificial (bulk) viscosity term is calculated according to the rule:
$$v_j^n=(\stackrel{~}{H}\stackrel{~}{g}^{1/3})_j^n,q_{j1/2}^n=b^2\stackrel{~}{\rho }_{j1/2}^n(\mathrm{\Delta }x)^2(v_j^nv_{j1}^n)^2,$$
$`B9`$
if $`v_j^nv_{j1}^n<0`$, and
$$q_{j1/2}^n=0,$$
if $`v_j^nv_{j1}^n>0`$. The artificial viscosity is very useful in the calculations, since it dampen the small scale (unphysical) modes. We used $`b=2`$ in our calculations, and checked that smaller values of $`b`$ lead to the same motion, but with additional small scale oscillations of the dynamical variables. After calculation of all these variables we advance the position matrix and its time derivative to the next time level:
$$\stackrel{~}{๐ฏ}_j^{n+1}=\stackrel{~}{๐ฏ}_j^n+\mathrm{\Delta }\tau ((\stackrel{~}{๐}^T\stackrel{~}{g})_j^n\frac{\stackrel{~}{P}}{x}+3(\stackrel{~}{๐}^T\stackrel{~}{๐
})_j^n+(\stackrel{~}{๐}\stackrel{~}{๐})_j^n),$$
$`B10`$
where
$$\frac{\stackrel{~}{P}}{x}=\frac{(\stackrel{~}{P}_{j+1/2}^n+q_{j+1/2}^n\stackrel{~}{P}_{j1/2}^nq_{j1/2}^n)}{\mathrm{\Delta }x},$$
and
$$\stackrel{~}{T}_j^{n+1}=\stackrel{~}{T}_j^n+\mathrm{\Delta }\tau \stackrel{~}{๐ฏ}_j^{n+1}.$$
$`B11`$
Note, that the scheme (B10-B11) conserves the circulation, and it also conserves the angular momentum provided the tidal term is โswitched offโ. To calculate the change of the energy density and the pressure, we calculate the kinetic energy term $`\frac{\dot{๐}\dot{๐}^T}{2}`$, and the potential energy term $`tr(๐
)`$ for the both time levels, and also the density $`\stackrel{~}{\rho }_{j1/2}^n`$ for the $`n+1`$ time level. Then we have
$$\stackrel{~}{ฯต}_{j1/2}^{n+1}=\stackrel{~}{\rho }_{j1/2}^{n+1}(\frac{\stackrel{~}{ฯต}_{j1/2}^n}{\stackrel{~}{\rho }_{j1/2}^n}(\mathrm{\Delta }E_k+\mathrm{\Delta }E_p)+\mathrm{\Delta }\tau Q),$$
$`B12`$
where
$$\mathrm{\Delta }E_k=(\frac{\dot{๐}\dot{๐}^T}{2})_j^{n+1}(\frac{\dot{๐}\dot{๐}^T}{2})_j^n\mathrm{\Delta }E_p=3tr(๐
_j^{n+1}๐
_j^n),$$
$`B13`$
and
$$Q=3\frac{(\stackrel{~}{H}_j^n\stackrel{~}{g}_j^n(\stackrel{~}{P}_{j+1/2}^n+q_{j+1/2}^n)\stackrel{~}{H}_{j1}^n\stackrel{~}{g}_{j1}^n(\stackrel{~}{P}_{j1/2}^n+q_{j1/2}^n))}{\mathrm{\Delta }x}+(\stackrel{~}{๐}\stackrel{~}{๐}\stackrel{~}{๐}^T)_j^n.$$
$`B14`$
The eq. (B12) tells that the scheme is conservative with respect to the law of energy conservation.
The boundary conditions for our problem are
$$๐_1^n=๐_1^n=0,\stackrel{~}{\rho }_{1/2}^n=\stackrel{~}{\rho }_{3/2}^n,$$
$`B15`$
and
$$\stackrel{~}{P}_J^n=0.$$
$`B16`$
The initial static stellar configurations are calculated by a shooting method described by e.g. Khokhlov 1991.
The size of the time step must follow from the stability analysis of our scheme, which is rather complicated. Therefore we constrain our time step by the condition:
$$\mathrm{\Delta }\tau =\frac{\alpha \mathrm{\Delta }x}{c_s},$$
$`B17`$
where
$$c_s=\frac{\mathrm{\Delta }x}{\stackrel{~}{g}\sqrt{\gamma (\stackrel{~}{P}+q)\stackrel{~}{\rho }(\stackrel{~}{๐}\stackrel{~}{๐}^T)}},$$
$`B18`$
is the velocity of propagation of a small perturbation (with respect to the mass coordinate) calculated in analytical linear approximation, and $`\alpha `$ is a parameter. We used $`\alpha =1/15`$ for all our calculations except the case $`n=2`$, $`\eta =0.1`$, where $`\alpha =1/60`$ has been used, and the case $`n=1.5`$, $`\eta =0.0894`$, where we used $`\alpha =1/120`$. In the case of the evolution of a pulsating star our condition (B17) is reduced to the well known von Neumann-Richtmyer condition provided $`\alpha =1/\sqrt{(4\pi )}`$ (e.g Richtmyer $`\&`$ Morton, 1967, p. 297-298).
We tested our scheme against Sedovโs analytical solution for a motion of an expanding spherical gas cloud, and also the problem of small oscillations in a polytropic star. In the both cases we found very good agreement between the analytical theory and our scheme.
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# 1 Fractional Time Evolution
## 1 Fractional Time Evolution
A series of recent investigations has found that, in a suitable long time limit, the macroscopic time evolution $`T^t`$ of a physical observable $`X(t)`$ is given as a convolution of the form
$$T_\alpha ^tX(t_0)=\underset{0}{\overset{\mathrm{}}{}}X(t_0s)h_\alpha \left(\frac{s}{t}\right)\frac{ds}{t}$$
(1)
where $`t0`$, $`0<\alpha 1`$ and
$$h_\alpha (x)=\frac{1}{x\alpha }H_{11}^{10}\left(\frac{1}{x}|\begin{array}{cc}(0,1)& \\ (0,1/\alpha )& \end{array}\right)$$
(2)
is defined through its Mellin transform
$$_0^{\mathrm{}}H_{11}^{10}\left(x|\begin{array}{cc}(0,1)& \\ (0,1/\alpha )& \end{array}\right)x^{s1}๐x=\frac{\mathrm{\Gamma }(s/\alpha )}{\mathrm{\Gamma }(s)}$$
(3)
For $`\alpha =1`$
$$h_1(x)=\underset{\alpha 1^{}}{lim}h_\alpha (x)=\delta (x1)$$
(4)
and hence eq. (1) reduces to
$$T_1^tX(t_0)=X(t_0t)$$
(5)
a simple translation. For fixed $`\alpha `$ the operators $`T_\alpha ^t,t0`$ form a one-parameter semigroup on a Banach space (e.g. the space of all continuous functions vanishing at infinity).
Much interest in the semigroups $`T_\alpha ^t`$ derives from the fact that their infinitesimal generators are proprotional to
$$(_+^\alpha X)(t)=\underset{s0}{lim}\frac{T_\alpha ^sX(t)X(t)}{s}=\frac{1}{\mathrm{\Gamma }(\alpha )}\underset{0}{\overset{\mathrm{}}{}}\frac{X(t)X(ts)}{s^{\alpha +1}}๐s,$$
(6)
i.e. to fractional time derivatives of order $`\alpha `$ . This may be seen from Ref. (p. 302), and the fact that the kernels $`h_\alpha (s/t)/t`$ in eq. (1) are stable probability densities (see also ).
Derivatives of noninteger order $`0<\alpha 1`$ with respect to time are in this way found to furnish an important generalization for the infinitesimal generator of the macroscopic time evolution of physical observables.
Given the generality of the results above it is natural to consider specific examples by replacing the integer order time derivative with a fractional derivative. A particularly interesting example is given by diffusion and random walks . The fractional generalization of the familiar master equation for random walks turns out to be intimately related to the theory of continuous time random walks with power-law tails in the waiting time density . It is the purpose of this note to exhibit an exact relationship between a fractional master equation and continuous time random walks .
Let me conclude this introduction with the remark that only a special subclass of separable continuous time random walks with algebraically decaying waiting time distributions corresponds exactly to a generalized master equation with Riemann-Liouville-type fractional time derivative. It would be interesting to specify precisely whether, and under which conditions, other waiting time densities with a power-law tail lead to the same result .
## 2 Fractional Master Equation
Consider a random walk and let $`p(\stackrel{}{r},t)`$ denote the probability density to find the walker at the position $`\stackrel{}{r}^d`$ at time $`t`$ if it was at the origin $`\stackrel{}{r}=0`$ at time $`t=0`$. The positions $`\stackrel{}{r}^d`$ may be discrete or continuous. The conventional master equation for the random walk reads
$$\frac{}{t}p(\stackrel{}{r},t)=\underset{\stackrel{}{r}^{}}{}w(\stackrel{}{r}\stackrel{}{r}^{})p(\stackrel{}{r}^{},t)$$
(7)
with initial condition $`p(\stackrel{}{r},0)=\delta _{\stackrel{}{r}0}`$. To establish the fractional generalization of this initial value problem one must respect the nonlocal nature of the fractional derivatives. This leads to the fractional master equation
$$_{0+}^\alpha p(\stackrel{}{r},t)=\delta _{\stackrel{}{r}0}\frac{1}{\mathrm{\Gamma }(1\alpha )t^\alpha }+\underset{\stackrel{}{r}^{}}{}w(\stackrel{}{r}\stackrel{}{r}^{})p(\stackrel{}{r}^{},t)$$
(8)
where $`0<\alpha 1`$, $`t0`$, and the initial condition $`p(\stackrel{}{r},0)=\delta _{\stackrel{}{r}0}`$ has been incorporated. Note that the fractional transition rates $`w(\stackrel{}{r})`$ have units of (1/time$`)^\alpha `$. They obey the relation $`_\stackrel{}{r}w(\stackrel{}{r})=0`$. Applying a fractional integral of order $`\alpha `$ to equation (8) yields an integral equation
$$p(\stackrel{}{r},t)=\delta _{\stackrel{}{r}0}+\frac{1}{\mathrm{\Gamma }(\alpha )}_0^t(tt^{})^{\alpha 1}\underset{\stackrel{}{r}^{}}{}w(\stackrel{}{r}\stackrel{}{r}^{})p(\stackrel{}{r}^{},t^{})dt^{}$$
(9)
reminiscent of the integral equation for continuous time random walks .
## 3 Continuous Time Random Walks
The basic integral equation for separable continous time random walks describes a random walker in continous time without correlation between its spatial and temporal behaviour, and reads
$$p(\stackrel{}{r},t)=\delta _{\stackrel{}{r}0}\mathrm{\Phi }(t)+_0^t\psi (tt^{})\underset{\stackrel{}{r}^{}}{}\lambda (\stackrel{}{r}\stackrel{}{r}^{})p(\stackrel{}{r}^{},t^{})dt^{}.$$
(10)
Here, as in (9), $`p(\stackrel{}{r},t)`$ denotes the probability density to find the walker at the position $`\stackrel{}{r}^d`$ at time $`t`$ if it started from the origin $`\stackrel{}{r}=0`$ at time $`t=0`$. $`\lambda (\stackrel{}{r})`$ is the probability for a displacement $`\stackrel{}{r}`$ in each single step, and $`\psi (t)`$ is the waiting time distribution giving the probability density for the occurrence of a time interval $`t`$ between two consecutive steps. The transition probabilities obey $`_\stackrel{}{r}\lambda (\stackrel{}{r})=1`$. The function $`\mathrm{\Phi }(t)`$ in eq. (10) is the survival probability at the initial position. It is related to the waiting time distribution through
$$\mathrm{\Phi }(t)=1_0^t\psi (t^{})๐t^{}.$$
(11)
Note that in equation (10) the system is prepared with the walker at position $`\stackrel{}{r}=0`$, and it develops from there according to $`\psi (t)`$.
## 4 Relation between equations (9) and (10)
The similarity between equations (9) and (10) suggests that the former is a special case of the latter. To show that this is true let
$$\psi (u)=\text{L}\{\psi (t)\}(u)=_0^{\mathrm{}}e^{ut}\psi (t)๐t$$
(12)
denote the Laplace transform of $`\psi (t)`$ and write
$$\lambda (\stackrel{}{q})=\text{F}\{\lambda (\stackrel{}{r})\}(\stackrel{}{q})=\underset{\stackrel{}{r}}{}e^{i\stackrel{}{q}\stackrel{}{r}}\lambda (\stackrel{}{r})$$
(13)
for the Fourier transform of $`\lambda (\stackrel{}{r})`$. Then the Fourier-Laplace transform $`p(\stackrel{}{q},u)`$ of the solution to (10) is given as
$$p(\stackrel{}{q},u)=\frac{1}{u}\frac{1\psi (u)}{1\psi (u)\lambda (k)}=\frac{\mathrm{\Phi }(u)}{1\psi (u)\lambda (\stackrel{}{q})}$$
(14)
where $`\mathrm{\Phi }(u)`$ is the Laplace transform of the survival probability.
Similarly the fractional master equation (9) can be solved in Fourier-Laplace space with the result
$$p(\stackrel{}{q},u)=\frac{u^{\alpha 1}}{u^\alpha w(\stackrel{}{q})}$$
(15)
where $`w(\stackrel{}{q})`$ is the Fourier transform of the kernel $`w(\stackrel{}{r})`$ in(9). Eliminating $`p(\stackrel{}{q},u)`$ between (14) and (15) gives the result
$$\frac{1\psi (u)}{u^\alpha \psi (u)}=\frac{\lambda (\stackrel{}{q})1}{w(\stackrel{}{q})}=C$$
(16)
where $`C`$ is a constant. The last equality holds because the left hand side of the first equality is $`\stackrel{}{q}`$-independent while the right hand side is independent of $`u`$.
From (16) it is seen that the fractional master equation characterized by the kernel $`w(\stackrel{}{r})`$ and the order $`\alpha `$ corresponds to a special class of space time decoupled continuous time random walks characterized by $`\lambda (\stackrel{}{r})`$ and $`\psi (t)`$. This correspondence is given precisely as
$$\psi (u)=\frac{1}{1+Cu^\alpha }$$
(17)
and
$$\lambda (\stackrel{}{q})=1+Cw(\stackrel{}{q})$$
(18)
with the same constant $`C`$ appearing in both equations. Not unexpectedly the correspondence defines the waiting time distribution uniquely up to a constant while the structure function is related to the Fourier transform of the transition rates.
To invert the Laplace transformation in (17) and exhibit the form of the waiting time density $`\psi (t)`$ in the time domain one rewrites eq. (17)
$$\psi (u)=\frac{u^\alpha }{C}\frac{1}{1+\frac{u^\alpha }{C}}=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1)^k}{C^{k+1}}u^{\alpha (k+1)}$$
(19)
and inverts the series term by term. This yields the result
$$\psi (t;\alpha ,C)=\frac{t^{\alpha 1}}{C}\underset{k=0}{\overset{\mathrm{}}{}}\frac{1}{\mathrm{\Gamma }(\alpha k+\alpha )}\left(\frac{t^\alpha }{C}\right)^k.$$
(20)
First one notes that for $`\alpha =1`$ the result reduces to
$$\psi (t;1,C)=\frac{1}{C}\mathrm{exp}(t/C)$$
(21)
the familiar exponential form. For $`0<\alpha <1`$ the result is recognized as
$$\psi (t;\alpha ,C)=\frac{t^{\alpha 1}}{C}E_{\alpha ,\alpha }\left(\frac{t^\alpha }{C}\right).$$
(22)
where the function
$$E_{\alpha ,\beta }(z)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{z^k}{\mathrm{\Gamma }(\alpha k+\beta )}$$
(23)
is the generalized Mittag-Leffler function .
The asymptotic behaviour of $`\psi (t)`$ for $`t0`$ is readily obtained by noting that $`E_{\alpha ,\alpha }(0)=1`$. Hence $`\psi (t)`$ behaves as
$$\psi (t)t^{1+\alpha }$$
(24)
for $`t0`$. For $`0<\alpha <1`$ the waiting time density becomes singular at the origin.
The asymptotic behaviour for large waiting times $`t\mathrm{}`$ is obtained from the asymptotic series expansion for the Mittag-Leffler function
$$E_{\alpha ,\beta }(z)=\underset{n=1}{\overset{N}{}}\frac{z^n}{\mathrm{\Gamma }(\beta \alpha n)}+O(|z|^N)$$
(25)
valid for $`|`$arg$`(z)|<(1(\alpha /2))\pi |`$ and $`z\mathrm{}`$. It follows from this that $`E_{\alpha ,\beta }(x)x^2`$ for $`x\mathrm{}`$ and hence that
$$\psi (t)t^{1\alpha }$$
(26)
for large $`t\mathrm{}`$ and $`0<\alpha <1`$. This result shows that the waiting time distribution has an algebraic tail as it is usually assumed in the theory of continuous time random walks .
In summary it has been shown that the master equation with a Riemann-Liouville fractional time derivative of order $`\alpha `$ corresponds exactly to a continuous time random walk whose waiting time density is related to the generalized Mittag-Leffler function and exhibits a power-law tail. It would be interesting to know the precise conditions under which the more general class of separable continuous time random walks obeying eq. (26) โapproximatesโ this result.
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# Some Remarks Concerning the Feynman โIntegral over All Pathsโ Method
## Abstract
Suppose we have two nonequivalent but s-equivalent Lagrange functions, the question arises: are they both equally well fitted for the Feynman quantization procedure or do they lead to two different quantization schemes.
Dedicated to Roman S. Ingarden on the occasion of his 80th birthday
1. The goal of this note is to exhibit the following problem. It is well known that in the quantization prescription, based on the Feynman โintegral over all pathsโ the classical Lagrange function is used in the exponent of the integrand of the Feynman integral. The physical content of a dynamical system is, however, mainly characterized by the equations of motion of this systems; the Lagrange function, if such one exists at all for these equations, plays a secondary rรดle, as there can be many nonequivalent Lagrange functions linked to equations of motion (Euler Lagrange Equations), yielding the same set of solutions - so called s-equivalent equations.
The question arises: suppose we have two nonequivalent but s-equivalent Lagrange functions, are they both equally well fitted for the Feynman quantization procedure or do they lead to two different quantization schemes.
2. To begin with let us consider the case of one classical particle in a (1+1)-dimensional space-time and the largest set of s-equivalent Lagrange functions, corresponding to the equation of motion of this particle. We do not need to specify the form of this equation; to each equation written in the normal form, viz.
$$\ddot{x}=f(x,\dot{x},t)$$
(1)
corresponds always a Lagrange function . The inverse problem for the case of (1+1) dimensions was treated extensively by many scientists , .
It is known that the most general form of an autonomous Lagrange function, s-equivalent to a given autonomous Lagrange function $`L(x\dot{x})`$, the form of which we do not specify, is
$$L^{}=\dot{x}\underset{c}{\overset{\dot{x}}{}}G(x,u)๐u\mathrm{\Sigma }(H)$$
(2)
where
$$G(x,\dot{x})\frac{d\mathrm{\Sigma }(H)}{dH}\frac{^2L}{\dot{x}^2},$$
(3)
$$H\dot{x}\frac{L}{\dot{x}}L,$$
(4)
and $`\mathrm{\Sigma }(z)`$ is an arbitrary differentiable function of $`z`$. The constant $`c`$ is so chosen that the integral on the r.h.s. of (2) does not diverge<sup>1</sup><sup>1</sup>1we could even assume $`c=c(x)`$.. The Hamilton function reads
$$H^{}=\dot{x}\frac{L^{}}{\dot{x}}L^{}=\mathrm{\Sigma }(H)+const..$$
(5)
The Lagrange function $`L^{}`$ for different choices of $`\mathrm{\Sigma }`$, assuming $`\frac{d\mathrm{\Sigma }(z)}{dz}`$ is not a constant, are not equivalent to each other as well as to $`L`$; in other words they do not differ from each other by a function $`\frac{d\mathrm{\Phi }(t)}{dt}`$.
To make things more specific let us now specify the original Lagrange function $`L`$ as well as $`\mathrm{\Sigma }`$ and $`c`$, viz.
$$L=\frac{1}{2}\dot{x}^2V(x),$$
(6)
$$\mathrm{\Sigma }(H)=\frac{1}{2}H^2,$$
(7)
$$c=0.$$
(8)
Then
$$H=\frac{1}{2}\dot{x}^2+V(x),$$
(9)
$$L^{}=\frac{1}{24}\dot{x}^4+\frac{1}{2}\dot{x}^2V\frac{1}{2}V^2,$$
(10)
$$H^{}=\frac{1}{2}H^2=\frac{1}{8}\dot{x}^4+\frac{1}{2}\dot{x}^2V+\frac{1}{2}V^2,$$
(11)
and<sup>2</sup><sup>2</sup>2With the notation $`\dot{x}=z`$ we have
$$\frac{dz}{dp^{}}=\frac{1}{H}=(2H^{})^{\frac{1}{2}},$$
$`(A)`$
$$\frac{d^2z}{dp^2}=\frac{z}{H^3}.$$
For $`z`$ becoming large $`\frac{dz}{dp^{}}`$ vanishes like $`z^2`$ and $`\frac{d^3z}{dp^2}`$ like $`z^5`$. Using the canonical Hamilton equation
$$z=\frac{H^{}}{p^{}}$$
and taking into account (A) we get the equation for $`H^{}`$, viz.
$$\frac{^2H^{}(x,p^{})}{p^2}\frac{1}{\sqrt{2}}\frac{1}{\sqrt{H^{}(x,p^{})}}=0.$$
$`(B)`$ The particular solution of (B) independent of $`x`$ reads
$$\widehat{H}^{}=\left(\frac{81}{32}\right)^{\frac{1}{3}}p^{\frac{4}{3}}$$
which corresponds to large $`p^{}`$ and $`\dot{x}`$ and discarding $`V(x)`$. The application of the first order perturbative procedure for small $`V`$ and $`\frac{dV}{dx}`$ as well as the use of canonical Hamilton equations yields
$$\dot{x}=\left(6p^{}\right)^{\frac{1}{3}}2\left(6p^{}\right)^{\frac{1}{3}}V$$
and
$$H^{}=\left(\frac{81}{32}\right)^{\frac{1}{3}}p^{\frac{4}{3}}\left[\left(\frac{9}{2}\right)^{\frac{1}{2}}V+a\right]p^{\frac{2}{3}}$$
where $`a`$ is a small number.
$$p^{}\frac{L^{}}{\dot{x}}=\frac{1}{6}\dot{x}^3+\dot{x}V(x).$$
(12)
Relation (12) is an algebraic equation of third degree with respect to
$$\dot{x}(p^{},x)=\dot{x}(p^{},x).$$
For
$$p^{}=0\text{and}V(x)>0$$
(13)
we have three roots of (12)
$$\dot{x}_1=0,\dot{x}_{2,3}=\pm i\sqrt{6V(x)}.$$
(14)
For obvious reasons we choose the real solution. In case $`V`$ is not always positive but it is bounded from below we may change $`V`$ in (6) by adding to it a properly chosen constant so that $`V`$ is then always positive.
The solution of (12) reads
$$\dot{x}=\frac{p^{}}{V}\frac{1}{6}\frac{1}{V}\left(\frac{p^{}}{V}\right)^3+\frac{1}{12}\frac{1}{V^2}\left(\frac{p^{}}{V}\right)^5\frac{1}{18}\frac{1}{V^3}\left(\frac{p^{}}{V}\right)^7+o\left(\left(\frac{p^{}}{V}\right)^9\right).$$
(15)
Notice that the few first terms of (15) coincide with
$$\frac{p^{}}{V}\left[1+\frac{1}{6}\mathrm{ln}\left(1\frac{1}{V}\left(\frac{p^{}}{V}\right)^2\right)\right].$$
(16)
For large $`\dot{x}`$ and $`p^{}`$
$$\dot{x}=\left(6p^{}\right)^{\frac{1}{3}}.$$
(17)
We have
$$H^{}=\frac{1}{2}V^2+\frac{1}{2}\frac{p_{}^{}{}_{}{}^{2}}{V}\frac{1}{24}\frac{p_{}^{}{}_{}{}^{4}}{V^4}+o\left(\left(\frac{p^{}}{V}\right)^8\right).$$
(18)
3. Let us now investigate the quantal case of one particle presented in the language of Feynmanโs approach.
It is well known , that in case the Hamiltonian function consists of two terms from which one depends only on $`p`$ and the other one only on $`x`$, the formula of Feynmanโs โintegral over all pathsโ with the classical Lagrange function in the exponent of the integral can be recovered from standard quantum mechanical approach.
To remind the Reader on this procedure let us consider the Hamiltonian function (9),<sup>3</sup><sup>3</sup>3We put the mass of the particle equal to one ($`m=1`$)., viz.
$$H=\frac{1}{2}p^2+V(x).$$
(19)
Starting from the first principles of Quantum Mechanics we have for the transition amplitude
$$\varphi (x^{},t_2|x,t_1)=x^{}|\mathrm{exp}\{i\widehat{H}(t_2t_1)\}|x,$$
(20)
where $`|`$ and $`|`$ denote the bra - and ket - states resp. and $`\widehat{H}`$ is the Hamilton operator
$$\widehat{H}\frac{1}{2}\widehat{p}^2+V(x),\widehat{p}=i\frac{}{x}.$$
(21)
We may write (20) as follows
$`x^{}|\mathrm{exp}\{i\widehat{H}t\}|x`$ $`=`$ $`\underset{\genfrac{}{}{0pt}{}{\mathrm{\Delta }t0}{\mathrm{\Delta }tn=t}}{lim}{\displaystyle ๐x_{n1}\mathrm{}๐x_1x^{}|e^{i\widehat{H}\mathrm{\Delta }t}|x_{n1}x_{n1}|\mathrm{}}`$ (24)
$`\mathrm{}|x_1x_1|e^{i\widehat{H}\mathrm{\Delta }t}|x.`$
If we use the formula
$$e^{(a+b)t}=\underset{n\mathrm{}}{lim}\left(e^{a\frac{t}{n}}e^{b\frac{t}{n}}\right)^n$$
(25)
then
$`x^{}|\mathrm{exp}\{i\widehat{H}t\}|x`$ $`=`$ $`\underset{\genfrac{}{}{0pt}{}{\mathrm{\Delta }t0}{\mathrm{\Delta }tn=t}}{lim}{\displaystyle ๐x_{n1}\mathrm{}๐x_1x^{}|e^{i\frac{\widehat{p}^2}{2}\mathrm{\Delta }t}|x_{n1}x_{n1}|\mathrm{}}`$ (28)
$`\mathrm{}|x_1x_1|e^{i\frac{\widehat{p}^2}{2}\mathrm{\Delta }t}|xe^{iV(x)t}.`$
Further we have
$`x^{}|e^{i\frac{\widehat{p}^2}{2}\mathrm{\Delta }t}|x`$ $`=`$ $`{\displaystyle ๐px^{}|e^{i\frac{\widehat{p}^2}{2}\mathrm{\Delta }t}|pp|x}`$ (29)
$`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle ๐pe^{i\frac{p^2}{2}\mathrm{\Delta }t}e^{ip(x^{}x)}}.`$ (31)
as
$$p|x=\left(\frac{1}{2\pi }\right)^{\frac{1}{2}}e^{ipx}.$$
(32)
Notice that
$`{\displaystyle \frac{i\mathrm{\Delta }t}{2}}p^2ip(x^{}x)`$ $`=`$ $`{\displaystyle \frac{i\mathrm{\Delta }t}{2}}\left(p^2+{\displaystyle \frac{2}{\mathrm{\Delta }t}}p(x^{}x)+{\displaystyle \frac{1}{(\mathrm{\Delta }t)^2}}(x^{}x)^2\right)`$ (35)
$`+{\displaystyle \frac{i}{2}}{\displaystyle \frac{(x^{}x)^2}{\mathrm{\Delta }t}}.`$
Consequently
$`x^{}|e^{i\frac{\widehat{p}^2}{2}\mathrm{\Delta }t}|x`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle ๐p\mathrm{exp}\left\{\frac{i\mathrm{\Delta }t}{2}\left(p+\frac{x^{}x}{\mathrm{\Delta }t}\right)^2\right\}\mathrm{exp}\left\{\frac{i}{2}\frac{(xx^{})^2}{\mathrm{\Delta }t}\right\}}`$ (36)
$`=`$ $`\left(2\pi i\mathrm{\Delta }t\right)^{\frac{1}{2}}\mathrm{exp}\left\{{\displaystyle \frac{i}{2}}\left({\displaystyle \frac{x^{}x}{\mathrm{\Delta }t}}\right)^2\mathrm{\Delta }t\right\}`$ (38)
where we used the saddle point method to evaluate
$$\frac{1}{2\pi }๐p\mathrm{exp}\left\{\frac{i\mathrm{\Delta }t}{2}\left(p+\frac{x^{}x}{\mathrm{\Delta }t}\right)^2\right\}=(2\pi i\mathrm{\Delta }t)^{\frac{1}{2}}.$$
(39)
Taking into account (28) and (38) we get eventually
$`\varphi (\underset{ยฏ}{\text{x}}^{},t_2|\underset{ยฏ}{\text{x}},t_1)`$ $`=`$ $`\underset{\genfrac{}{}{0pt}{}{n\mathrm{}}{n\mathrm{\Delta }t=t_2t_1}}{lim}{\displaystyle \underset{j=1}{\overset{n1}{}}}{\displaystyle ๐x_j\underset{k=1}{\overset{n}{}}(2\pi i\mathrm{\Delta }t)^{\frac{1}{2}}}`$ (42)
$`\mathrm{exp}\left\{i\left[\left({\displaystyle \frac{x_kx_{k1}}{\mathrm{\Delta }t}}\right)^2V(x)\right]\mathrm{\Delta }t\right\}`$
where $`x_nx^{}`$, $`x_0x`$. Thus in the exponent in (42) we have, indeed,
$$i\underset{t_1}{\overset{t_2}{}}L(x(t),\dot{x}(t))๐t,$$
(43)
is conjectured at the start.
The procedure presented above can not be applied in case of $`H^{}`$ and $`L^{}`$ given by (18) and (10) resp. as
$$\widehat{H}^{}=\frac{1}{2}V^2+\frac{1}{6}\left(\frac{1}{V}\widehat{p}^2+\widehat{p}\frac{1}{V}\widehat{p}+\widehat{p}^2\frac{1}{V}\right)+\mathrm{},$$
(44)
is a power series in expressions of type $`\frac{1}{V^m}\widehat{p}^l`$, $`\widehat{p}^l\frac{1}{V^m}`$, $`l,m=1,2,\mathrm{}`$ and $`\widehat{p}`$ and $`x`$ can not be separated. So a new quantization prescription is needed.
It is also not at all clear whether $`L^{}`$, given by (10), inserted into the exponent of the integral instead of $`L`$ in (43) yields the same physical results as using $`L`$ of (6). It seems rather that it leads to different value of the transition amplitude and to a different kind of quantization.
The question to be answered is: what are the limitations in using the Feynman rule for the โintegral over all pathsโ. Unfortunately, I do not feel to be able to give an answer to it. Thus the problem remains open, at least for me.
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# THE INTERACTION FORCE BETWEEN ROTATING BLACK HOLES AT EQUILIBRIUM 11footnote 1journal reference:Theor. Math. Phys. v.116, p.1024(1998)
## 1 Introduction
In the previous work , we studied the axially symmetric stationary vacuum solution of Einstein equations describing the equilibrium configuration of rotating black holes. The equilibrium configuration was understood to mean a stationary solution possessing a disconnected event horizon. All such solutions satisfy a certain boundary-value problem for a system of elliptic nonlinear equations that can be easily obtained from , where the regularity conditions for the symmetry axis and the event horizon were first formulated. Our investigation showed that this boundary-value problem reduces to a matrix Riemann problem with a rational conjugation matrix. We also noted that solutions corresponding to stationary equilibrium configuration of black holes are likely to develop a conical singularity on the symmetry axis, which is the reason why the black hole do not fall on each other. We note that for rotating black holes, the presence of conical singularity for arbitrary values of the solution parameters is not quite obvious from the physical standpoint, because it can be assumed that the rotation of black holes leads to repulsion forces that can compensate the gravitational attraction.
The conical singularity itself allows us to introduce the notion of the interaction force between the black holes. We consider this in more detail in the next section. The analytic expression for the interaction force between Schwarzschild black holes has long been known . The main result of the present work is our generalization of this to the case of rotating black holes.
## 2 The boundary-value problem
In cylindrical coordinates, the axially symmetric metric is
$$ds^2=Vdt^2+2Wdtd\varphi +Xd\varphi ^2+\frac{X}{\rho ^2}e^\beta (d\rho ^2+dz^2),$$
where the metric coefficient depend only on $`\rho `$ and $`z`$. The Einstein equations are then divided into two system of nonlinear equations, the first of which can be brought to the form
$$(\rho g_{,\rho }g^1)_{,\rho }+(\rho g_{,z}g^1)_{,z}=0,g=\left(\begin{array}{cc}V& W\\ W& X\end{array}\right),detg=\rho ^2,$$
$`(2.1)`$
and the second one to the form
$$(\mathrm{ln}(\frac{X}{\rho ^2}e^\beta ))_{,\zeta }=\frac{i}{2\rho }\left(1+\mathrm{tr}(\rho g_{,\zeta }g^1)^2\right),_\zeta =\frac{1}{2}(_zi_\rho ).$$
$`(2.2)`$
Let the event horizon have $`N`$ disconnected components, let $`z_1,\mathrm{}z_{2N}`$ be the $`z`$-coordinates of the intersection points of the event horizon and the symmetry axis, and let $`\mathrm{\Omega }_i`$ be the angular velocity of the $`i`$th black hole. In the neighborhood of the $`i`$th black hole, we choose the coordinate system
$$\rho ^2=(\lambda ^2m_i^2)(1\mu ^2),m_i=\frac{z_{2i}z_{2i1}}{2},$$
$$z\frac{z_{2i}+z_{2i1}}{2}=\lambda \mu ,|\mu |1,$$
in which the regularity condition of the symmetry axis and the event horizon can be written as
$$\left(\begin{array}{cc}1& \mathrm{\Omega }_i\\ 0& 1\end{array}\right)g\left(\begin{array}{cc}1& 0\\ \mathrm{\Omega }_i& 1\end{array}\right)=\left(\begin{array}{cc}(\lambda ^2m_i^2)\widehat{V}(\lambda ,\mu )& \rho ^2\widehat{W}(\lambda ,\mu )\\ \rho ^2\widehat{W}(\lambda ,\mu )& (1\mu ^2)\widehat{X}(\lambda ,\mu )\end{array}\right),$$
$`(2.3)`$
where $`\widehat{X}`$ and $`\widehat{V}`$ are smooth functions not equal to zero. Using (2.3), we can easily obtain the boundary conditions for the system (2.1)
$$\rho g_{,\rho }g^1=\left(\begin{array}{cc}0& O(1)\\ 0& 2\end{array}\right),\rho 0,z\mathrm{\Gamma },$$
$`(2.4a)`$
$$\widehat{\mathrm{\Omega }}_i\rho g_{,\rho }g^1\widehat{\mathrm{\Omega }}_i^1=\left(\begin{array}{cc}2& 0\\ O(1)& 0\end{array}\right),\rho 0zI_i,\widehat{\mathrm{\Omega }}_i=\left(\begin{array}{cc}1& \mathrm{\Omega }_i\\ 0& 1\end{array}\right).$$
$`(2.4b)`$
$$\rho g_{,z}g^1=O(1),\rho 0,zR.$$
$`(2.4c)`$
where $`\mathrm{\Gamma }`$ is the symmetry axis consisting of $`N+1`$ connected components,
$$\mathrm{\Gamma }=\mathrm{\Gamma }_j=RI_i(j=1,\mathrm{},N+1,i=1,\mathrm{},N),I_i=(z_{2i1},z_{2i}).$$
The symbol $`O(1)`$ denotes uniformly bounded functions on the corresponding interval. We impose the condition at infinity
$$W=\rho ^2O(\frac{1}{r^3}),X=\rho ^2(1+O(\frac{1}{r})),r=\sqrt{\rho ^2+z^2}.$$
$`(2.5)`$
Boundary conditions (2.4) and (2.5) are sufficient to construct solutions of the first group of Einstein equations (2.1).
We now consider the properties of solutions of system (2.2). Using (2.3), we can easily verify that for $`\rho =0`$ and $`z\mathrm{\Gamma }`$, we have
$$_z\beta =0,\beta |_{\mathrm{\Gamma }_i}=b_i,$$
$`(2.6)`$
where $`b_i`$ are some constants. The condition for the solution to be regular on the symmetry axis ($`\frac{X^{,a}X_{,a}}{4X}1`$ on the symmetry axis) is satisfied if and only if Eq.(2.3) is supplemented by the additional requirement that
$$b_i=0.$$
$`(2.7)`$
For asymptotically flat solutions (those for which $`\beta 0`$ as $`r\mathrm{}`$) condition (2.7) is automatically satisfied on the first and the last components of the symmetry axis but cannot be satisfied on the remaining components (see ). In the general case, the last statement is not yet proved. In this work, we evaluate the constant $`b_i`$, but the corresponding analytic expression is unfortunately too complicated to analyze the solvability of Eq.(2.7).
The $`b_i`$ parameters also have their own interpretation: the quantity
$$F_i=\frac{1}{4}(e^{b_i/2}1)$$
$`(2.8)`$
can be considered as the interaction force between the black holes. Here and what follows, we assume that $`b_1=b_{N+1}=0`$. We now comment on the origin of this interpretation. If condition (2.7) is not satisfied, the curvature tensor becomes singular in the neighborhood of the corresponding symmetry axis components. We assume that the singular points are filled with โmatterโ, which is precisely what prevents the black hole falling on each other. The tension in $`z`$ direction per unit surface is $`\mathrm{T}(e_z,e_z)`$ where $`\mathrm{T}`$ is the energy-momentum tensor and $`e_z`$ is the normalized vector parallel to $`/z`$. It is natural to define the interaction force of the black holes as the integral
$$_{S_\epsilon }\mathrm{T}(e_z,e_z)๐s,$$
where $`S_\epsilon `$ is the 2-surface coordinatized by $`\rho \varphi `$(with $`\rho \epsilon )`$ and $`ds`$ is the corresponding area element. Evaluating this integral, we obtain Eq. (2.8) .
To conclude this section, we give a brief derivation of representation for $`b_i`$ . We bring Eq.(2.2) to the form
$$_\rho \beta =\frac{\rho }{2}\left((_\rho \mathrm{ln}\frac{X}{\rho ^2})^2(_z\mathrm{ln}\frac{X}{\rho ^2})^2+\frac{(_\rho Y)^2(_zY)^2}{X^2}\right),$$
$`(2.9a)`$
$$_z\beta =_z\mathrm{ln}X(\rho _\rho \mathrm{ln}X2)+\frac{\rho _zY_\rho Y}{X^2},$$
$`(2.9b)`$
where
$$dY=\frac{1}{\rho }(XdWWdX),(d\rho =dz,dz=d\rho ).$$
Let $`C_\epsilon `$ be a curve connecting $`\mathrm{\Gamma }_{i+1}`$ and $`\mathrm{\Gamma }_i`$:
$$\rho =\epsilon \mathrm{sin}\tau ,z\frac{z_{2i}+z_{2i1}}{2}=\sqrt{\epsilon ^2+m_i^2}\mathrm{cos}\tau ,\mathrm{\hspace{0.33em}\hspace{0.33em}0}\tau \pi .$$
Then
$$b_{i+1}b_i=_{C_\epsilon }๐\beta .$$
$`(2.10)`$
The left-hand side of (2.10) is independent of $`\epsilon `$. Taking $`\epsilon `$ to zero and using (2.3) and (2.9), we obtain that
$$b_{i+1}b_i=2\left(\mathrm{ln}\widehat{X}(z_{2i})\mathrm{ln}\widehat{X}(z_{2i1})\right),$$
$`(2.11)`$
where $`\widehat{X}(z_{2i})=\widehat{X}(m_i,1)`$ and $`\widehat{X}(z_{2i1})=\widehat{X}(m_i,1)`$. In the next section, we use Eq.(2.11) to evaluate $`b_i`$.
## 3 The Riemann problem
We showed in that for every solution of boundary-value problem (2.4), (2.5) there exists a piecewise analytic matrix $`\chi (\omega )`$ satisfying the conjugation condition on the imaginary axis
$$\chi _{}(\omega )=\chi _+(\omega )\left(\begin{array}{cc}1& 0\\ 0& \omega \end{array}\right)T(k)\left(\begin{array}{cc}1& 0\\ 0& 1/\omega \end{array}\right),k=z+(\omega ^2\rho ^2)/2\omega $$
$`(3.1a)`$
and normalized at infinity by
$$\chi (\omega )I,\omega \mathrm{}.$$
$`(3.1b)`$
The rational matrix $`T(k)`$ is defined as
$$T(k)=\widehat{D}_{N+1}T_N\widehat{D}_N\mathrm{}T_1\widehat{D}_1,$$
$`(3.2)`$
where
$$T_j=\left(\begin{array}{cc}1\frac{2M_j}{kz_{2j1}}& 4M_j\mathrm{\Omega }_j\\ \frac{2L_j}{(kz_{2j})(kz_{2j1})}& 1+\frac{2M_j}{kz_{2j}}\end{array}\right),\widehat{D}_j=\left(\begin{array}{cc}1& D_j\\ 0& 1\end{array}\right).$$
$`(3.3)`$
The $`M_j`$ $`L_j`$ parameters have the physical interpretation of the respective full mass and angular momentum of $`j`$th black hole and are related by
$$M_j=m_j+2\mathrm{\Omega }_jL_j,\mathrm{\hspace{0.33em}\hspace{0.33em}2}m_j=z_{2j}z_{2j1}.$$
$`(3.4)`$
An additional requirement that the matrix $`T(k)`$ be symmetric leads to $`2N+1`$ nonlinear algebraic equations on the parameters $`D_j,L_j`$ and $`\mathrm{\Omega }_j`$, which determine $`D_j`$ and $`L_j`$ as functions of $`\mathrm{\Omega }_j`$ and $`z_j`$. In particular,
$$\underset{j=1}{\overset{N+1}{}}D_j+\underset{j=1}{\overset{N}{}}4\mathrm{\Omega }_jM_j=0.$$
$`(3.5)`$
Let $`d_1=D_1,d_{i+1}=d_i+D_{i+1}+4M_i\mathrm{\Omega }_i`$ and
$$t_{2j1}(k)=C_{2j1}^1\left(\begin{array}{cc}1& 0\\ \frac{1}{2\mathrm{\Omega }_j(kz_{2j1})}& 1\end{array}\right)C_{2j1},t_{2j}(k)=C_{2j}^1\left(\begin{array}{cc}1& \frac{\mathrm{\Omega }_j}{2(kz_{2j})}\\ 0& 1\end{array}\right)C_{2j},$$
where
$$C_{2j1}=\left(\begin{array}{cc}1& d_j\\ 0& 1\end{array}\right),C_{2j}=\left(\begin{array}{cc}0& \mathrm{\Omega }_j\\ 1/\mathrm{\Omega }_j& 4M_j\end{array}\right)C_{2j1}.$$
Then applying (3.5) and easily verified identity
$$T_j\left(\begin{array}{cc}1& d_j\\ 0& 1\end{array}\right)t_{2j1}^1t_{2j}^1=\left(\begin{array}{cc}1& d_j4M_j\mathrm{\Omega }_j\\ 0& 1\end{array}\right),$$
we derive that
$$T(k)=t_{2N}(k)t_{2N1}(k)\mathrm{}t_2(k)t_1(k).$$
$`(3.6)`$
We now note that the matrix
$$\left(\begin{array}{cc}1& 0\\ 0& \omega \end{array}\right)t_{2j}(k)t_{2j1}(k)\left(\begin{array}{cc}1& 0\\ 0& 1/\omega \end{array}\right)$$
has no singularity at $`\omega =0`$ and tends to identity matrix as $`\omega \mathrm{}`$,
$$\left(\begin{array}{cc}1& 0\\ 0& \omega \end{array}\right)t_{2j}(k)t_{2j1}(k)\left(\begin{array}{cc}1& 0\\ 0& 1/\omega \end{array}\right)|_{\omega =0}=\left(\begin{array}{cc}1& 8M_j(2M_j\mathrm{\Omega }_j+d_j)/\rho ^2\\ 0& 1\end{array}\right).$$
$`(3.7)`$
Therefore, the same properties are shared by the matrix
$$\left(\begin{array}{cc}1& 0\\ 0& \omega \end{array}\right)T(k)\left(\begin{array}{cc}1& 0\\ 0& 1/\omega \end{array}\right),$$
and Riemann problem (3.1) is correctly posed for any $`D_j,L_j,\mathrm{\Omega }_j`$ satisfying (3.5).
The only singularities of the conjugation matrix are simple poles at points
$$\omega _i^\pm =(z_iz)\pm \sqrt{(z_iz)^2+\rho ^2},\omega _i^{}=\rho ^2/\omega _i^+.$$
It now follows from Eq.(3.1) that $`\chi _\pm (\omega )`$ are rational functions with simple poles at the points $`\omega _i^\pm `$ ($`\omega _i^+>0`$, $`\omega _i^{}<0`$),
$$\chi _\pm (\omega )=I+\underset{j=1}{\overset{2N}{}}\frac{A_j^\pm }{\omega \omega _j^\pm },$$
$`(3.8)`$
with $`A_j^\pm `$ being independent of $`\omega `$. It follows from the unimodularity of $`T(k)`$ that $`det\chi _\pm (\omega )=1`$, and $`detA_j^\pm =0`$. Let $`\xi _j`$ be the eigenvector of $`A_j^{}`$; then
$$A_j^{}\xi _j=0,A_j^{}=a_j\xi _j^\sigma ,\xi _j^\sigma =(\xi _j^2,\xi _j^1),$$
$`(3.9)`$
with $`a_j`$ being a column vector. We further note that $`\chi _{}(\omega )`$ is unimodular if and only if there exists a constant $`v_i`$ such that
$$a_i=v_i\underset{\omega \omega _i^{}}{lim}\chi _{}(\omega )\xi _i,$$
$`(3.10)`$
Equation (3.9) and (3.10) yield a system of linear equations for $`a_i`$,
$$a_i=v_i\xi _i+v_i\underset{ik}{}\frac{\xi _k^\sigma \xi _i}{\omega _i^{}\omega _k^{}}a_k.$$
$`(3.10a)`$
The parameters $`v_i`$ and $`\xi _i`$ can be easily determined from conjugation matrix.
We define the matrix $`\mathrm{G}_i(k)=t_{i1}(k)\mathrm{}t_1(k)`$ and the vectors
$$c_{2j}=C_{2j}^1\left(\begin{array}{c}1\\ 0\end{array}\right),c_{2j1}=C_{2j1}^1\left(\begin{array}{c}0\\ 1\end{array}\right).$$
Then
$$\xi _i=\left(\begin{array}{cc}1& 0\\ 0& \omega _i^{}\end{array}\right)\mathrm{G}_i^1(z_i)c_i,$$
$`(3.11)`$
and also
$$v_i=\frac{\omega _i^{}v_i^0}{u_i\omega _i^{}v_i^0\xi _i^\sigma \left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)\xi _i},$$
$`(3.12)`$
where
$$v_{2j}^0=\frac{\mathrm{\Omega }_j}{\omega _{2j}^{}\omega _{2j}^+},v_{2j1}^0=\frac{1}{\mathrm{\Omega }_j}\frac{1}{\omega _{2j1}^{}\omega _{2j1}^+}$$
and
$$u_i=1+(\omega _i^{}\omega _i^+)v_i^0c_i^\sigma _k\mathrm{G}_i(z_i)\mathrm{G}_i^1(z_i)c_i.$$
We also note that
$$\xi _i^\sigma =\omega _i^{}c_i^\sigma \mathrm{G}_i(z_i)\left(\begin{array}{cc}1& 0\\ 0& 1/\omega _i^{}\end{array}\right)$$
The approach adopted here to solve the Riemann problem is similar to one used in . As is clear from what follows, to construct an exact solution to nonlinear system (2.1) we can take $`T`$ to be any symmetric matrix. This method for deriving exact solutions is evidently different from the method used in ; however, it does not lead to new solutions of (2.1). In particular, the solution that we investigate here is in the class of $`2N`$-solitons in Minkowski space background .
We now show how the solution of (2.1) can be reconstructed from the solution of Riemann problem (3.1). Let
$$D_1=_z\frac{2\omega ^2}{\omega ^2+\rho ^2}_\omega ,D_2=_\rho +\frac{2\omega \rho }{\omega ^2+\rho ^2}_\omega .$$
Because $`D_1k=D_2k=0`$, it follows from (3.1) that โlogarithmicโ derivatives $`D\mathrm{\Psi }\mathrm{\Psi }^1`$ of the functions
$$\mathrm{\Psi }_\pm (\omega )=\chi _\pm (\omega )\left(\begin{array}{cc}1& 0\\ 0& \omega \end{array}\right)$$
have no other singularities in addition to the simple poles at $`\pm i\rho `$. Therefore,
$$D_1\mathrm{\Psi }_{}(\omega )\mathrm{\Psi }_{}^1(\omega )=\frac{i\rho _\omega \mathrm{\Psi }_{}(i\rho )\mathrm{\Psi }_{}^1(i\rho )}{\omega i\rho }+\frac{i\rho _\omega \mathrm{\Psi }_{}(i\rho )\mathrm{\Psi }_{}^1(i\rho )}{\omega +i\rho }$$
$`(3.13a)`$
$$D_2\mathrm{\Psi }_{}(\omega )\mathrm{\Psi }_{}^1(\omega )=\frac{\rho _\omega \mathrm{\Psi }_{}(i\rho )\mathrm{\Psi }_{}^1(i\rho )}{\omega i\rho }+\frac{\rho _\omega \mathrm{\Psi }_{}(i\rho )\mathrm{\Psi }_{}^1(i\rho )}{\omega +i\rho }$$
$`(3.13b)`$
We now take into account that $`\overline{T}(\overline{\omega })=T(\omega )`$ and define the real matrix
$$g=\chi _{}(0)\left(\begin{array}{cc}1& 0\\ 0& \rho ^2\end{array}\right).$$
$`(3.14)`$
From (3.13) with $`\omega =0`$, one obtains
$$_\zeta gg^1=_\omega \mathrm{\Psi }_{}(i\rho )\mathrm{\Psi }_{}^1(i\rho ),\overline{}_\zeta gg^1=_\omega \mathrm{\Psi }_{}(i\rho )\mathrm{\Psi }_{}^1(i\rho ).$$
$`(3.15)`$
The compatibility condition for Eqs.(3.13) is given by system (2.1), and therefore, Eq.(3.14) is a solution of (2.1).
From here on, we assume that the parameters $`d_i`$ and $`M_i`$ are determined from the requirement that $`T(k)`$ be symmetric. Then
$$\underset{i=1}{\overset{N}{}}M_i(d_i+2\mathrm{\Omega }_iM_i)=0,(\underset{k\mathrm{}}{lim}k(T_{21}(k)T_{12}(k))=0)$$
$`(3.16)`$
and $`\chi _{}(0)=\chi _+(0)`$ (see (3.7)). Using the last identity and the uniqueness of the solution to the Riemann problem, we obtain the reduced problem
$$\chi _{}(\omega )=g\stackrel{~}{\chi }_+^1(\rho ^2/\omega )\left(\begin{array}{cc}1& 0\\ 0& 1/\rho ^2\end{array}\right)(\stackrel{~}{T}(\rho ^2/\omega )=T(\omega )),$$
$`(3.17)`$
where $`\stackrel{~}{T}`$ denotes the transposed matrix. Taking $`\omega `$ either to zero or to infinity in (3.17), we see that $`g`$ is a symmetric matrix.
We now investigate the properties of solutions of the Riemann problem as $`\rho 0`$. Let $`z\mathrm{\Gamma }_{m+1}`$, then
$$\omega _i^+2(z_iz),\omega _i^{}0,i2m+1,\rho 0$$
$`(3.18a)`$
$$\omega _i^+0,\omega _i^{}2(z_iz),i2m,\rho 0$$
$`(3.18b)`$
Using the formula $`\omega _i^+=\rho ^2/\omega _i^{}`$, we can easily verify that all the coefficients of linear system (3.10) also well defined for $`\rho =0`$, and it is therefore natural to assume the existence of the limit $`lim_{\rho 0}a_i`$. This assumption ensures the existence of all the limits used in what follows.
Introduce the matrix
$$\chi _1(\omega )=\chi _{}(\omega )\left(\begin{array}{cc}1& 0\\ 0& \omega \end{array}\right)\mathrm{G}_{2m+1}^1\left(\begin{array}{cc}1& 0\\ 0& 1/\omega \end{array}\right)=\chi _+(\omega )\left(\begin{array}{cc}1& 0\\ 0& \omega \end{array}\right)\mathrm{Q}_{2m+1}\left(\begin{array}{cc}1& 0\\ 0& 1/\omega \end{array}\right)$$
$$=I+\underset{i=1}{\overset{2m}{}}\frac{B_i^1}{\omega \omega _i^+}+\underset{i=2m+1}{\overset{2N}{}}\frac{B_i^1}{\omega \omega _i^{}},$$
$`(3.19)`$
where
$$\mathrm{Q}_i(k)=t_{2N}(k)\mathrm{}t_i(k),T(k)=\mathrm{Q}_i(k)\mathrm{G}_i(k).$$
$`(3.20)`$
We can see from (3.18) that
$$\underset{\rho 0}{lim}\chi _1(\omega )=I+\frac{B}{\omega }.$$
$`(3.21)`$
Further, let
$$\chi _2(\omega )=\chi _+(\omega )\left(\begin{array}{cc}1& 0\\ 0& \omega \end{array}\right)\stackrel{~}{\mathrm{G}}_{2m+1}\left(\begin{array}{cc}1& 0\\ 0& 1/\omega \end{array}\right)$$
In view of the identities
$$\left(\begin{array}{cc}1& 0\\ 0& \omega \end{array}\right)\stackrel{~}{t}_{2j1}\stackrel{~}{t}_{2j}\left(\begin{array}{cc}1& 0\\ 0& 1/\omega \end{array}\right)|_{\omega =0}=I,$$
$$\left(\begin{array}{cc}1& 0\\ 0& \omega \end{array}\right)\stackrel{~}{t}_{2j1}\stackrel{~}{t}_{2j}\left(\begin{array}{cc}1& 0\\ 0& 1/\omega \end{array}\right)|_{\omega =\mathrm{}}=\left(\begin{array}{cc}1& 8M_j(d_j+2M_j\mathrm{\Omega }_j)\\ 0& 1\end{array}\right),$$
we have
$$\chi _+(0)=\chi _2(0),\chi _2(\mathrm{})=\left(\begin{array}{cc}1& 8_{j=1}^mM_j(d_j+2M_j\mathrm{\Omega }_j)\\ 0& 1\end{array}\right).$$
$`(3.22)`$
In terms of $`\chi _1`$ and $`\chi _2`$, Riemann problem (3.1) can be rewritten as
$$\chi _2(\omega )=\chi _1(\omega )\left(\begin{array}{cc}1& 0\\ 0& \omega \end{array}\right)\mathrm{Q}_{2m+1}^1\stackrel{~}{\mathrm{G}}_{2m+1}\left(\begin{array}{cc}1& 0\\ 0& 1/\omega \end{array}\right),$$
$`(3.23)`$
and reduced problem (3.17) as
$$\chi _2(\omega )=g\stackrel{~}{\chi }_1^1(\rho ^2/\omega )\left(\begin{array}{cc}1& 0\\ 0& 1/\rho ^2\end{array}\right)=B^0+\underset{i=1}{\overset{2m}{}}\frac{B_i^2}{\omega \omega _i^{}}+\underset{i=2m+1}{\overset{2N}{}}\frac{B_i^2}{\omega \omega _i^+}.$$
$`(3.24)`$
Taking $`\rho 0`$ in (3.23) and taking into account (3.21), we see that
$$\chi _2(\omega )=B^0+\underset{i=1}{\overset{2N}{}}\frac{B_i^2}{\omega 2(z_iz)}=(I+\frac{B}{\omega })\left(\begin{array}{cc}1& 0\\ 0& \omega \end{array}\right)\gamma _{m+1}(\omega )\left(\begin{array}{cc}1& 0\\ 0& 1/\omega \end{array}\right),$$
$`(3.25)`$
where
$$\gamma _{m+1}(\omega )=\mathrm{Q}_{2m+1}^1(z+\omega /2)\stackrel{~}{\mathrm{G}}_{2m+1}(z+\omega /2).$$
$`(3.26)`$
The function $`\chi _2(\omega )`$ is nonsingular at $`\omega =0`$. This is possible only if
$$B=\left(\begin{array}{cc}0& \frac{\gamma _{m+1}^{12}(0)}{\gamma _{m+1}^{22}(0)}\\ 0& 0\end{array}\right)$$
Finally,
$$\chi _{}(0)=\left(\begin{array}{cc}1/\gamma _{m+1}^{22}(0)& \left(\gamma _{m+1}^{22}(0)_\omega \gamma _{m+1}^{12}(0)\gamma _{m+1}^{12}(0)_\omega \gamma _{m+1}^{22}(0)\right)/\gamma _{m+1}^{22}(0)\\ 0& \gamma _{m+1}^{22}(0)\end{array}\right)$$
$`(3.26)`$
Recall that $`\chi _{}(0)=\chi _+(0)=\chi _2(0)`$. We now note that $`\gamma _{m+1}(0)`$ is a symmetric real matrix with the unit determinant and also that $`\gamma _{m+1}^{22}(0)1`$ as $`z\mathrm{}`$; therefore,
$$\gamma _{m+1}^{22}(0)>0.$$
$`(3.27)`$
Since
$$\frac{\widehat{X}}{\lambda ^2m_i^2}|_{\mathrm{\Gamma }_{i+1}}=\gamma _{i+1}^{22}(0,z),\frac{\widehat{X}}{\lambda ^2m_i^2}|_{\mathrm{\Gamma }_i}=\gamma _i^{22}(0,z),$$
and $`\lambda ^2m_i^2(zz_{2i})(zz_{2i1})`$ as $`\rho 0`$, we find
$$\widehat{X}(z_{2i})=m_i\mathrm{res}_{z_{2i}}\gamma _{i+1}^{22}(0,z),\widehat{X}(z_{2i1})=m_i\mathrm{res}_{z_{2i1}}\gamma _i^{22}(0,z).$$
$`(3.28)`$
It follows from (3.27) that
$$\mathrm{res}_{z_{2i}}\gamma _{i+1}^{22}(0,z)>0,\mathrm{res}_{z_{2i1}}\gamma _i^{22}(0,z)<0.$$
Equations (2.11) and (3.28) give the sought-for analytic expression for the constants $`b_i`$ and the interaction force (2.8) of black holes.
## 4 Interaction force between two black holes
For two black holes, the symmetry axis has three connected components. We normalize $`\beta `$ such that $`b_1=b_3=0`$, then
$$F=\frac{1}{4}\left(\frac{\widehat{X}(z_2)}{\widehat{X}(z_1)}1\right),$$
$`(4.1)`$
and
$$\widehat{X}(z_2)=m_1\mathrm{res}_{z_2}\gamma _2^{22}(0,z),\widehat{X}(z_1)=m_1\mathrm{res}_{z_1}\gamma _1^{22}(0,z).$$
$`(4.2)`$
Expression (4.1) depends on $`d_i,\mathrm{\Omega }_i,L_i,M_i`$ and $`z_i`$. These parameters are not independent and must satisfy the system of nonlinear algebraic equations
$$\mathrm{res}_{z_i}T_{12}(k)=\mathrm{res}_{z_i}T_{21}(k),M_i=m_i+2\mathrm{\Omega }_iL_i.$$
$`(4.3)`$
It can be easily seen from (3.16) that
$$d_1=2M_1\mathrm{\Omega }_1+M_2d,d_2=2M_2\mathrm{\Omega }_2M_1d.$$
$`(4.4)`$
Using (4.3) and (4.4), we can eliminate parameters $`d_1,d_2,\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$ from the expression for the interaction force (4.1). Unfortunately, the formula thus obtained is extremely cumbersome. The situation is greatly simplified in one particular case, however. Namely, let the black holes rotate in opposite directions and have the same total and irreducible masses:
$$M_1=M_2=M,m_1=m_2=m,\mathrm{\Omega }_1=\mathrm{\Omega }_2=\mathrm{\Omega }.$$
We set
$$z_1=m_1,z_2=m_1,z_3=Rm_2,z_4=R+m_2.$$
$`(4.5)`$
Then
$$L=2M^2(m+M)\frac{R+2M}{R2M}\mathrm{\Omega },M=m+2\mathrm{\Omega }L,L_1=L_2=L,d=0$$
$`(4.6)`$
and
$$F=\frac{M^2}{R^24M^2}.$$
$`(4.7)`$
By definition, $`R2m`$. Using (4.6), we can easily show that $`2MR`$ (where the equality takes place only in the limiting case of $`R=2m`$); therefore , $`F>0`$ and tends to infinity as $`R2m`$. We also observe that the interaction force in the case under consideration is exactly the same as for two Schwarzschild black holes with identical masses (equal to $`M`$). The interaction force of two nonrotating black holes ($`\mathrm{\Omega }_1=\mathrm{\Omega }_2=0`$) is given by
$$F=\frac{m_1m_2}{R^2(m_1+m_2)^2}.$$
$`(4.8)`$
It is easy to see that in the cases we have considered, $`F`$ tends to the Newtonian limit as $`R\mathrm{}`$. This result can be generalized; the estimate
$$F=\frac{M_1M_2}{R^2}+O(\frac{1}{R^4})$$
$`(4.9)`$
holds as $`R\mathrm{}`$. To derive (4.9), it suffices to use the approximate solution of (4.3):
$$d=\frac{2L_2M_1+2L_1M_2}{M_1M_2}\frac{1}{R^2}+4\left(L_1+L_2+\frac{2L_1M_2^2+2L_2M_1^2}{M_1M_2}\right)\frac{1}{R^3}+O(\frac{1}{R^4}),$$
$$\mathrm{\Omega }_i=\frac{L_ia_i}{2M_i^2(M_i+m_i)}$$
and
$$a_1=1\frac{4M_2}{R}+\frac{8M_2^2}{R^2}+O(\frac{1}{R^3}),a_2=1\frac{4M_1}{R}+\frac{8M_1^2}{R^2}+O(\frac{1}{R^3}).$$
We now discuss the general properties of constraints (4.3) on the angular velocities and angular momenta. We especially note that the equality $`L_2=\mathrm{\Omega }_2=0`$ cannot be satisfied for any values of the remaining parameters with $`L_10`$. This means that if $`L_2=0`$ and the second black hole does not contribute to the total angular momentum, an approaching observer would still see this black hole rotating with a certain angular velocity. We consider this behavior quite natural. Indeed, the notion of the angular velocity is related to the behavior of geodesics in the neighborhood of black hole: every test particle approaching the black hole is involved in its rotation, and the angular velocity of this rotation becomes equal to the angular velocity of the black hole when the particle reaches the event horizon. That the angular velocity is nonzero is therefore determined by the fact that the first black hole involves the second one in its rotation. This explanation corresponds to the behavior of $`\mathrm{\Omega }_2`$ at large distances: if $`L_2=0`$ then $`\mathrm{\Omega }_2=O(1/R^3)`$.
We return to the exact solution (4.6). The foregoing discussion allows us to take the angular momenta of each black hole as the independent parameters. We then have $`\mathrm{\Omega }0`$ and $`Mm`$ as $`R2m`$ which means that as black holes approach each other, they loose the full mass and โslow downโ each others rotation.
Evidently, the realistic solution for the configuration of two black holes cannot be stationary. As we have seen however, the stationary solutions admits a sufficiently realistic physical interpretation and demonstrate, at least qualitatively, the effects caused by the interaction of black holes. In our opinion, this is because the Einstein equations determine not only the gravitational field but also the law of motion of material bodies .
This work was supported by RFBR grant No. 96-01-00548.
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# Imaging of Ultraluminous Infrared Galaxies in the Near-UV
## 1 Introduction
Many recent Hubble Space Telescope (HST) observations have shown the prevalence of clustered star-formation in interacting galaxies (Whitmore et al. 1995, Meurer et al. 1995). Other observations have shown that Ultraluminous Infrared Galaxies (ULIGs; objects with infrared luminosities, $`L_{\mathrm{ir}}`$ <sup>1</sup><sup>1</sup>1$`L_{\mathrm{ir}}`$ L(8โ1000ยตm) is computed using the flux in all four IRAS bands according to the prescription given in Perault (1987); see also Sanders & Mirabel (1996). Throughout this paper we use $`H_\mathrm{o}`$ = 75 km s<sup>-1</sup>Mpc<sup>-1</sup>, q<sub>o</sub> = 0.5 (unless otherwise noted)., greater than $`10^{12}L_{\mathrm{}}`$) often have extended morphologies consistent with being advanced mergers (Sanders et al. 1988a, Murphy et al. 1996, Clements et al. 1996). Surace et al. (1998, hereafter Paper I) showed that these same systems have evidence for โknotsโ of star-formation with characteristic diameters of $``$100 pc distributed over their central several kpc and along their extended tidal features. Estimated ages for these knots based on optical and near- infrared colors were in the range of 10โ800 Myrs. In at least 25% of ULIGs, particularly those with โwarmโ mid-infrared colors and Seyfert optical spectra, it is also possible that some fraction of the optical and hence U light is emitted by an AGN.
With the exception of integrated photometry taken with single beam photometers for a few systems (Neugebauer et al. 1987, Young et al. 1996), only two of the ULIGs have been examined at short wavelengths before (Trentham et al. 1999). This is primarily due to former assumptions about the dusty nature of the ULIGs. Since the bulk of the luminosity in ULIGs is emitted at 25โ100 $`\mu `$m as a result of dust absorption and reradiation, the assumption has always been that the galaxies themselves are intrinsically dusty and that extinction must therefore be very high. In short, there would be nothing to see at short wavelengths. However, ground- based and HST imaging (Paper I, Surace 1998, Surace & Sanders 1999; hereafter Paper II) has revealed the presence of luminous clustered star formation at optical and near-infrared wavelengths. Paper II developed a broad-band multi-color technique at optical and near-infrared wavelengths in order to measure and correct for the presence of a foreground dust screen. These observations have shown that the apparent line-of-sight extinction is quite low (typically A$`{}_{\mathrm{V}}{}^{}<`$ 2 magnitudes), implying that many of the active star forming regions have already ejected any surrounding dust or have particularly favorable scattering geometries. Additionally, millimeter-wave observations of some of these systems have indicated directly that the dust is strongly centrally concentrated and localized, and hence the more extended starburst regions may be relatively unobscured (Bryant & Scoville 1996, Scoville et al. 1997). Line-of-sight extinction is therefore not a serious impediment to detecting the ULIG star-forming regions at short wavelengths.
There are several advantages to working at shorter wavelengths. As noted in Paper I, there is an ambiguity in interpreting broad-band colors of starbursts with burst ages of 10โ100 Myrs, as their colors in all filters redder than B are degenerate during this time period. This can be partially disentangled through the use of U-band imaging; (U$``$B) colors evolve almost monotonically throughout this period in the life of the starburst. Although reddening will still prevent a precise starburst age estimate since it is most severe at short wavelengths, the (U$``$B) colors, in conjunction with longer wavelength data ( Surace et al. 1998 \[Paper I\], Surace et al. 1999 \[Paper II\]), will help constrain the upper age limits of the starbursts more accurately. Furthermore, the youngest stars, which dominate the luminosity of the starburst, have spectral energy distributions (SEDs) which peak at ultraviolet wavelengths. Ultraviolet observations therefore provide the greatest possible contrast between compact young star- forming regions and the redder, extended underlying galaxy. In the evolutionary model of Sanders et al. (1988a), warm ULIGs are a more advanced evolutionary state than cool ULIGs, and this should be reflected in the upper limits of the knot ages.
Finally, there have been many recent advances in the study of high redshift galaxies. Particularly important are deep submillimeter surveys, which have discovered a significant population of high-redshift galaxies which are luminous at far-infrared wavelengths (Smail et al. 1997, Hughes et al. 1998), and which are inferred on the basis of their submillimeter luminosity to be high-redshift counterparts of the local ULIGs. Optical and near-infrared observations of these high-redshift objects correspond to rest-frame UV, and hence a better understanding of the UV properties of local ULIGs is required in order to interpret high-redshift results.
We present U imaging data for selected ULIGs drawn from two complete samples of โwarmโ AGN-like ULIGs and โcoolโ starburst-like ULIGs which have been previously well-studied at other wavelengths (Sanders et al. 1988ab, Kim et al. 1998, Paper I, II, Surace 1998), with the primary goal of characterizing these systems at near-UV wavelengths in order to further understanding of high-redshift systems. U and optical colors are used to derive starburst knot ages assuming synthetic starburst models. The estimated ages, reddenings, and luminosities are used to derive the contribution of the extended starburst regions to the known bolometric luminosities of the galaxy systems. The history of star-formation in these systems and the implications of these new observations for high-redshift systems are discussed.
## 2 Data
### 2.1 Sample Selection
The ultraluminous infrared galaxies were chosen from two small, nearly complete samples. The first sample consists of the โwarmโ ULIGs first described by Sanders et al. (1988b); more detailed optical and near-infrared imaging can be found in Surace et al. (1998) and Surace & Sanders (1999). These systems were selected using a mid-infrared color criterion known to select active galaxies ($`f_{25}/f_{60}>0.2`$) <sup>2</sup><sup>2</sup>2The quantities $`f_{12}`$, $`f_{25}`$, $`f_{60}`$, and $`f_{100}`$ represent the IRAS flux densities in Jy at 12 ยตm, 25 ยตm, 60 ยตm, and 100 ยตm respectively. and which has been used previously to select AGN from the IRAS Point Source Catalog (de Grijp et al. 1985, 1987). Studies of several small but complete samples of warm ULIGs have shown that many of these objects have a point-like optical appearance on the Palomar Sky Survey and that they exhibit broad (i.e. Seyfert 1) optical emission lines, characteristics that have led them to be referred to as โinfrared QSOsโ (e.g. Low et al. 1988, Sanders et al. 1988b). These warm objects, which account for $``$ 20โ25 % of the total population of ULIGs discovered by IRAS (Kim 1995), appear to represent a critical transition stage in the evolution of the larger population of โcoolโ ULIGs into optical QSOs according to the unification scheme proposed by Sanders et al. (1988a). Since they were selected to be transition objects similar to optical QSOs, the luminosity criterion used was L$`{}_{\mathrm{bol}}{}^{}>10^{12}`$ L , hence the inclusion of two systems with L$`{}_{\mathrm{ir}}{}^{}<10^{12}`$ L(Sanders et al. 1988b).
The second sample consists of ULIGs with complementary โcoolโ colors ($`f_{25}/f_{60}<0.2`$) as detailed by Surace et al. (2000). These systems are more similar to the majority of ULIGs (representing roughly 75% of the ULIG population; Kim 1995) and are thus similar to the Bright Galaxy Sample ULIGs described by Sanders et al. (1988a). They typically have HII or LINER spectra, have diffuse nuclear regions with optical/near-infrared colors similar to young stars obscured by A$`{}_{\mathrm{V}}{}^{}`$1-4 magnitudes (Surace 1998), and a high percentage of double nuclei. They therefore appear to represent an earlier stage in the merger process. Together these two samples span the entire range of interactions and far-infrared colors observed in ULIGs and in the evolutionary scheme of Sanders et al. (1988ab) represent two temporally distinct populations of the same type of physical object. Since the ultraviolet emission originates predominantly in massive young stars which have formed on timescales similar to that of the merger process, if the evolutionary hypothesis is correct then an examination of the warm and cool ULIGs should provide evidence for differences in their ultraviolet properties.
Because of spatial constraints and the limited telescope resources available, not all of the objects in the two complete samples could be observed. A total of 20 objects were observed: 11 (of a complete sample of 12) โwarmโ ULIGs and 9 (of a sample of 14) cool ULIGs. Tables 1 & 2 list the objects actually observed in the two samples, along with observation details.
### 2.2 Observations
The data were taken on December 26โ29, 1997, March 23โ24, 1998 and November 12โ13, 1998 at the f/31 focus of the UH 2.2m telescope on Mauna Kea using the Orbit 2048$`\times `$2048 camera. The Orbit camera is a backside- illuminated, thinned CCD that is UV-flooded for normal operation; itโs quantum efficiency at U (based on laboratory tests) is as high as 80%. At the f/31 focus this camera is oversampled, and hence the data was binned 2$`\times `$2 on chip yielding a pixel scale of 0.09โณ pixel<sup>-1</sup>. The UH U filter (which is noticeably bluer than the standard Johnson U filter; $`\lambda _{central}=`$3410ร
, $`\mathrm{\Delta }_\lambda `$=320ร
;Wainscoat 1996) was used for all observations. Total exposure times were typically 30โ90 minutes, the long exposure times being necessitated by the narrowness of the U filter. A sequence of five or more frames was taken in order to allow post-processing rejection of the high cosmic ray incidence. These frames were dithered to further decorrelate structured pixel-to-pixel noise in the CCD itself (such as blocked rows, flat- fielding errors, etc.). The CCD bias pattern was removed by subtracting a median bias frame constructed from 20โ30 bias frames (zero-time exposures), and pixel-to-pixel response variations were corrected by dividing the data by a high S/N median twilight flat. In both cases the expected poisson noise in the calibration frames was below 0.5%. On the second run this camera developed severe dark current instabilities. This was calibrated by taking 4โ5 hours of dark current (closed shutter integration) data at the beginning and end of each night. The high S/N dark current images were made for each half of the night, scaled to the exposure times used for the science data, and subtracted from them. In most cases these dark integrations had the same exposure times as used for the actual science observations, thus helping eliminate time-dependent dark current variations. This strategy was effective with no visible dark current pattern after processing. The individual frames were then rotated to the normal orientation (northeast at upper left), aligned using IMALIGN in IRAF, and medianed using IMCOMBINE and an algorithm that rejects cosmic rays based on the CCD noise characteristics.
The data were flux calibrated by observing optical standard stars (Landolt 1983, 1992) with a wide range of colors and at different airmasses, allowing the derivation of accurate zero-point offsets and color terms relating the U filter to the standard Johnson U-band filter. This also allowed the atmospheric extinction to be measured and corrected; this can be as high as 0.2โ0.6 magnitudes per airmass (the observed mean was 0.27) at such short wavelengths. In general, the photometric calibration appears to be accurate to 0.04 magnitudes. Possibly the largest uncertainty is the conversion from U to Johnson U-band; the relatively large color term can introduce uncertainties as high as 0.05โ0.1 magnitudes, depending on the underlying SED of the observed feature. Throughout this paper quoted magnitudes will be on the Johnson U-band calibration, using the color correction :
m<sub>U</sub> = -2.5 Log$`\left(\frac{adu}{exptime}\right)`$ \+ ZP$`_U^{}`$ \- (airmass$``$extinction$`_U^{}`$) + 0.13$``$(B-I)
(1)
where the first term is the instrumental magnitude, the second is the instrumental magnitude zeropoint, the third is the atmospheric extinction, and the last term is the color dependency which is appropriate for the combination of the Orbit camera and the UH U filter. No K-corrections have been applied to any magnitude reported here since the unknown SED at U prevents an accurate determination of such corrections. (B$``$I) is derived from previous optical observations (Surace et al. 1998, Surace 1998) and is accurately known (with values of (B$``$I) typically just greater than unity). For the integrated galaxy magnitudes we used the integrated galaxy colors for the color term correction, while for the star-forming regions we used the local colors.
Fluxes were measured using aperture photometry. Circular and polygonal apertures were constructed using SKYVIEW. The background was estimated using a median of the data directly outside the apertures, excluding any bright small-scale structure. Integrated magnitudes were derived from large aperture measurements. In all cases the galaxies have been deeply imaged at I-band (Surace 1998). The apertures were chosen to be large enough to encompass the entire known extent of the galaxies. The knots were identified based on spatial coincidence with previously identified star-forming knots seen in optical and near-infrared images. In many cases the spatial resolution is too poor to resolve all of the known knots. In these cases, apertures were chosen to be large enough to encompass the entire extent of each group of unresolved knots (and whose combined photometry are listed with a โ/โ in Table 3). Aperture corrections were derived from the observed point spread function (PSF) and have been applied to the magnitudes reported here.
Calibration of the PSF is problematic at U because most stars are not very luminous at such short wavelengths, thereby decreasing the likelihood of achieving a PSF S/N comparable to that of the target objects. In some cases it was possible to estimate the PSF directly from the combined science images in the same manner as in Paper II. DAOPHOT was used to iteratively construct an empirical PSF based on all the observed stars in the combined data, resulting in a high S/N PSF. In the remaining cases where not enough stars were present the PSF was derived from the standard star observations taken immediately before and after each set of science data. Most of the image degradation at short wavelengths is due to high-order atmospheric turbulence, and hence the quality of the off-axis tip/tilt guiding is largely irrelevant to the achieved PSF; the observed flux standards should therefore provide adequate PSF calibrators. Unfortunately, because of the poor resolution at short wavelengths, aperture corrections tend to be quite large (0.5โ1 mag) for sub-arcsecond apertures, and hence knowledge of the PSF is critical. In most cases, in order to derive accurate colors the previous B and I- band imaging had to be re-evaluated in apertures larger than those previously used in Surace (1998) in order to compensate for seeing-induced aperture effects. This was particularly necessary for comparison with the WFPC2 data. Total cumulative uncertainties, including those in calibration, measurement, and aperture corrections, are $``$0.2 magnitudes.
## 3 Results
### 3.1 Morphology
Figures 1 & 2 present the U observations of the 11 warm and 9 cool ULIGs, respectively, and Appendix A provides detailed descriptions. Many of the systems show considerable structure in their extended tidal features. In general, the morphology of the ULIGs at U is very similar to their morphology at B and to a lesser extent at I-band (Paper I). This is not surprising, as the extinction at U (A<sub>U</sub>=1.53A<sub>V</sub>) is only slightly larger than that at B (A<sub>B</sub>=1.45 A<sub>V</sub>), yet is much larger than I (A<sub>I</sub>=0.48 A<sub>V</sub>; Rieke & Lebofsky 1985). Therefore, dust extinction effects are unlikely to significantly alter the relative morphology between U and B. Any features near the I-band detection limit, however, will be completely obscured at U by even just a few magnitudes of extinction. Because of the generally poor spatial resolution at U (0.7โ1.5โณ), any features within a radius of $``$1.5-2 kpc of the luminous point-like nuclei of some of the warm ULIGs will be undetectable.
Of the โwarmโ ULIGs, the two double nucleus systems Mrk 463 (3.8 kpc nuclear separation) and IRAS 08572+3915 (6.2 kpc nuclear separation) have the most spectacular U features. In both cases up to half a dozen U knots are scattered along the lengths of the tidal tails. The extra-nuclear (diameter greater than 2.5 kpc) regions account for 65% and 83% of the total U flux, respectively, where this number includes both the compact knots and the extended tidal emission. Similarly, the single active nucleus galaxies I Zw 1, Mrk 1014 and Mrk 231 have U knots embedded in apparent tidal features. IRAS 12071$``$0444 and IRAS 15206+3342 have evidence for marginally resolved knots in their inner 10 kpc diameter nuclear regions. IRAS 05189$``$2524, Pks 1345+12, and IRAS 01003$``$2238 have no detectable knots, although this may be due to our limited spatial resolution. The one warm ULIG which shows no U knots at all and should, based on the U -band spatial resolution and detection limits, is IRAS 07598+6508. In 7 out of 11 cases more than 50% of the U flux originates within a radius of 4 kpc of the nucleus, which more than anything reflects that 4 of the sample are acknowledged to be QSOs, with UV-luminous active nuclei.
The โcoolโ ULIGs have an equally varied U morphology. More than half of the cool ULIGs have more than 50% of their U -band flux emitted in their inner 4 kpc radii. The single nucleus systems IRAS 00091$``$0738, UGC 5101, Mrk 273, and IRAS 23365+3604 have long, well-developed tidal tails which are easily at I-band (Surace et al. 2000), yet none of these tails has the high surface brightness U knots found in the warm systems IRAS 08572+3915 and Mrk 463 and even the extended tidal structure itself is barely detected at U . The double nucleus systems IRAS 12112+0305 (2.5 kpc nuclear separation) and IRAS 14348$``$1447 (5.5 kpc nuclear separation), however, clearly have U knots located in their tails and also in their nuclear regions. This indicates that the appearance of knots in the tidal structure is more a function of dynamical age as identified by the appearance of double nuclei than it is of mid-IR color. Finally, the very widest separation system examined, IRAS 01199$``$2307 (24 kpc), has no evidence for U knots in its barely detectable tails.
This dichotomy in the appearance or lack of knots along the extended tidal features may indicate differences in the star-formation rate in the tails during different phases of the merger evolution. Since every merger system must pass through a double-nucleus phase prior to when the two nuclei have coalesced into a single nucleus, the double nucleus systems are dynamically younger than the single-nucleus systems. However, knots embedded in the extended tidal tail structure are visible in 5 out 7 of the double nucleus systems, but only 2 out of 13 of the single nucleus systems (Table 4). Based on Monte Carlo simulations, if the probability of having extranuclear knots were the same for both samples, the likelihood of our finding as few single nucleus systems and as many double nucleus systems with knots is less than 1%, regardless of the lifetime of the knots vs. the merger lifetime. This appearance of compact star-forming knots in the tails of every merger system with nuclear separations less than 10 kpc, but not in single- nucleus post-merger systems, may indicate that the tails undergo a relatively unified burst of compact, clustered star formation during the 100 Myr period prior to the merging of the nuclei, and that no further local bursts occur. Conversely, many of the single nucleus systems (e.g. UGC 5101, IRAS 23365+3604) do have complex structure in their nuclear regions. Speculatively, this could be due to gas depletion in the tails; the denser nuclear environment may be able to fuel the bursts for much longer, or to fuel a sequence of many bursts as was indicated in Paper I. Fading of the starburst knots is very rapid at U ; the knots fade from their peak intensity (which occurs at an age of 30 Myrs for an instantaneous burst โ see below) by nearly 5 magnitudes by the time they are only 300 Myrs old. Fading by only 3 magnitudes would make most of the luminous knots observed undetectable. Since we detect no knots even more luminous than those seen here, it is therefore likely that we are detecting the knots near the peak of their luminosity and that most of the knots we observe at U are between 10-100 Myrs in age, which is borne out by an analysis of their (U$``$B,B$``$I) colors (see below). Since the nuclear regions of the more dynamically advanced systems (i.e. Mrk 273, UGC 5101) continue to have luminous U knots, either these knots have an intrinsically higher peak luminosity (thus allowing them to be detectable for a longer period of time), the time-scale of the burst is much longer than the instantaneous burst modeled here, or there is an on-going generation of new knots. Note that these knots are shown below to be physically unrelated to the source of the majority of the bolometric luminosity and hence this result does not indicate anything about the evolution of the primary luminosity source in ULIGs.
An alternative explanation to the dichotomy in the appearance of star-forming knots is that two entirely different populations of ULIGs exist. The above explanation implies that the ultraluminous systems remain ultraluminous throughout the merger process. Because we have selected our targets based on a luminosity criterion, it is possible that we have actually selected systems that are ultraluminous during the early merger phase and are accompanied by many star-forming extranuclear knots, and a different population that is ultraluminous during the late merger phase and which has no knots. One way to test this would be to search for relic knots in the tails of the single nucleus systems. If knots with colors similar to older stars were found in the single nucleus systems then this would demonstrate that during their earlier double-nucleus phase they appeared similar at U to the double-nucleus systems we are observing. However, we lack the sensitivity to detect these older knots.
Figures 3 & 4 present surface brightness profiles for the warm and cool ULIGs, respectively, at U . The profiles were centered on the U flux centroid corresponding to the major โnucleiโ (the brightest feature at B-band) identified previously in Paper I and Surace (1998). Circular azimuthally- averaged isophotes were fitted using the IRAF/STSDAS task ELLIPSE. In each case the centers and ellipticities were held fixed. The radial profiles were characterized in two ways.
First, the radial profiles were fit with a Sรฉrsic (1968) surface brightness law:
$$I(r)=I_eexp(b[(r/R_e)^{1/n}1])+I_{back};b=1.9992n0.3271$$
(2)
which for n=4 is a de Vaucouleurs profile, and for n=1 is an exponential disk profile. With this definition of b, R<sub>e</sub> and I<sub>e</sub> are the half-light radius and intensity. The fitting was done with the IGOR data analysis package. In order to minimize the effects of seeing, which dominate the shape of the profile at small radii, the profiles were fit only at radii greater than the FWHM of the observed PSF (typically 1โณ). The Sรฉrsic profile has many degrees of freedom, and so additional constraints had to applied. I<sub>back</sub> was set equal to the observed background sky value far from the galaxy. I<sub>e</sub> was constrained to be greater than zero and less than the maximum observed surface brightness. Finally, R<sub>e</sub> was estimated by integrating the total light in successive radii until the half-light radius was found. Initially, the fitting algorithm allowed R<sub>e</sub> to vary as a free parameter; however, if it arrived at a clearly non-physical value (i.e. less than 100 pc or greater than 20 kpc), then R<sub>e</sub> was held fixed at the measured value. Even so, the derived value of n was very sensitive to the input parameters. This is primarily due to the distance of the ULIGs and the limited spatial resolution of the U observations. Zheng et al. (1998) found, based on high spatial resolution HST data at I-band, that typical values of R<sub>e</sub> for a small sample of ULIGs were 1โ6 kpc, with a median near 3 kpc. Unfortunately, this is similar to the radii over which the U radial profiles have been fit (typically 2โ10 kpc). The Sรฉrsic profile has an inflection near R<sub>e</sub>, and as a result n is poorly constrained (Figure 5). Furthermore, profiles with fixed parameters other than n are very similar at large radii for n $`>2`$. It is the profile at very small radii, which we cannot measure easily from the ground, that constrains the fit best. Only systems with apparently single nuclei were fit. This was because the uncertainties in estimating parameters in the presence of a second nucleus resulted in an unstable, and probably meaningless, value of n. The observed value of n is certain to within $`\pm `$0.3. Previous work has shown that n is usually slightly underestimated (Marleau & Simard 1998).
The Sรฉrsic index n was found to lie between 0.8 (IRAS 01003-2238) and 3.5 (IRAS 05189-2524), with a median value of 1.8. There is nothing to suggest a trend between the index n and far-infrared color. Both warm and cool galaxies have nearly the same range and median; Kolomogoroff-Smirnov statistics indicate that we cannot reject the null hypothesis that the two samples were drawn from the same parent sample (Press et al. 1992). The value of n$`=`$1.8 is intermediate between an exponential disk and a de Vaucouleurs profile with the same half-light radius, although it will more closely resemble the de Vaucouleurs profile than the exponential disk. The derived value of R<sub>e</sub> is contaminated by residual light from the QSO nucleus in I Zw 1, Mrk 1014, IRAS 07598+6508, and Mrk 231. In the remaining galaxies R<sub>e</sub> ranges from 1.9 kpc (IRAS 00091-0738) to 9.3 kpc (Mrk 273), with a median of 5.3 kpc. This is larger than that found by Zheng et al. (1998).
Additionally, for each galaxy, an r<sup>-1/4</sup> law profile was fit to the outer isophotes, beginning at a radius equal to the FWHM of the PSF. An exponential disk model was also fit. Similar to the result found above with the Sรฉrsic law, it is apparent that these two models fit the ULIGs to a varying degree (see Figures 3 & 4). Those dominated by a bright central nucleus (I Zw 1, Mrk 1014, IRAS 07592+6508, and Mrk 231) are all extremely well fitted by a de Vaucouleurs (r<sup>-1/4</sup> ) profile, while an exponential disk fits poorly. In some cases, both the exponential disk and the de Vaucouleurs profile fit the outer isophotes equally well (e.g. IRAS 15206+3342, PKS 1345+12), while in other cases, the outer radii are well-fit by the r<sup>-1/4</sup> profile and the inner regions by the exponential disk. Since the inner regions are more affected by dust extinction and resolution effects it is unlikely that this really represents the underlying stellar distribution. There is no obvious trend between far-infrared colors (i.e., โwarmnessโ or โcoolnessโ) and the type of best-fit model; it is not the case that all warm ULIGs have r<sup>-1/4</sup> profiles, while the cool have exponential disks, which might have been expected if the two samples had different dynamical ages.
Since the morphologies of the systems are consistent with the remnants of major mergers, this is the expected result since relaxation of the disks results in their assuming a profile approximated by the r<sup>-1/4</sup> law. Significant deviations from an r<sup>-1/4</sup> law are still seenโ specifically, they decrease more rapidly at large radii than might be expected for a de Vaucouleurs profile. Given the relatively sharp, discontinuous edges of the tidal features, which are the most extended parts of the merger systems, this turnover in the brightness profile at large radii probably represents itโs truncation past the maximum extent of the tidal debris. This falloff in intensity is also similar to the behavior expected for an exponential disk, and may be due to incomplete relaxation of the merger remnant (Toomre 1977).
Similar results have been found for merging galaxies in general (Schweizer 1982, Wright et al. 1990, Stanford & Bushouse 1991), and more specifically for ULIGs at R & K- band (Wright et al. 1991, Kim 1995). The latter is not surprising, since it is now known that most ULIGs are merging systems. What is perhaps surprising is that this result is true even at U and over a large range of radii extending to as far as 20 kpc. While longer wavelength observations (e.g. near-infrared) trace primarily the older stellar population and hence most of the stellar mass, the U observations should trace higher mass stars and sites of intense local star formation. However, an examination of Table 3 indicates that most of the compact features observed at U represent only a fraction of the total U flux and hence the U low surface brightness emission continues to trace the general galactic stellar population. This is similar to the result of Meurer et al. (1995). U traces stars with ages of at most a few hundred Myrs, which is similar to the dynamical merger timescale. The existence of an extended smooth stellar population at U may indicate that the ULIGs are currently undergoing spheroidal component growth, similar to that observed by Fanelli et al. (1997) and Smith et al. (1996) in the far-ultraviolet in two nearby starburst galaxies.
### 3.2 Colors and Derived Starburst Ages
#### 3.2.1 Modeled Colors
Table 3 presents the photometric data for the U observations, all of which have been transformed to the Johnson U-band from the observed U filter. The colors of the observed features were compared to various modeled colors. The U colors provide essentially no discriminant between AGN and starburst activity. The synthetic optical QSO model introduced in Paper II predicts (U$``$B, B$``$I) colors for a QSO at the sample median redshift of ($``$0.78, 0.67) and the UVSX sample of Elvis et al. (1994) has colors of (U$``$B),(B$``$I) = ($``$0.82,1.12)$`\pm `$(0.16,0.19), which are similar to a modeled instantaneous starburst with an age of 10 Myrs. As discussed in Paper II, long wavelength emission ($`>`$ 2 $`\mu `$m) therefore remains the best possible broad-band color discriminant between AGN and recent star formation activity. Throughout this analysis we will continue to consider the U emission as a result of either starburst activity or AGN depending on the results found at other wavelengths (Surace 1998). Larson & Tinsley (1978) find typical (U$``$B) colors for non- interacting galaxies in the RC2 in the range of $``$0.2 to $`+`$0.6, and Guiderdoni et al. (1988) find typical values of (U$``$B) for spirals of 0.03. Such a wide span of colors is probably due to the variable star formation histories of the โnormalโ galaxies involved. Perhaps a more meaningful comparison can be made to the (U$``$B) colors of individual stellar spectral types: (U$``$B) = $``$1.15, $``$1.06, $``$0.02, 0.07, 0.05, 0.47, and 1.28 for spectral classes O5, B0, A0, F0, G0, K0, and M0, respectively (for stars on the main sequence โ the corresponding giant stars are somewhat redder; Allen 1973). Only OB stars can directly contribute to (U$``$B) colors significantly less than 0.
The derived colors are compared to evolutionary synthesis models. Figure 4 illustrates the (U$``$B, B$``$I) colors derived from two of the most popular models: BC95 (first described in Bruzual & Charlot 1993; the updated version is BC95) and Starburst 99 (Leitherer et al. 1999). The BC95 models represent isolated, bare stellar ensembles (photospheres), whereas the Starburst 99 models also include nebular continuum emission, which has an effect at both very short ($`<`$ 4000ร
) and very long ($`>`$ 1.5 $`\mu `$m) wavelengths (Leitherer & Heckman 1995). BC95 extends to much greater ages than Starburst 99, but it is unlikely that any of the star-forming regions observed here are so old. Neither model includes the effects of dust, which will both redden the intrinsic colors and also deplete the nebular continuum. In this paper most of the ages will be based on the BC95 models, since it is easier to understand how processes such as reddening affect the observed colors when applied directly to stars without additional complications such as nebular continuum which may not be well-understood.
The age estimates based on color diagrams of this sort are relatively insensitive to the exact shape of the IMF, which in this case is a Salpeter initial mass function (IMF) with upper and lower mass cutoffs of 125 and 0.1 M, respectively.. This is because the colors of a stellar (starburst) ensemble are dominated by the most luminous stars since what is actually observed is a luminosity-weighted average of all the stars in the burst (Leitherer 1996). It would require very extreme changes in the mass index of the IMF in order for late-type stars to dominate the colors by sheer numbers. The low mass end of the IMF primarily acts as a reservoir for the mass of the burst; most of the luminosity is emitted by only a small ($``$ 10%) mass fraction of the burst. Truncation of the lower end of the IMF, therefore, primarily adjusts the luminosity per unit mass of the burst, and thus the rate at which the gas supply is consumed. For example, truncating the IMF at 3 M in the BC95 model above only increases the bolometric luminosity per unit mass by a factor of 5 and leaves the color evolution almost unchanged until roughly 1 Gyr (at which time the entire stellar population has left the main sequence and therefore becomes very red; Charlot et al. (1993)). Similarly, adjusting the upper mass cutoff from 120 to 60 M leaves the color evolution unchanged after the first few million years due to the very rapid evolution of the upper end of the IMF. This model is thus likely to be a good one for a typical long-lived stellar population extending from the most massive stars to the least massive. Changing the upper and lower mass cutoffs is likely to only change the bolometric luminosity per unit mass, and even then only by a factor of a few.
The short wavelength colors are most strongly affected by dust and metallicity. Figure 7 presents different metallicity models from Starburst 99 (BC95 assumes solar metallicity). Shown are solar metallicity, twice solar metallicity, and a 0.05 solar metallicity case. Note that the solar and twice-solar metallicity models have nearly the same colors. Given that the local ULIGs appear to be collisions of mature, gas-rich spiral galaxies (Sanders et al. 1988a), it is not implausible that the gas is of near solar-metallicity. Unlike the situation in the near- infrared (i.e., as in Paper II) direct thermal dust emission cannot affect the short wavelength colors as these wavelengths correspond to temperatures far above the dust sublimation temperature. This leaves reddening as the dominant source of color variance.
Two dust reddening models are considered. The first is a simple foreground screen model, where all the dust lies on a path between the emitting source and the observer. The reddening law of Rieke & Lebofsky (1985) was used. This is illustrated in Figure 8 with a solid line vector corresponding to an extinction of 1 magnitude at V-band. Ages for the starbursts are derived by dereddening their observed colors until they intersect the modeled colors for a solar metallicity starburst. Fortunately, objects can only be made more red by extinction, and thus the colors can define upper limits to the stellar ages. However, while this may be a usable extinction model in the near-infrared, at wavelengths as short as U scattering may play an increasingly important role. Witt et al. (1992) have shown that for some scattering geometries it is difficult to achieve effective extinctions at U- band much greater than a few magnitudes even when the line of sight extinction is as high as 25 magnitudes. The โdusty galaxyโ model of Witt et al. (1992) in which dust and stars uniformly fill a spherical volume was used, and which was also used in the multi-color analysis of Surace et al. (1999). This is motivated by the observed morphology of star-formation in the nuclear regions of these galaxies (Surace et al. 1998,1999). Hill et al. (1995) find that the Witt et al. (1992) โdusty nucleusโ model provides a better fit to ultraviolet observations of HII regions in M81. However, this geometry is more appropriate to individual HII regions and not to the galaxy as a whole. In any case, for the maximum possible effective extinctions allowed by the U observations the dusty nucleus and dusty galaxy models have nearly identical broadband color properties. Figure 8 illustrates the reddening path of the โdusty galaxyโ model in units of total optical depth at V-band of $`\tau `$=0, 5, 10, and 15. Qualitatively it is very similar to the foreground extinction model. However, the reddening vector for the dusty galaxy model is slightly flatter due to the decreased effective optical depth at U as a result of the increased scattered light component relative to B and I-band. As a result, derived ages based on this model will typically be 20% higher than those based on a foreground screen. All ages quoted here are upper limits based on the foreground dust extinction model, since it is simpler in application.
#### 3.2.2 Observed Colors
The integrated (U$``$B) colors of the ULIGs span a very large range from -0.5 (IRAS 12112+0305) to 1.1 (UGC 5101), with a median of 0.4. Such an enormous range in integrated colors is similar to the Larson & Tinsley (1978) result, and probably does not constrain the galaxiesโ stellar populations significantly. Due to the wide range in colors, another approach must be used. It is more meaningful to address the colors of the individual compact high surface brightness features observed in the ULIGs, since these features correspond to the young starburst component of the galactic stellar population. Figure 8 shows the colors of the compact, high surface brightness nuclear regions of the ULIGs presented in Table 3.
The lack of resolved U knots among many of the cool ULIGs forces an examination of their โnuclearโ colors. In Surace (1998) it was shown that the โnuclearโ regions (defined by a diameter of 2.5 kpc) typically had optical/near- infrared colors similar to young star-forming regions with contributions from hot dust and mild (A<sub>V</sub>= 2โ4 magnitudes) of reddening. The median nuclear colors for the cool ULIGs observed at U are (U$``$B, B$``$I)=(0.1, 2.1), with (U$``$B) ranging from 1.2 to $``$0.6. The median age corresponding to the colors of the individual nuclei are 44 Myrs and 35 Myrs for the BC95 and LH99 models, respectively. Although it is not possible to discriminate between AGN and starburst activity based on UBI colors alone, the lack of a K excess in the cool ULIGs (Surace 1998) supports the interpretation of the blue (U$``$B) colors as being indicative of young stars. The presence of young stars in the cool ULIG nuclei (and in particular the single nucleus systems such as UGC 5101) indicates the presence of on-going young star formation, in contrast to the tidal tails โif a burst has occurred there it has already faded to the point of undetectability. In several cases U knots are resolvable within a radius of 4$``$5 kpc of the nucleus. In IRAS 12112+0305 the knots flanking the northern nucleus as well as the knots in the southern arc have dereddened ages of 34 and 6 Myrs, respectively. The southern knot in Mrk 273 has colors most similar to a 1 Gyr old burst (although this seems unlikely due to the fading discussed above), while the western knot appears to be about 200 Myrs. Finally, IRAS 23365+3604 has knots ranging between 3โ60 Myrs.
Examining the โwarmโ ULIGs is more problematic, since there is evidence for AGN activity in the nuclei of all of these systems (Sanders et al. 1988b, Papers I & II), and since AGN have (U$``$B,B$``$I) colors very similar to those of a young starburst. It is therefore likely that any examination of the nuclear (2.5 kpc diameter) regions will simply reflect the presence of AGN light. Instead, we break the warm ULIGs into several subcategories.
The double nucleus systems have the most complex morphology. Of the 3 systems with double nuclei, IRAS 08572+3915 and Mrk 463 each have many U knots, while Pks 1345+12 seems to have few. Although this latter case may be a result of the low S/N of our observations, the HST data (Paper I) also shows little in the way of knot structure as well. The knots in the tails of IRAS 08572+3915 appear to have dereddened ages of 5โ10 Myrs, while those in Mrk 463 vary from 4-90 Myrs, although the median is only 6 Myrs. The stellar galactic nuclei (based on imaging at other wavelengths) IRAS 08572+3915e and Mrk 463w have (U$``$B, B$``$I) colors of ($``$0.3, 1.6) and ($``$0.5,1.5) yielding dereddened ages of 45 and 25 Myrs, respectively. Similarly, it was shown in Paper II that the optical emission from IRAS 08572+3915w is predominantly starlight, since the AGN only begins to contribute at near-infrared wavelengths. For it we derive a starburst age of 37 Myrs. Pks 1345+12 has integrated colors so red (U$``$B, B$``$I)=(0.8, 2.5) as to be similar to a very old, late-type stellar population, although this could also be the result of A<sub>V</sub> = 2โ3 magnitudes of intervening dust.
The single nucleus systems have simpler morphology with fainter tidal features. The integrated colors of IRAS 01003$``$2238 indicate a maximum age of 40 Myrs, although at least some of the knots observed with WFPC2 are at most 5 Myrs in age (Paper I). Of the single-nucleus non-QSO systems, Mrk 231, IRAS 12071$``$0444, and IRAS 05189$``$2524 have the most well-developed tidal structure, but only the former two have any evidence for resolvable starburst knots. This is probably due to the lower spatial resolution at U ; there is evidence in the WFPC2 data (Paper I) for an extended starburst region in the nucleus of IRAS 05189$``$2524, but which cannot be resolved here. The estimated age of the southern, presumably stellar, region of IRAS 12071$``$0444 is 125 Myrs. The northern region in the vicinity of the putative active nucleus contains several knots of star formation which we cannot resolve here, and which may be quite young. IRAS 15206$``$3342โs western half has colors of ($``$0.4, 1.5), which are similar to those of a 20 Myr old burst. The eastern component shows a peculiarly large U excess ($``$0.9, 1.6), which may be indicative of contamination by the putative active nucleus (Paper I). The same is likely to be true of the nuclear colors of IRAS 05189$``$2524, whose dereddened colors are similar to a 20 Myr old starburst. Mrk 231โs โhorseshoeโ has dereddened colors similar to a starburst 225 Myrs old. The extended structure in the QSO I Zw 1 appear similar to the spiral arms of a spiral host galaxy. Mrk 1014 has several knots embedded in its tidal tails, 2 of which were detected with WFPC2. Unfortunately, they are sufficiently faint that their colors are not well determined; they are likely to be in the age range 3โ30 Myrs. Lastly, the extended features to the south and east of IRAS 07598+6502 were not obviously detected; however, these features were of fairly low effective surface brightness due to the seeing. A careful examination of the region as a whole indicates (U$``$B) similar to a stellar population of 5โ10 Myrs.
It is apparent from Figure 8 that most of the โwarmโ ULIGs have bluer nuclear colors than the โcoolโ ULIGs. However, their dereddened ages are extremely similar. Thus, while the redder colors might be indicative of greater age, it is more likely that they indicate a greater degree of dust obscuration, either by virtue of more line-of-sight dust or by having a dust geometry less favorable to scattering and thus a greater effective optical depth. This would be consistent with the evolutionary scenario proposed by Sanders et al. (1988a), in which dust is cleared away from the AGN in the warm ULIGs; presumably the same process would clear dust away from the star formation regions as well. The broad-band colors do not show a detectable maximum age difference between the two samples of ULIGs (although this is poorly constrained ). This requires that the dust clearing time be relatively short compared to the timescale of the nuclear burst ($``$ 100 Myrs) in order for the dust geometry to be the dominant color determinant above and beyond the effects of stellar aging.
Figure 9 illustrates the nuclear colors of single vs. double nucleus ULIGs. Unfortunately, several systems known to have double nuclei at other wavelengths either could not be resolved or had confused detections at one or more wavelengths (Pks 1345+12, IRAS 22491-1808, IRAS 12112+0305, IRAS 0119-2307). As a result the density of points in the figure is deceptive. Although the single nucleus systems are of necessarily greater dynamical age, this is not clearly reflected in their nuclear colors which have similar dereddened ages for both types of systems. This would argue against the evolutionary scenario, since it implies that dynamical age is not linked to dust reddening, which appears from the above to be linked to the mid-infrared colors of the ULIGs. However, the blue double nucleus point are strongly affected by IRAS 08572+3915 and Mrk 463. Since exclusion of these would produce different results, it is likely that there are too few galaxies in the sample to as yet draw meaningful conclusions.
### 3.3 Luminosities
#### 3.3.1 Observed Total U-band Luminosities
The total luminosities $`\nu `$L<sub>ฮฝ</sub> at U-band were computed using the conversion that a magnitude zero star has a flux density F<sub>ฮฝ</sub> of 1810 Jy in the U filter (Bessell 1979). The average rest-frame luminosity (i.e. the total amount of energy in units of 3.9$`\times `$10<sup>33</sup> ergs s<sup>-1</sup> in the U-band filter) is $`\nu `$L<sub>ฮฝ</sub>= 10<sup>10.7</sup> L, while the median is 10<sup>10.0</sup> L. This high average value is strongly influenced by the nuclear emission of the QSOs I Zw 1, Mrk 1014, and IRAS 07598+6508. Excluding these three systems yields an average that is also $`\nu `$L<sub>ฮฝ</sub>= 10<sup>10.0</sup> L. In terms of absolute magnitude, the median is M<sub>U</sub>= $``$20.8 and the average is $``$21.3. These results are consistent with the values extrapolated from the far-UV observations of IRAS 22491$``$1808 and IRAS 12112+0305 (Trentham et al. 1999). Given that infrared-selected QSO host galaxies are typically similar in luminosity to their nuclei (McLeod et al. 1994, Surace 1998), the ULIG QSOs are likely to have hosts similar in luminosity to the mean luminosity of the ULIGs.
#### 3.3.2 Contribution of U-band structure to the bolometric luminosity
We consider the probable bolometric luminosity of the observed knots and nuclear regions. The U-band bolometric correction (BC) is relatively robust against changes in the IMF mass cutoffs, since the energetics are dominated by the most massive stars. Indeed, adjustments to the IMF mass cutoffs, while changing the bolometric luminosity per unit mass by a factor of a few, have little effect on color as explained above, and hence the bolometric correction remains unchanged. Throughout the starburst age range of 0-1 Gyr, the bolometric correction at U is predicted by both BC95 and Starburst99 to vary only from about 0.5 to 2, and is very nearly 1 during the span of 10-600 Myrs for the modeled starburst. The likely bolometric luminosity of the star formation observed at U was determined by applying the bolometric correction applicable to a given measured emission region based on itโs estimated dereddened upper age limit, which was in turn based on its UBI colors. This was then increased by dereddening by the amount indicated by the UBI colors (typically A<sub>V</sub> = 1โ1.5 magnitudes). These estimates were all based on a foreground dust extinction model.
The cool ULIG bolometric luminosities due to current star-formation were calculated by considering the emission inside the nuclear (2.5 kpc diameter) regions which are known optically to host complex star-forming knots which we cannot resolve here, as well as any observed knots outside this region. The estimated nuclear L<sub>bol</sub> ranges from 10<sup>9.7</sup>L (UGC 5101) to 10<sup>11.3</sup>L (IRAS 00091$``$0738), with a median of 10<sup>10.7</sup>L and 75% lying in the range 10<sup>10.4</sup>โ10<sup>10.9</sup>L . The star-formation observed at U thus accounts for approximately 3% of the bolometric luminosity in a typical cool ULIG, assuming an average L<sub>bol</sub> of 10<sup>12.3</sup>L(Sanders et al. 1988a). This is extremely similar to the result derived by Surace (1998) from optical/near-infrared wavelengths.
The warm ULIGsโ star-formation budget was tallied by adding up all the emission in the observed knots as well as non-AGN nuclear regions (i.e., IRAS 08572+3915e) and all the nuclear regions where star-formation is believed to dominate at U based on WFPC2 B-band imaging (e.g., IRAS 08572+3915w, IRAS 12071$``$0444). The warm ULIG bolometric luminosities vary from 10<sup>9.9</sup> (IRAS 08572+3915) to 10<sup>11.7</sup>L(IRAS 15206+3342), with a median of 10<sup>10.6</sup>L. Again, this is similar to results found at other wavelengths. The warm ULIGs (not including 3c273) have a mean L<sub>bol</sub>=10<sup>12.26</sup>L; the average warm ULIG has a contribution of $``$2% by star formation observed at U to the bolometric luminosity, although this number is as high as 30% in at least one case (IRAS 15206+3342).
Estimation of the bolometric luminosity is more problematic in the presence of mixed stars and dust with scattering. Witt et al. (1992) point out that while the effective optical depth at short wavelengths like U may be quite low due to scattering, the true optical depth may be many times higher. As a result, while the galaxy may be quite luminous at U , a significant amount of longer wavelength emission is absorbed and reradiated in the thermal infrared due to the high total optical depth. In this case, the colors of the ULIGs are consistent with young stars mixed with dust mixed and having a total V-band optical depth of 5โ20. The โdereddeningโ correction will increase the bolometric luminosity by a factor 4โ10$`\times `$ greater than that of the foreground screen alone (Witt et al. 1992), and the observed star formation could contribute a sizable fraction of the known bolometric luminosity of the systems. This is unlikely, however. If the observed extended star formation did contribute a significant fraction of the bolometric luminosity, then the galaxies would be extended at every wavelength in a manner similar to that at U . However, it is known that the ULIGs become increasingly more compact in the near-infrared (Surace et al. 1999, Scoville et al. 1999) and ultimately are unresolved from the ground at mid-infrared wavelengths (Soifer et al. 1999). Since the latter are near the peak of the spectral energy distribution, this implies that the star-forming knots seen at U are not significant contributors to the bolometric luminosity, and that the dusty galaxy model is not a good approximation to the actual dust distribution.
## 4 Implications for High-Redshift Systems
Our 3410ร
observations of low-redshift ULIGs are applicable to observations of high-z galaxies at optical and near-infrared wavelengths. In particular, rest-frame U observations correspond to R-band observations at z=0.88 and H-band at z=3.9. Recent observations of the Hubble Deep Field (HDF) with ISO (Rowan-Robinson et al. 1998) have shown that the preponderance of strong mid-IR emitters are apparently interacting, star- forming galaxies. Deep imaging surveys with the Sub-mm Common User Bolometer Array (SCUBA) instrument on the James Clerk Maxwell telescope (JCMT) on Mauna Kea have recently discovered large numbers of sub-mm luminous high redshift galaxies (Smail et al. 1997). The projected star- formation rates for these galaxies are as high as 100 M/year; the ULIGs are the only local analog of starbursts with such high star formation rates (Hughes et al. 1998, Barger et al. 1998). Like the high redshift sub-mm galaxies, ULIGs are also found in interacting systems, and have SEDS with extremely strong far-IR components. Therefore, a more thorough understanding of the properties of ULIGs is likely to be the best opportunity for understanding the properties of these systems at high redshift.
The local ULIGs are observed to have angular diameters as large as 30 kpc. If the high redshift sub-mm galaxies are mergers of L galaxies similar to the low redshift ULIGs, then they will still be resolvable from the ground with diameters of 5โณ at z=1 and 6โณ at z=3 for a q<sub>0</sub>=0.5 universe (Peebles 1993). For q<sub>0</sub>=0 these are 4โณ and 3.5โณ. This is similar in many cases to the observed angular sizes of the optical counterparts of the sub-mm galaxies (Smail et al. 1998). The observed surface brightness of the galaxies is independent of q<sub>0</sub> and varies as (1+z)<sup>4</sup>. The apparent surface brightness of an ULIG at z=0.1 decreases by 2.6 magnitudes at z=1 and 5.6 magnitudes at z=3. The extended tidal structures in the ULIGs have typical surface brightnesses of 21-23 magnitudes per square arcsecond at U . The equivalent depth at z=1 in R-band ($``$ 25โ27) is easily reached in less than an hour by Keck (Hogg et al. 1997). The brightest U extended tidal features at z=3 will have surface brightnesses approaching H=27, and are probably unattainable by ground-based telescopes. However, the new generation of 8m class telescopes, equipped with adaptive optics, will have angular resolutions of approximately 40 milliarcseconds at H. This will provide a physical spatial resolution at high redshift comparable to or higher than that achievable for local ULIGs now from the ground at U . The star- forming clusters, which have typical HST resolved sizes of 100 pc (Surace et al. 1998), will still be unresolved. The brightest of these unresolved knots, as well as the nuclei, will approach H=25โ26, and will be detectable by a large telescope in 1โ2 hours (due to the enormous increase in point source detectability afforded by high spatial resolution imaging in background-limited situations). If the star- forming galaxies now being seen at high redshift are a result of collisions and mergers, then future optical and near-infrared observations of sub-mm luminous star-forming galaxies at high redshift are likely to appear similar in appearance to the U observations of ULIGs, with tidal debris and clustered star formation clearly visible. The Smail et al. (1998) results do actually indicate that 6/16 optically identified high-z galaxies are interacting systems similar in appearance to local ULIGs. The ULIG mergers appear in many cases to also fuel an AGN; the rate of apparent AGN activity in high redshift submm- luminous galaxies might indicate whether the merger process, while efficient in fueling a black hole, is also capable of creating new ones. Current estimates of the rate of optically detected AGN in high redshift sub-mm galaxies are $``$ 20%. This is similar to the fraction of โwarmโ AGN-like systems seen in complete samples of ULIGs in the local universe (Kim 1995).
Observers will attempt to morphologically classify their observations of high redshift galaxies. Initial observations at high redshift may be of relatively poor spatial resolution, making it difficult to discern the high surface brightness features seen in many of the images of local ULIGs. Examinations of the radial profiles of objects at high redshift can help discern if they are โelliptical-likeโ or โdisk-likeโ. As shown in ยง3.1, if the high redshift systems are similar to those examined here, then they will predominantly have radial profiles intermediate between the two. However, it may be possible to image high redshift star-forming galaxies and determine if their detailed morphologies are consistent with the local merger ULIGs observed here. Preliminary analyses of high-redshift galaxies believed to be undergoing active star-formation (for example, in the Hubble Deep Field) do show evidence for r<sup>-1/4</sup> profiles, possibly due to the same merger activity seen in the local ULIGs.
## 5 Conclusions
1. Many ULIGs are luminous at U (median $`\nu `$L<sub>ฮฝ</sub>= 10<sup>10.0</sup> L) and have evidence for considerable small-scale structure. This structure is almost identical to that seen at optical wavelengths.
2. Many of the double-nucleus systems have luminous knots of U emission in their tidal tails, which are absent in the single-nucleus systems. This may indicate that the clustered star-formation in the tails is predominantly brief and occurs early in the merger process prior to the coalescence of the nuclei.
3. The star-forming knots in the cool and warm ULIGs have similar upper age limits (10โ100 Myrs), although those in the warm ULIGs are generally consistent with less dust obscuration. This is consistent with greater dust clearing if the warm ULIGs are older systems than the cool.
4. The bolometric luminosity of the star-formation observed at U is generally only a very small fraction ($``$ 3%) of the known bolometric luminosity, a result nearly identical to that derived at optical and near-infrared wavelengths. This implies that even over a span of a factor of 20 in visual extinction, the star-forming structure remains similar. Most of the clustered star-formation we can detect at optical and even at near-infrared wavelengths probably occurs in knots that are either not heavily extinguished along the line of sight or have a favorable scattering geometry. The source of most of the bolometric luminosity remains completely hidden at these wavelengths, a result directly supported by mid-infrared imaging (Soifer et al. 2000).
5. The radial profiles of the ULIGs at U appear to be similar to a Sรฉrsic profile with index n = 1.8. This is an intermediate form between the exponential disk and the r<sup>-1/4</sup> law. This is the expected result for a partially relaxed merger remnant.
6. If they are similar to the local ULIGs, high-redshift ULIGs observed at optical wavelengths will have resolvable tidal structure and star- forming features which can be detected with ground-based telescopes.
We thank Gerry Luppino and John Tonry for the construction of the Orbit camera, without whose remarkably high quantum efficiency these observations could not have been made. We also thank John Dvorak and Chris Stewart for their skill at operating the telescope. We thank Lee Armus and B.T. Soifer for providing valuable discussions on the ULIGs, Aaron Evans for his useful comments on early drafts of this text, and Alan Stockton for advocating the value of U-band observations. We thank an anonymous referee whose comments helped strengthen this presentation. J.A.S. was supported by NASA grant NAG 5-3370 and by the Jet Propulsion Laboratory, California Institute of Technology, under contract with NASA. DBS was supported in part by JPL contract no. 961566.
## Appendix A Notes on Individual Objects
The following notes are provided as a description of the individual targets and to aid in understanding the relationship between the U emission and features seen at other wavelengths.
IRAS 00091$``$0738 โ the U emission originates primarily in the stellar โnucleusโ, with very little emission coming from the dense tidal tail to the south.
I Zw 1 โ the U emission originates almost entirely (90%) in the Seyfert 1 nucleus. Both tidal arms can be traced at U and closely mirror the B-band morphology. In particular, the two condensations in the NW arm and the bright blue condensation in the SW arm are both easily detected at U .
IRAS 01003$``$2238 โ no features are discernible, which is expected from the small size scale of the extended structure seen in the WFPC2 images (Surace et al. 1998) and the poor ground-based resolution. Itโs colors are similar to a starburst of age 10 Myrs.
IRAS 01199$``$2307 โ both nuclei are detected at U , but the emission is dominated by the SW nucleus by roughly 3.5:1. The tidal arms are undetected at U .
Mrk 1014 = PG 0157+001 โ the eastern tidal arm is very prominent at U , breaking into many smaller clumps and knots. The QSO nucleus accounts for about 86% of the total U emission. The integrated host galaxy color is U$``$B=0.72, similar to that of old stars.
IRAS 05189$``$2524 โ the U emission essentially mirrors that at all other wavelengths (Surace et al. 1998, 1999). The EW and NS tidal loops are easily seen. There are no apparent knots, although there may be some undetected in the resolved, extended nucleus.
IRAS 07598+6502 โ the star-forming knots detected at B are undetected at U . The nucleus accounts for 94% of the emission.
IRAS 08572+3915 โ large knots of emission trace the leading edges of all the tidal structures for their entire length. Many small knots are seen in the vicinity of the putative NW active nucleus. The SW nucleus dominates at short wavelengths, perhaps indicating that a strong starburst is taking place there. Both nuclei have (U$``$B) $`<`$0.25 in a 2.5 kpc diameter aperture; this, combined with their (B$``$I) colors from Surace et al. (1998), indicates the presence of stars that cannot be much more than 100 Myrs old in both nuclei.
UGC 5101 โ the U emission comes from a complex extended region in the core of the merger system. The tidal tails are barely detectable.
IRAS 12071$``$0444 โ clearly extended in the same manner as the WFPC2 images. Nearly all the compact U emission arises from the northern nucleus seen at B. The southern complex of knots is clearly detected at U .
IRAS 12112+0305 โ all of the tidal structure seen at B & I is seen at U . The condensations in the center of the southern tidal arc are clearly visible. The galactic โnucleusโ in the northern part of the system is dominated by a compact knot of emission at U , with another knot located on the short tidal arc extending to the west. Most notable by itโs absence is the star-like nucleus in the center of the system, which has completely disappeared between U and B-band, implying a line-of- sight extinction greater than 3 A<sub>V</sub>. This may explain the appearance of the galaxy shortward of 2200 ร
, since the Trentham et al. (1999) HST/FOC data indicate almost no detectable emission from the two nuclei themselves.
Mrk 231 โ The U emission comes almost entirely from the active nucleus and from the southern star-forming โhorseshoeโ. The horseshoe has (U$``$B)=0, yielding a maximum age of a few hundred Myrs. The corkscrewing structure extending from the NW of the nucleus and wrapping around through the west to the extreme south of the nucleus is also detected at U .
Mrk 273 โ Although the tidal tails are clearly detected at U , nearly all the emission arises in the central few kpc, which contains several large knots of U emission which are not seen in the WFPC2 I-band images (Surace 1998) or in CFHT K-band images (Knapen et al. 1997) In particular, the extremely prominent southern U knot does not appear in either of these images. The โnucleusโ is the luminous source to the north.
Pks 1345+12 โ this galaxy is very red ( U$``$B=0.73, B$``$I=2.54) and hence difficult to detect with high S/N at U . However, it appears that most of the U emission arises in the vicinity of the western nucleus.
Mrk 463 โ the eastern nucleus is distinctly elongated N-S, and the western nucleus is elongated E-W exactly as expected from the WFPC2 images of Paper I. Nearly all of the U emission in the eastern nucleus comes from the position of the northern โknotโ, and not the southern infrared โnucleusโ (Surace & Sanders 1999), which is consistent with its interpretation as an AGN ionization feature. Although the eastern nucleus is physically smaller, and hence has a higher peak surface brightness, it is the western nucleus that dominates in overall luminosity at U . All of the โknotsโ identified in the tidal ring structure are extremely luminous at U . Their colors indicate ages less than 10 Myrs.
IRAS 14348$``$1447 โ again, the U emission mirrors that at B. The blue knots seen at B surrounding the southern nucleus contribute heavily to the U emission in that nucleus, while the northern nucleus has a chain of knots in the base of the northern tail that contribute to the U emission.
IRAS 15206+3342 โ closely resembles the WFPC2 optical images (Paper I). Most of the U emission seems to originate in the compact nuclei in the southern half of the primary โstringโ of knots, consistent with their blue colors.
IRAS 15250+3609 โ the high surface brightness central galaxy component is well detected at U , as are the tidal features. Surprisingly, what appeared to be two separate features at I-band now appear to be a single arc bending away from the galaxy to the NE. The nearly circular arc to the south is undetected. The U emission is emitted by a compact source in the southern half of the central galaxy nucleus.
IRAS 22491$``$1808 โ most of the U emission originates in the western nucleus and northwestern tidal arm, and corresponds to the star clusters seen at optical wavelengths. Trentham et al. (1999) find that most of the emission shortward of 2200 ร
originates in the arm.
IRAS 23365+3604 โ the U emission closely mirrors that at B. All of the knots detected at B are luminous at U , and have derived (U$``$B) $`<`$0, indicating probable ages of 10 Myrs or less. The tails are almost undetectable; the most prominent extended emission at U is the arclike extension of the disk to the NW of the nucleus. Overall, the total integrated galaxy colors are similar to those of an old stellar population.
Figures 1 & 2 are images and are available in JPEG format from this preprint archive.
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# Entrapment of a Network of Domain Walls
\[
## Abstract
We explore the idea of a network of defects to live inside a domain wall in models of three real scalar fields, engendering the $`Z_2\times Z_3`$ symmetry. The field that governs the $`Z_2`$ symmetry generates a domain wall, and entraps the hexagonal network formed by the three-junctions of the model of two scalar fields that describes the remaining $`Z_3`$ symmetry. If the host domain wall bends to the spherical form, in the thin wall approximation there may appear non-topological structures hosting networks that accept diverse patterns. If $`Z_3`$ is also broken, the model may generate a buckyball containing sixty junctions, a fullerene-like structure. Applications to cosmology are outlined.
PACS numbers: 11.27.+d, 11.30.Er, 98.80.Cq, 47.54.+r
\]
Domain walls appear in diverse branches of physics, as for instance in systems of condensed matter that present ferromagnetic , ferroelectric and other properties , and also in cosmology . They arise in systems with at least two isolated degenerate minima, and in field theory they usually live in three spatial dimensions as bidimensional objects, seen as immersions into $`(3,1)`$ dimensions of static solutions of $`(1,1)`$ dimensional models that engender the $`Z_2`$ symmetry. The standard domain wall presents no internal structure, but there are models where they may entrap field configurations that engender non-trivial behavior. This idea follows as in Refs. , and in the more recent Refs. . Other investigations include for instance supersymmetry , supergravity , and the recent applications to polymers .
Although domain walls may be dangerous to cosmological applications, they have found their way into cosmology as for instance seeds for the formation of non-topological structures. This possibility appears in Refs. , where the discrete symmetry is changed to an approximate symmetry, or in Ref. , with the discrete symmetry biased so that domains of distinct but degenerate vacua spring unequally. The non-topological structures may be stable, but now stability requires the presence of conserved charges, of bosonic and/or fermionic origin.
Domain walls may also be of interest when they host non-trivial structures. We illustrate this point using the model introduced in the first work of Ref. , describing the pair of fields $`(\varphi ,\chi )`$ via the superpotential $`W(\varphi ,\chi )=\varphi +(1/3)\varphi ^3+r\varphi \chi ^2`$. Here $`r`$ is a parameter, real and dimensionless, that couples the two fields. The system is described by a quartic potential, and we use natural units, working with dimensionless space and time variables, and fields. The equations of motion for field configurations $`\varphi =\varphi (z)`$ and $`\chi =\chi (z)`$ are solved by solutions of the first order differential equations $`d\varphi /dz=1+\varphi ^2+r\chi ^2`$ and $`d\chi /dz=2r\varphi \chi `$. In this model, the sector connecting the minima $`(\pm 1,0)`$ is a BPS sector, with energy density or tension $`t=4/3`$. For $`r>0`$ this BPS sector admits two different types of static solutions: the one-field solutions $`\varphi _1(z)=\mathrm{tanh}(z)`$ and $`\chi _1=0`$, and the two-field solutions $`\varphi _2(z)=\mathrm{tanh}(2rz)`$ and $`\chi _2(z)=a(r)/\mathrm{cosh}(2rz)`$, with $`a^2(r)=1/r2`$, valid for $`0<r<1/2`$. We are working in $`(3,1)`$ space-time dimensions, so the one-field solution represents a standard domain wall, while the two-field solution appears as a domain wall having internal structure. As $`z`$ varies in $`(\mathrm{},\mathrm{})`$, in configuration space the vectors $`(\varphi _k,\chi _k)`$, $`k=1,2`$, describe a straight line segment $`(k=1)`$ and an elliptic arc $`(k=2)`$, resembling light in the linearly and elliptically polarized cases, respectively. These solutions also appear in condensed matter, and there they are named Ising and Bloch walls, respectively . They appear as solutions of the anisotropic $`XY`$ system, a model used to describe ferromagnetic transition in magnetic systems. The Bloch walls can be seen as chiral interfaces, and may be used to describe more complex phenomena, as for instance in the applications where chirality is also broken .
In Ref. the idea of nesting a network of defects inside a domain wall has been presented. This possibility may appear in models with three real scalar fields, engendering the $`Z_2\times Z_3`$ symmetry. In the present work we offer a model that contains the basic mechanisms behind this idea. The model will ultimately lead to the scenario of a domain wall hosting a network of defects, which may have direct interest to physics, as we show when we explore the pattern of the nested network in the case we allow the underlying $`Z_2\times Z_3`$ symmetry to be broken.
We first develop the idea of a domain wall hosting a network of defects, in a model described by three real scalar fields, with the (dimensionless) potential,
$`V(\sigma ,\varphi ,\chi )`$ $`=`$ $`{\displaystyle \frac{2}{3}}\left(\sigma ^2{\displaystyle \frac{9}{4}}\right)^2+\left(r\sigma ^2{\displaystyle \frac{9}{4}}\right)(\varphi ^2+\chi ^2)`$ (2)
$`+(\varphi ^2+\chi ^2)^2\varphi (\varphi ^23\chi ^2)`$
Here $`r`$ couples $`\sigma `$ to the pair of fields $`(\varphi ,\chi )`$. This potential is polynomial, and contains up to the fourth order power in the fields. Thus, it behaves standardly in $`(3,1)`$ space-time dimensions. Also, it presents discrete $`Z_2\times Z_3`$ symmetry. We set $`(\varphi ,\chi )(0,0)`$, to get the projection $`V(\sigma ,0,0)V(\sigma )=(2/3)(\sigma ^29/4)^2`$. The projected potential presents $`Z_2`$ symmetry, and can be written with the superpotential $`W(\sigma )=(2\sqrt{3}/9)\sigma ^3(3\sqrt{3}/2)\sigma `$, in the form $`V=(1/2)(dW/d\sigma )^2`$. The reduced model supports the explicit configurations $`\sigma _h(z)=\pm (3/2)\mathrm{tanh}(\sqrt{3}z)`$. The tension of the host wall is $`t_h=3\sqrt{3}=(3/2)m_h`$, where $`m_h`$ represents the mass of the elementary $`\sigma `$ meson. Also, the width of the wall is such that $`l_h1/\sqrt{3}`$.
The potentials projected inside $`(\sigma 0)`$ and outside $`(\sigma \pm \mathrm{\hspace{0.17em}3}/2)`$ the host domain wall are $`V_{in}(\varphi ,\chi )`$ and $`V_{out}(\varphi ,\chi )`$. Inside the wall we have
$`V_{in}(\varphi ,\chi )`$ $`=`$ $`(\varphi ^2+\chi ^2)^2\varphi (\varphi ^23\chi ^2)`$ (4)
$`{\displaystyle \frac{9}{4}}(\varphi ^2+\chi ^2)+{\displaystyle \frac{27}{8}}`$
This potential engenders the $`Z_3`$ symmetry, and there are three global minima, at the points $`v_1^{in}=(3/2)(1,0)`$ and $`v_{2,3}^{in}=(3/4)(1,\pm \sqrt{3})`$, which define an equilateral triangle. Outside the wall we get
$`V_{out}(\varphi ,\chi )`$ $`=`$ $`(\varphi ^2+\chi ^2)^2\varphi (\varphi ^23\chi ^2)`$ (6)
$`+{\displaystyle \frac{9}{4}}(r1)(\varphi ^2+\chi ^2)`$
$`V_{out}`$ also engenders the $`Z_3`$ symmetry, but now the minima depend on $`r`$. We can adjust $`r`$ such that $`r>9/8`$, which is the condition for the fields $`\varphi `$ and $`\chi `$ to develop no non-zero vacuum expectation value outside the host domain wall, ensuring that the model supports no domain defect outside the host domain wall. The restriction of considering quartic potentials forbids the possibility of describing the $`Z_3`$ portion of the model with the complex superpotential used in ; see also Ref. .
We investigate the masses of the elementary $`\varphi `$ and $`\chi `$ mesons. Inside the wall they degenerate to the single value $`m_{in}=3\sqrt{3/2}`$. Outside the wall, for $`r>9/8`$ they also degenerate to a single value, $`m_{out}(r)=3\sqrt{(r1)/2}`$, which depends on $`r`$. We see that $`m_{out}(r=4)=m_{in}`$. Also, $`m_{out}(r)>m_{in}`$ for $`r>4`$, and $`m_{out}(r)<m_{in}`$ for $`r`$ in the interval $`(9/8,4)`$.
We study linear stability of the classical solutions $`\sigma =\sigma _h(z)`$ and $`(\varphi ,\chi )=(0,0)`$. The fields $`\varphi `$ and $`\chi `$ vanish classically, and their fluctuations $`(\eta _n,\xi _n)`$ decouple. The procedure leads to two equations for the fluctuations, that degenerate to the single Schrรถdinger-like equation
$$\frac{d^2\psi _n(z)}{dz^2}+\frac{9}{2}V(z)\psi _n(z)=w_n^2\psi _n(z)$$
(7)
Here $`V(z)=1+r\mathrm{tanh}^2\sqrt{3}z`$. This equation is of the modified Pรถschl-Teller type, and can be examined analytically. The lowest eigenvalue is $`w_0^2=(3/2)\sqrt{6r+1}6`$. There is instability for $`r<5/2`$, showing that the host domain wall with $`(\varphi ,\chi )=(0,0)`$ is unstable and therefore relax to lower energy configurations, with $`(\varphi ,\chi )(0,0)`$ for $`r<5/2`$. Inside the host domain wall the sigma field vanishes, and the model is governed by the potential $`V_{in}(\varphi ,\chi )`$, which consequently may allow the presence of non-trivial $`(\varphi ,\chi )`$ configurations. The host domain wall entraps the system described by $`V_{in}(\varphi ,\chi )`$ for the parameter $`r`$ in the interval $`(9/8,5/2)`$. In this interval we have $`m_{out}<m_{in}`$, showing that it is not energetically favorable for the elementary $`\varphi `$ and $`\chi `$ mesons to live inside the wall for $`r(9/8,5/2)`$. The model automatically suppress backreactions of the $`\varphi `$ and $`\chi `$ mesons into the defects that may appear inside the host domain wall.
In Ref. the potential inside the wall was shown to admit a network of domain walls, in the form of a hexagonal array of domain walls. In the thin wall approximation the network may be represented by the solutions
$`\varphi _1`$ $`=`$ $`{\displaystyle \frac{3}{8}}+{\displaystyle \frac{9}{8}}\mathrm{tanh}\left({\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{27}{8}}}(y+\sqrt{3}x)\right)`$ (8)
$`\chi _1`$ $`=`$ $`{\displaystyle \frac{3}{8}}\sqrt{3}{\displaystyle \frac{3}{8}}\sqrt{3}\mathrm{tanh}\left({\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{27}{8}}}(y+\sqrt{3}x)\right)`$ (9)
and by $`(\varphi _k,\chi _k)`$, obtained by rotating the pair $`(\varphi _1,\chi _1)`$ by $`2(k1)\pi /3`$, for $`k=2,3`$. We identify the space $`(\varphi ,\chi )`$ with $`(x,y)`$, so rotations in $`(\varphi ,\chi )`$ also rotates the plane $`(x,y)`$ accordingly. The energy or tension of the individual defects in the network is given by, in the thin wall approximation $`t_n=(27/8)\sqrt{3/2}=(9/8)m_{in}`$. In the nested network, the width of each defect obeys $`l_n\sqrt{8/27}`$. This shows that $`l_h/l_n=3/2\sqrt{2}`$, and so the host domain wall is slightly thicker than the defects in the nested network. In the thin wall approximation, the potential $`V_{in}(\varphi ,\chi )`$ allows the formation of three-junctions as reactions that occur exothermically, and the nested array of thin wall configurations is stable. In Fig. 1 we depict the hexagonal network of defects inside the domain wall, in the thin wall approximation. The dashed lines show equilateral triangles, that belong to the dual lattice. Both the hexagonal network and the dual triagular network are composed of equilateral polygons, a fact that follows in accordance with the $`Z_3`$ symmetry.
We now explore the breaking of the $`Z_2\times Z_3`$ symmetry of the model. The simplest case refers to the breaking of the $`Z_3`$ symmetry, without breaking the remaining $`Z_2`$ symmetry. We consider the case of breaking the internal $`Z_3`$ symmetry in the following way. We take for instance the vacuum state $`v_1^{in}=(3/2)(1,0)`$, and change its position to a location farther from or closer to the other minima of the system, increasing or decreasing the angle between two of the three defects. We illustrate this situation in Fig. 2. We do this without removing the degeneracy of the three minima. This mechanism changes $`Z_3Z_2`$, so we refer to it as the minimal breaking.
We notice that the energy of the defect depends on the distance between the two minima the defect connects, and goes with the cube of it. If the $`Z_3`$ symmetry is broken to $`Z_2`$, the three-junction is stable for $`\mathrm{cos}(\theta /2)=t_n/2t_n^{}`$, where $`\theta `$ is the angle between the modified defects, with tensions changed from $`t_n`$ to $`t_n^{}`$. This implies that $`t_n^{}>t_n`$ for $`\theta >2\pi /3`$, and $`t_n^{}t_n`$ for $`\theta 2\pi /3`$. If the vacuum state deviates significantly from its $`Z_3`$-symmetric position, we cannot neglect the correction to the energy of the defects, and this would changes the equilateral hexagons of Fig. 1 to non-equilateral hexagons. However, if the vacuum state deviates slightly from its $`Z_3`$-symmetric position, one may neglect the correction to the energy of the defects. In this case we are slightly breaking $`Z_3Z_2`$.
We now concentrate on breaking the $`Z_2`$ symmetry of the host domain wall. We can do this with the inclusion in the potential of a term odd in $`\sigma `$, that slightly removes the degeneracy of the two minima $`\sigma =\pm 3/2`$. Thus, the host domain wall bends trying to involve the local minimum, the false vacuum. To stabilize the non-topological structure we include charged fields into the system. The way one couples the charged fields is not unique, but if we choose to add fermions, we can couple them to the $`\sigma `$ field in a way such that the projection with $`(\varphi ,\chi )(0,0)`$ may leave the model supersymmetric. This is obtained with the superpotential $`W(\sigma )`$, with the Yukawa coupling $`d^2W/d\sigma ^2=(4/3)\sqrt{3}\sigma `$. In this case massless fermions bind to the host domain wall, and contribute to stabilize the non-topological defect that emerges with the breaking of the $`Z_2`$ symmetry.
The breaking of the $`Z_2`$ symmetry can be done breaking or not the remaining $`Z_3`$ symmetry of the model. We examine these two possibilities supposing that the host domain wall bends under the assumption of spherical symmetry, becoming a non-topological defect with the standard spherical shape. This is the minimal surface of genus zero, and according to the Euler theorem we can only tile the spherical surface with three-junctions as a regular polygonal network in the three different ways: with 4 triangles, or 6 squares, or yet 12 pentagons. These three cases preserve the $`Z_3`$ symmetry of the original network, locally, at the three-junction points. However, if locally one slightly breaks the $`Z_3`$ symmetry of the network to the $`Z_2`$ one, the three-junctions can now tile the spherical surface with 12 pentagons and 20 hexagons. We think of breaking the $`Z_3`$ symmetry minimally, to the $`Z_2`$ symmetry, through the same mechanism presented in Fig. 2. Thus, if the symmetry is broken slightly we can consider the defect tensions as in the regular hexagonal network.
The tiling with 12 pentagons and 20 hexagons generates a spherical structure that resembles the fullerene, the buckyball composed of sixty carbon atoms. We visualize the symmetries involved in the spherical structures thinking of the corresponding dual lattices, which are triangular lattices, but in the three first cases the triangles are equilateral, while in the fourth case they are isosceles. We recall that regular heptagons introduce negative curvature, so they cannot appear when the genus zero surface is minimal. However, they may for instance spring to generate higher energy states from the fullerene-like structure, locally roughening the otherwise smooth spherical surface.
We write the energy of the non-topological structure as $`E_{nt}^n=E_{nt}^s+E_n`$, where $`E_{nt}^s`$ stands for the energy of the standard non-topological defect, and $`E_n`$ is the portion due to the nested network. We use $`E_{nt}^s=E_q+E_h`$, which shows the contributions of the charged fields and of the host domain wall, respectively. We have $`E_h=St_h`$ and $`E_n=Ndt_n`$, where $`S`$ is the area of the spherical surface, and $`N`$ and $`d`$ are the number and length of segments in the nested network. In the thin wall approximation the radius $`R`$ of the non-topological structure should obey $`Rl_h`$, to make such structure much larger than the characteristic width of the host domain wall. We introduce the ratio
$$\frac{E_{nt}^n}{E_{nt}^s}=\left(1+\frac{N}{1+r_q}\right)\frac{t_n}{t_h}\frac{d}{S}$$
(10)
with $`r_q=E_q/E_h`$. The non-topological structure nests a network of defects, which modifies the scenario one gets with the standard domain wall. The modification depends on the way one couples charged bosons and fermions to the $`\sigma ,\varphi `$, and $`\chi `$ fields. However, if the $`Z_3`$ symmetry is locally broken to the $`Z_2`$ one, the most probable defect corresponds to the fullerene or buckyball structure. But if the $`Z_3`$ symmetry is locally effective, there may be three equilateral structures, the most probable arising as follows. We consider the simpler case of plane polygonal structures, identifying the tetrahedron ($`i=3`$), cube ($`i=4`$), and dodecahedron ($`i=5`$). We introduce $`R_{ij}`$ as the energy ratio for the $`i`$ and $`j`$ structures. We get
$$R_{ij}=\frac{1+r_q+\frac{9\sqrt{2}}{16h_i}}{1+r_q+\frac{9\sqrt{2}}{16h_j}},i,j=3,4,5$$
(11)
Here $`h_3,h_4`$, and $`h_5`$ stand for the radius of the incircle of the triangle, square, and pentagon, respectively. Energy favors the tetrahedron, which is self-dual because the network and its dual are the very same triangular lattice. There are two other configurations, the octahedron, dual to the cube, and the icosahedron, dual to the dodecahedron. They do not appear in the $`Z_2\times Z_3`$ model because they require four- and five-junctions, respectively.
The present work can be extended in several directions. For instance, in the $`Z_2\times Z_3`$ model, if the host domain wall bends cylindrically, one may get to nanotube-like configurations . Moreover, we could consider other models, presenting the $`Z_2\times Z_k`$ symmetry $`(k=4,5,6)`$, leading to $`k`$-junctions. This would allow to tile the plane with squares $`(k=4)`$, or triangles $`(k=6)`$, and the spherical surface with triangles, as the octahedron $`(k=4)`$ or the icosahedron $`(k=5)`$. These investigations show a new way of modifying the standard domain wall, and this may find direct application in condensed matter and in cosmology . In cosmology, for instance, the internal network modifies the wall energy and changes its cosmological evolution. We follow Ref. to see that the ratio between the volume pressure and the surface tension that govern the wall evolution now changes according to $`(p_V/p_T)^n=(p_V/p_T)^s/(1+\alpha t_n)`$; $`\alpha >0`$ is a numerical factor, and $`t_n0`$ recovers the standard case. This shows that the ending scenario now depends not only on the way the symmetry is broken, but also on the modification the internal network introduces, enlarging the possibility of the wall dominating the energy density before the volume pressure can act. Another line could follow Ref. , asking how fermions could be coupled such that the corresponding zero modes that inhabit the wall could also bind to the internal network. Such mechanism would lead to a scenario where the nested fermions could form a one-dimensional gas inside the nested network, changing the way the soliton star evolves. We also mention the investigation concerning pattern formation in cosmology, as in the recent work . Our investigation provides other possible scenarios for pattern formation in the early universe.
We thank C. Furtado, F. Moraes, J. R. S. Nascimento, and R. F. Ribeiro for discussions, and CAPES, CNPq, and PRONEX for partial support.
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