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# BOSE-EINSTEIN CORRELATIONS IN HEAVY-ION PHYSICS AND ELECTRON-POSITRON COLLISIONS
## 1 Theoretical Overview
Bose-Einstein correlations (BEC) are a phase-space phenomenon: Symmetrization of the multiparticle wave function affects the measured $`n`$-particle coincidence spectra and leads to an enhancement relative to the corresponding product of independent 1-particle spectra, if the emitted particles are close in phase-space (i.e. they occupy the same elementary phase-space cell). The spatial length of the elementary phase-space cells is limited by the geometric size of the source of particles with the considered momentum. The larger this size, the narrower these cells are in momentum space. By tuning the relative momenta and watching the onset of BEC effects one can thus measure the spatial length of the elementary phase-space cells and thereby the size of the source.
Wigner Functions. A description of BEC effects among $`n`$ particles thus involves the $`n`$-particle phase-space density. Since we are discussing a quantum mechanical phenomenon, we are not talking about a classical phase-space density (which has directly a probabilistic interpretation), but about the Wigner density (which is positive definite only when averaged over many elementary phase-space cells). If the particles are emitted independently, the (unsymmetrized) $`n`$-particle Wigner density factorizes, and all $`n`$-particle coincidence cross sections are expressible through the single-particle Wigner function $`S(x,p)`$. The assumption of independent particle emission is justifiable in heavy ion collisions where the many unobserved particles serve as a reservoir for all kinds of conserved quantities. In $`e^+e^{}`$ collisions this is much less obvious and needs to be tested experimentally.
Correlation Function. As long as the source has sufficiently low phase-space density that multi-particle symmetrization effects are dominated by two-particle exchange terms, the two-particle correlation function $`C(\text{tensysevensyfivesy}𝒒,\text{tensysevensyfivesy}𝑲)`$, defined as the ratio of the 2-particle coincidence spectrum $`P_2(\text{tensysevensyfivesy}𝒑_a,\text{tensysevensyfivesy}𝒑_b)`$ and the product of single-particle spectra $`P_1(\text{tensysevensyfivesy}𝒑_a)P_1(\text{tensysevensyfivesy}𝒑_b)`$ with $`\text{tensysevensyfivesy}𝒒=\text{tensysevensyfivesy}𝒑_a\text{tensysevensyfivesy}𝒑_b`$ and $`\text{tensysevensyfivesy}𝑲=(\text{tensysevensyfivesy}𝒑_a+\text{tensysevensyfivesy}𝒑_b)/2`$, is given by $`^\mathrm{?}`$<sup>a</sup><sup>a</sup>aWe here neglect Coulomb final state interactions since methods are known to correct the data for them.$`^\mathrm{?}`$
$$C(\text{tensysevensyfivesy}𝒒,\text{tensysevensyfivesy}𝑲)=𝒩\left(1+\frac{|_xS(x,K)e^{iqx}|^2}{_xS(x,p_a)_xS(y,p_b)}\right)=𝒩\left(1+\frac{P_1(\text{tensysevensyfivesy}𝑲)^2}{P_1(\text{tensysevensyfivesy}𝒑_a)P_1(\text{tensysevensyfivesy}𝒑_b)}\left|\frac{_xS(x,K)e^{iqx}}{_xS(x,K)}\right|^2\right).$$
(1)
Here $`_xd^4x`$, $`q^0=E_aE_b`$, $`K^0=(E_a+E_b)/2`$, and
$$P_1(\text{tensysevensyfivesy}𝒑)=_xS(x,p)\mathrm{with}p^0=E_p=\sqrt{m^2+\text{tensysevensyfivesy}𝒑^2}.$$
(2)
The normalization $`𝒩`$ depends on the multiplicity distribution via $`^\mathrm{?}`$ $`𝒩=n(n1)/n^2`$. In heavy ion collisions usually $`𝒩1`$. Due to the mass-shell constraint $`^\mathrm{?}`$ $`q^0=\text{tensysevensyfivesy}𝜷\text{tensysevensyfivesy}𝒒`$ (where $`\beta =\text{tensysevensyfivesy}𝑲/K^0\text{tensysevensyfivesy}𝑲/E_K`$ is the velocity of the particle pair) the Fourier transform in (1) is not invertible: the separation of temporal and spatial aspects of the emission function $`S(x,K)`$ requires additional model assumptions which must be provided by a physical picture of the time evolution of the source until freeze-out.$`^\mathrm{?}`$
The Reduced Correlator. While (1) goes to $`2𝒩`$ at $`\text{tensysevensyfivesy}𝒒=\mathrm{\hspace{0.17em}0}`$, real correlation functions usually approach a smaller value $`𝒩(1+\lambda )`$ with $`\lambda (\text{tensysevensyfivesy}𝑲)<1`$. Possible reasons are partial phase coherence in the source and decay contributions from long-lived resonances.$`^\mathrm{?}`$ To account for this one rewrites (1) as
$$C(\text{tensysevensyfivesy}𝒒,\text{tensysevensyfivesy}𝑲)=𝒩\left(1+\lambda (\text{tensysevensyfivesy}𝑲)𝒦(\text{tensysevensyfivesy}𝒒,\text{tensysevensyfivesy}𝑲)\right)=𝒩\left(1+\lambda (\text{tensysevensyfivesy}𝑲)\frac{P_1(\text{tensysevensyfivesy}𝑲)^2}{P_1(\text{tensysevensyfivesy}𝒑_a)P_1(\text{tensysevensyfivesy}𝒑_b)}𝒦_{\mathrm{red}}(\text{tensysevensyfivesy}𝒒,\text{tensysevensyfivesy}𝑲)\right).$$
(3)
The reduced correlator $`𝒦_{\mathrm{red}}(\text{tensysevensyfivesy}𝒒,\text{tensysevensyfivesy}𝑲)`$ is given by the last term in (1) which contains the information about the space-time structure of $`S(x,K)`$. To isolate it one constructs $`C(\text{tensysevensyfivesy}𝒒,\text{tensysevensyfivesy}𝑲)`$ from the measured 1- and 2-particle cross sections, applying the Coulomb correction, determines $`𝒩`$ and $`\lambda (\text{tensysevensyfivesy}𝑲)`$ from the limits $`q\mathrm{\hspace{0.17em}0}`$ and $`q\mathrm{}`$, divides by $`𝒩`$ and subtracts the 1, and finally divides the result by $`\lambda (\text{tensysevensyfivesy}𝑲)`$ and the measured ratio of single particle cross sections $`P_1(\text{tensysevensyfivesy}𝑲)^2/P_1(\text{tensysevensyfivesy}𝒑_a)P_1(\text{tensysevensyfivesy}𝒑_b)`$. For large sources like those in heavy ion collisions this ratio is close to unity,$`^\mathrm{?}`$ but for small sources like those in $`e^+e^{}`$ it can contribute significantly to the tensysevensyfivesy$`𝒒`$-dependence of $`C(\text{tensysevensyfivesy}𝒒,\text{tensysevensyfivesy}𝑲)`$; it is then important to divide it out before trying to extract the source size. So far we have seen no data analysis where this is done! Instead, one usually extracts the size directly from $`𝒦(\text{tensysevensyfivesy}𝒒,\text{tensysevensyfivesy}𝑲)`$, without dividing out the 1-particle spectra. As we will see, this can be quite misleading.
Source Radii from BEC. One usually characterizes $`^\mathrm{?}`$ the source function $`S(x,K)`$ by its norm, center and space-time variances (widths), all of which are generally functions of the momentum tensysevensyfivesy$`𝑲`$ of the emitted particles. In this “Gaussian approximation” the reduced correlator reads
$$𝒦_{\mathrm{red}}(\text{tensysevensyfivesy}𝒒,\text{tensysevensyfivesy}𝑲)=\mathrm{exp}\left[q^\mu q^\nu \stackrel{~}{x}_\mu \stackrel{~}{x}_\nu (\text{tensysevensyfivesy}𝑲)\right],$$
(4)
where $`\stackrel{~}{x}_\mu \stackrel{~}{x}_\nu =x_\mu x_\nu x_\mu x_\nu `$, with
$$x_\mu x_\nu (\text{tensysevensyfivesy}𝑲)=\frac{_xx_\mu x_\nu S(x,K)}{_xS(x,K)},$$
(5)
are the space-time variances of the emission function (effective source sizes). Different conventions for resolving the mass-shell constraint $`q^0=\text{tensysevensyfivesy}𝜷\text{tensysevensyfivesy}𝒒`$ and expressing (4) in terms of three indepenent components of $`q`$ lead to different Gaussian parametrizations for the correlator.$`^\mathrm{?}`$ The corresponding Gaussian width parameters, the “HBT (Hanbury Brown - Twiss) radii”, are then combinations of the variances $`\stackrel{~}{x}_\mu \stackrel{~}{x}_\nu (\text{tensysevensyfivesy}𝑲)`$ and thus functions of the pair momentum tensysevensyfivesy$`𝑲`$.
## 2 Bose-Einstein Correlations in Heavy Ion Collisions
Due to space reasons we will be very short – detailed discussions can be found elsewhere.$`^{\mathrm{?},\mathrm{?}}`$ For Pb+Pb collisions at the SPS it was found that the pion emitting source is a rapidly expanding fireball in approximate local thermal equilibrium which at decoupling has a temperature of about 100 MeV and expands nearly boost-invariantly in the longitudinal direction while the average transverse expansion velocity is a bit larger than half the light velocity. The collective expansion manifests itself in a strong and characteristic dependence of the space-time variances $`\stackrel{~}{x}_\mu \stackrel{~}{x}_\nu `$ of the effective source $`S(x,K)`$ on the pair momentum tensysevensyfivesy$`𝑲`$. This implies a corresponding tensysevensyfivesy$`𝑲`$-dependence of the HBT radii extracted from (4). The pion emission process lasts only for about 2-3 fm/$`c`$ but it doesn’t begin until at least 6-8 fm/$`c`$ after the collision. Freeze-out thus is a rather sudden process at the end of an extended rescattering and expansion stage. It is important to stress that the separation of longitudinal and transverse flow and access to the emission duration $`\stackrel{~}{t}^2`$ is only possible in a full-fledged 3-dimensional and tensysevensyfivesy$`𝑲`$-dependent analysis of the correlation function $`C(\text{tensysevensyfivesy}𝒒,\text{tensysevensyfivesy}𝑲)`$. Projections to lower dimensionality (e.g on $`q_{\mathrm{inv}}^2`$) lead to uncontrollable and unrecoverable loss of information.
## 3 Bose-Einstein Correlations in $`e^+e^{}`$ Collisions
As stated in Sec. 1, to compute Bose-Einstein correlations one needs information on the Wigner phase-space density of the source. Going simulation programs of particle production in high-energy $`e^+e^{}`$ collisions like PYTHIA, JETSET and HERWIG provide only momentum-space information on the produced particles. This is not enough to calculate BEC effects. Different methods have been suggested to provide the missing coordinate-space information, either directly or indirectly.$`^\mathrm{?}`$ We previously studied $`^\mathrm{?}`$ BEC in VNI which studies the time evolution of the collision in phase-space. Here we present some very early results based on a phase-space version $`^\mathrm{?}`$ of JETSET 7.4 which provides both the momenta and production coordinates for the produced particles. Our version of this code distributes the transverse distance of the production points from the central string axis according to a Gaussian with rms radius of 0.78 fm while S. Todorovova’s version $`^\mathrm{?}`$ puts the production points right on the string axis. This latter procedure is inconsistent with the uncertainty relation, and we found accordingly $`^\mathrm{?}`$ that it produces correlation functions which rise as a function of $`q_s`$ instead of decaying.
The algorithm for computing the correlation function from the positions and momenta of the generated pions is described elsewhere;$`^\mathrm{?}`$ we use the “classical” algorithm without wave packet smearing.$`^\mathrm{?}`$ In order to test the space-time structure of the events generated by JETSET and the BEC afterburner, we begin with a simple event topology ($`e^+e^{}Z^0q\overline{q}2`$ jets) and consider only directly produced pions, thus avoiding the multiscale problems associated with longlived resonance decays. We analyse the correlation function in a Cartesian coordinate system where the longitudinal ($`l`$ or L) axis is along the direction defined by the relative momentum of the initial $`q\overline{q}`$ pair ($``$ jet axis), the outward ($`o`$ or T) direction is defined by the transverse pair momentum $`\text{tensysevensyfivesy}𝑲_\mathrm{T}`$, and the sideward ($`s`$) axis points in the third direction.
The top two left panels of Fig. 1 show the correlator in the side direction. The reduced correlator $`𝒦_{\mathrm{red}}`$ is seen to be independent of $`K_\mathrm{T}`$ and always reproduces the input rms width of the string: $`R_s=r_{\mathrm{rms}}/\sqrt{2}=\mathrm{\hspace{0.17em}0.55}`$ fm. In contrast, $`𝒦`$ does depend on $`K_\mathrm{T}`$, and for small $`K_\mathrm{T}`$ it produces smaller HBT radii (0.31, 0.41, 0.46 and 0.50 fm at $`K_\mathrm{T}=\mathrm{\hspace{0.17em}0}`$, 0.3, 0.5 and 1.0 GeV, respectively). This effect is an artifact induced by the ratio of 1-particle spectra in (3); it matters since the real radius is so small, producing significant errors if not divided out. For $`R_o`$ and $`R_l`$, on the other hand, its effect is in our calculation nearly negligible: these radii come out much larger than $`R_s`$. This, however, points to another problem: longitudinal HBT radii of up to 5 fm are incompatible with the data which give only about 1 fm (see the experimental talks in this session)! The problem seems to be connected with the large emission time duration $`\mathrm{\Delta }\tau `$ of up to 3 fm/$`c`$ at low $`K_\mathrm{T}`$. This parameter, which reflects the proper time distribution of string breaking processes in JETSET, is not fixed by 1-particle spectra, but it is seen to seriously affect the 2-particle correlations. We are presently trying to fix this problem. At this moment we can only say that the version of JETSET used by us disagrees with experiment at the level of 2-particle correlations.
Ackowledgement: This work was supported by the Deutsche Forschungsgemeinschaft.
## References
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# Minimal irreversible quantum mechanics. The decay of unstable states.
## I Introduction
In previous works we have studied the time evolution of a quantum oscillator coupled to a dense, but discrete (finite) bath of harmonic oscillators. For such system an irreversible behavior has appeared as a consequence of averaging in time the evolution of the characteristic quantum-oscillator variables (macroscopic quantities), since time evolution splits in very different scales: One related to small fluctuations, which are erased by averaging, another one related with recurrence phenomena, which are far enough of laboratory observational times, and the last one connected with observable phenomena, which involves irreversibility.
In Ref. it has been shown how to pass to the continuous bath. A resonance, which can be isolated and leads to an evolution for a macroscopic period of time (the same period as in the discrete case), has arisen because of the energy of the quantum oscillator is embedded in the continuum. One particularity of this continuous limit is that two of the three time scales have become irrelevant since they have zero measure with respect to the remaining time scale. However, this last scale is not exclusively governed by the resonance (associated with an exponential decay) but also by contributions coming from the semibounded feature of the energy spectrum.
Leaving these contributions aside the main behavior of the quantum oscillator is an exponential decay towards the equilibrium with the bath, and can be described only with the contribution due to the resonance. The present work is the conclusion of the previous papers , where a novel method was used to work directly in the continuum, including the exponential decay law in quantum mechanics. In this paper we continue with the development of the formalism we have called “Minimal irreversible quantum mechanics,” where time asymmetry can be represented through the choice of a subspace of “admissible” or “regular” solutions of the evolution equation.
The main idea goes as follows. Usually rigorous quantum mechanics must be formulated in a Gel’fand triplet
$$𝒮𝒮^\times ,$$
(1)
where:
$`𝒮`$ is the space of “regular states” or test-functions space, corresponding to Schwarz-class wave functions, that are considered as the “physical” states.
$``$ is the space of “states,” or Hilbert space, introduced to extend the notion of probability to a larger space and to use the well-known spectral theory of Hilbert spaces. These states correspond to square-integrable wave functions.
$`𝒮^\times `$is the space of “generalized states”, or rigged Hilbert space, namely the space of linear (or antilinear) functionals over $`𝒮`$, which are essentially used to find the spectral expansion of the regular states (e.g. Fourier expansions).
Let $`K`$ be the Wigner or time-inversion operator. As usual the evolution Hamiltonian $`H`$ is time symmetric, i.e.
$$KHK^{}=H.$$
(2)
In the wave function representation the action of $`K`$ coincides with the complex conjugation, so it is defined over $`𝒮`$ by
$$K\phi (x)=\phi ^{}(x).$$
(3)
Thus
$$K:𝒮𝒮.$$
(4)
Therefore $`𝒮`$ is also time symmetric.
But the real universe and macroscopic objects have clearly time-asymmetric evolutions. Therefore the task of this paper and the preceding ones is to explain how this time asymmetry appears while the quantum mechanical laws of the universe (embodied in $``$) are time symmetric. The usual and successful explanation is based on coarse-graining: Macroscopic objects have a huge number of dynamical variables and we can measure and control only a small number of them, the so-called relevant variables. If we neglect the rest of the variables, the irrelevant ones, we obtain time-asymmetric evolution equations. Nevertheless in paper (according to the line of thought pioneered in Refs. , , ) we have follow a different way. We have developed a sort of minimal irreversible quantum description, which reproduces time asymmetry from the basic microscopic level directly, where the key point is the presence of resonances (and additional hypotheses we have extracted from Ref. ).
Obviously we want to obtain the standard results making minimal changes to the well established and usual quantum mechanics. If we change Eqs. (2) or (3), it is almost sure to find experimental problems. So the minimal modification is to change Eq. (4) defining a new test-functions space $`\mathrm{\Phi }_+𝒮`$ such that
$$K:\mathrm{\Phi }_+\mathrm{\Phi }_{}\mathrm{\Phi }_+.$$
(5)
In this way $`K`$ is not even defined over the space of regular states $`\mathrm{\Phi }_+`$ and a time-asymmetric evolution arises.
This can be done if we postulate, as we have done in Ref. ,<sup>*</sup><sup>*</sup>*This postulate has been motivated in cosmological-global considerations in Refs. and . that all the “regular” or “admissible” states belong to a space $`_+\theta (H_+^2)`$ and also to $`𝒮.`$ Then $`\mathrm{\Phi }_+\theta (H_+^2𝒮)`$ $`[`$the time inverted states belong to a space $`_{}\theta (H_{}^2)`$ and $`\mathrm{\Phi }_{}\theta (H_{}^2𝒮),`$ respectively\],where $`\theta `$ is the Heaviside step function that gives the restriction to the positive real energy axis and $`H_\pm ^2`$ are the Hardy class function spaces .As spaces $`_{}`$ and $`_+`$ are isomorphic they are normally called $``$ .
An “irreversible” quantum theory based on a Gel’fand triplet
$$\mathrm{\Phi }_\pm _\pm \mathrm{\Phi }_\pm ^\times ,$$
(6)
is feasible and it yields physical results, as the dominant experimental decay of unstable states, if the test-function space $`\mathrm{\Phi }_+`$ is so chosen. This will be valid for systems where the existence of resonances dominates the evolution for the relevant period of observational time. We have shown that, what it is done in the quoted papers , , and is essentially a minimal modification of the ordinary reversible quantum theory. In fact, from now on we will consider that:
$`\mathrm{\Phi }_+`$ is the space of “regular states” or test-functions space, that are considered as the “admissible” states.
$`_+`$ is the space of “states,” or Hilbert space. These states are again particular square-integrable wave functions. But in paper and in this work we consider that only $`\mathrm{\Phi }_+`$ contains the “admissible” states.
$`\mathrm{\Phi }_+^\times `$is the space of “generalized states,” or rigged Hilbert space, namely the space of linear (or antilinear) functionals over $`\mathrm{\Phi }_+,`$ which are essentially used to find the spectral expansion of the “regular states” in Sec. III.
The spaces with subscripts “–” contain the time-inverted states of the corresponding spaces with subscripts “+”.
Friedrichs model was studied using this approach. In this work, we show that this idea can be used to take a slightly different point of view in studying dissipation phenomena of quantum Brownian motion. This more complex model will force us to generalize the definition of space $`\mathrm{\Phi }_+`$ although the roles played by the characters in the triplets $`\mathrm{\Phi }_\pm _\pm \mathrm{\Phi }_\pm ^\times `$ will remain the same.
Brownian motion has been extensively studied in the literature (we will only quote those papers particularly relevant to our line of work). E.g., in Ref. it was shown that for a system composed by a finite number of linear interacting oscillators a dissipative behavior can be found in the limit of a dense system (continuous spectrum). But, in this work we are concerned directly with dense systems with continuous spectrum. The presence of this continuous spectrum allows us to study the decay processes using analytical properties familiar in scattering theory . The model is a very well-known and widely used system, consisting of a harmonic oscillator coupled to an infinite and continuous bath. In this paper, as in Refs. , the bath is composed by an infinite collection of harmonic oscillators and the interaction is modelled to be linear and characterized by the spectral weight, but otherwise arbitrary. We show that the oscillator reaches a final equilibrium state via a damped evolution which is mostly exponential. We also show that some deviations from this exponential decay law (for very short and very long times) appear, which are intimately related with the presence of a lower bound of the energy.
In Sec. II the whole system (single oscillator plus the bath) is described and the Hamiltonian is introduced.
In Sec. III we diagonalize (in normal modes) the Hamiltonian. In the process of diagonalization some problems emerge, such as the lost of the discrete part of the energy spectrum . We can bypass these problems, if we use our definition of “regular” states. Then we can perform an analytical continuation of the spectral decomposition of the Hamiltonian, promoting the energy to complex values. To reach a successful interpretation of the results we require to generalize the definition given in paper to the model we are now studying. The mathematical bases of this generalization are shown in Appendix B.The reader who is not familiar with rigged Hilbert spaces and functional analysis can see Refs. and . This appendix also contains a rigorous mathematical understanding of the problem.
In Sec. IV mixed states and their evolution law are considered.
In Sec. V we deal with a very particular initial condition: An oscillator in a zero-temperature bath. We find the reduced density operator and show that the equilibrium state is reached. We accurately describe the time evolution of the system and estimate the Zeno and Khalfin effects for very short and long times, respectively. These are the deviations from an exact exponential decay law. Finally we show that our solution satisfy a Lindblad master equation when discarding these deviations from the exponential behavior.
Finally in Sec. VI some questions concerning irreversibility, already considered in papers and , are discussed.
We state our conclusions in Sec. VII.
Three mathematical appendices complete this work.
## II Particle-bath model
The system is a Brownian particle represented by a harmonic oscillator with natural frequency $`\mathrm{\Omega }.`$ It is well known that for a finite bath it is not possible to prove convergence in an equilibrium state in the limit $`t\mathrm{}`$ because of the existence of recurrences . However, for large systems these recurrence times become extremely huge and we can eliminate them by passing to the limit of an infinite continuous bath. Therefore, in this paper, we consider the oscillator in contact with a bath, already modeled by a continuous set of harmonic oscillators with natural frequencies $`\omega .`$ The coupling between the system and the bath is assumed to be linear with strength $`g(\omega )`$. The Hamiltonian for the composite system, in terms of creation and annihilation operators, is
$$H=\mathrm{\Omega }a^{}a+_0^{\mathrm{}}𝑑\omega \omega b_\omega ^{}b_\omega +\lambda _0^{\mathrm{}}𝑑\omega g(\omega )\left(a^{}b_\omega +b_\omega ^{}a\right).$$
(7)
The first term corresponds to the system, the second to the bath, and the third one corresponds to the interaction between them. In order that the Hamiltonian would be positive definite we must require that $`g(0)=0`$ and
$$\mathrm{\Omega }>\lambda ^2_0^{\mathrm{}}𝑑\omega \frac{g^2(\omega )}{\omega }.$$
(8)
This is an important condition which selects the kind of spectral densities appropriated to lead to an irreversible evolution. For example, the ohmic case which is frequently used in the literature must be disregarded, unless a cutoff is used. (Operators $`b_\omega `$ and $`b_\omega ^{}`$ are rigorously defined in Appendix A).
The Fock basis is the tensor product of the Fock basis of the isolated harmonic oscillator and those of the bath, namely
$$|n,\omega _1\mathrm{}\omega _m=|n|\omega _1\mathrm{}\omega _m,$$
(9)
where $`|\omega _1\mathrm{}\omega _m`$ represents a state with $`m`$ quanta in the bath, each one with frequency $`\omega _j`$ $`(j=1,\mathrm{},m)`$.
The total number of quanta is conserved allowing us to solve the problem by sectors (block diagonalization). The one-particle sector is referred as Friedrichs’ model and contains the relevant information that we need to compute physical quantities (see Sec. V).
## III Normal modes of the Hamiltonian and analytic continuation
The linearity in the coupling term of $`H`$ allows us to easily find a new set of uncoupled harmonic oscillators (normal modes), such that
$$I=_0^{\mathrm{}}𝑑\omega \stackrel{~}{b}_\omega ^{}\stackrel{~}{b}_\omega ,$$
(10)
$$H=_0^{\mathrm{}}𝑑\omega \omega \stackrel{~}{b}_\omega ^{}\stackrel{~}{b}_\omega ,$$
(11)
where
$$\stackrel{~}{b}_\omega =\xi _\omega a+_0^{\mathrm{}}𝑑\omega ^{}\mathrm{\Phi }_\omega (\omega ^{})b_\omega ^{}.$$
(12)
From a straightforward calculation , using the Heisenberg equations of motion, we obtain the coefficients of the unitary change of variables which diagonalize the Hamiltonian, precisely
$$\mathrm{\Phi }_\omega (\omega ^{})=\delta (\omega \omega ^{})+\frac{\lambda \xi _\omega g(\omega ^{})}{(\omega \omega ^{}+i\epsilon )}$$
(13)
and
$$\xi _\omega =\frac{\lambda g(\omega )}{\alpha (\omega +i\epsilon )},$$
(14)
where
$$\alpha (z)=z\mathrm{\Omega }\lambda ^2_0^{\mathrm{}}𝑑\omega \frac{g^2(\omega )}{z\omega }.$$
(15)
This function, which is the inverse of the reduced resolvent of $`H`$ in the one-particle sector, is not entire because it has a cut along the positive real axis corresponding to the continuous spectrum of the Hamiltonian. If $`\alpha (z)0`$ for all $`z𝐂`$, except for a possible real and negative $`\omega _0`$ such that $`\alpha (\omega _0)=0,`$ an isolated solution appears, which is non-analytic in $`\lambda `$ . We do not consider this case henceforth, since we are interested in analytic solutions satisfying condition (8).
If $`\alpha (z)=0`$ has no real solution it is not possible to find an operator $`\stackrel{~}{a},`$ such that $`\stackrel{~}{a}a`$ for $`\lambda 0.`$ In this case we have lost the particle number operator corresponding to the discrete part of the spectrum of $`H`$ and we do not have the correct form of $`H`$ when $`\lambda 0`$ . This problem can be solved promoting the energy (or frequency) $`\omega `$ to be a complex variable $`z`$. We define $`\beta (z)\left[\alpha (z)\right]^1`$. It can be proved that $`\beta (z)`$ has the same analytic structure than the one of the coefficient $`S(z)`$ of the scattering matrix . $`\beta (z)`$ is a meromorphic function on a double Riemann sheet with a cut along $`[0,+\mathrm{})`$. $`\beta _\pm (\omega )=\left[\alpha _\pm (\omega )\right]^1\beta (\omega \pm i\epsilon )`$ are defined on the upper and lower half-planes of the first Riemann sheet $`R_I`$ (physical sheet), and have meromorphic continuations to the lower and upper half-planes, respectively, in the second sheet $`R_{II}`$ (unphysical sheet). For simplicity we consider $`g(z)`$ such that the analytic extension of $`\beta _+(z)`$ into the second sheet has a simple pole $`z_0=\omega _0\frac{i}{2}\gamma `$ \[$`\gamma >0`$ and $`\alpha _+(z_0)=0`$\] in the lower half-plane. Also $`\beta _{}(z)`$ has a simple pole $`z_0^{}`$ on the upper plane in $`R_{II}.`$
We can now study the meaning of $`z_0`$. From the role played by $`z_0`$ in the evolution equation we know that $`(\mathrm{Im}z_0)^1=\gamma ^1`$ is the mean life time of the unstable state $`|1,v=a^{}|0,v`$ and $`(\mathrm{Re}z_0)`$ is the shift of the bare frequency $`\mathrm{\Omega }`$ \[see and also Eq. (59)\]. But, $`z_0`$ is the root of $`\alpha _+(z)`$ and from Eq. (15) we can estimate, up to the second order in $`\lambda ,`$
$$z_0=\mathrm{\Omega }+\lambda ^2\mathrm{P}\underset{0}{\overset{\mathrm{}}{}}𝑑\omega \frac{g^2(\omega )}{\mathrm{\Omega }\omega }i\pi \lambda ^2g^2(\mathrm{\Omega }),$$
(16)
where $`\mathrm{P}`$ denotes the Cauchy principal part of the integral.<sup>§</sup><sup>§</sup>§The principal part comes from the well known identity between distributions
$`{\displaystyle \frac{1}{x+i\epsilon }}=\mathrm{P}{\displaystyle \frac{1}{x}}i\pi \delta (x),x𝐑`$
The mean life of the unstable state and the shift frequency are given by
$$\gamma =2\pi \lambda ^2g^2(\mathrm{\Omega })$$
(17)
and
$$\delta \mathrm{\Omega }=\lambda ^2\mathrm{P}\underset{0}{\overset{\mathrm{}}{}}𝑑\omega \frac{g^2(\omega )}{\mathrm{\Omega }\omega }.$$
(18)
which are well-known expressions in the theory of unstable systems , usually derived from the Fermi golden rule.
Regarding the coupling function of the form $`g(\omega )\omega ^n`$ we find that the ohmic case ($`n=1`$) without cutoff does not satisfy the positivity condition (8). If we call $`\gamma _{1/2}`$ the coefficient for the subohmic case ($`0<n<1`$) and $`\gamma _2`$ the coefficient for supraohmic case ($`n>1`$), it is easy to prove that
$$\gamma _2\gamma _1\gamma _{1/2}.$$
(19)
Now we will find a generalized partition of the identity $`I`$ and a generalized spectral decomposition of $`H`$ that recovers the discrete part of the spectrum . In order to do this let $`\mathrm{\Gamma }`$ be the curve of Fig. 1. It lays on $`R_I`$ for $`\beta _{}(z)`$ and on $`R_{II}`$ for $`\beta _+(z)`$. We define the analytic function of $`z𝐂`$
$$\alpha _\mathrm{\Gamma }(z)=z\mathrm{\Omega }\lambda ^2_\mathrm{\Gamma }𝑑z^{}\frac{g^2(z^{})}{zz^{}},$$
(20)
which generalizes Eq. (15).
To find the partition of the identity and a expansion of $`H`$ we will use in some adequate analyticity properties.Properties of this kind were already introduced in previous works . Thus we will define a space $`\mathrm{\Phi }_+_+`$ of states $`|\phi `$ such that the function $`0,\omega _1,\mathrm{},\omega _n|\phi =\phi _0(\omega _1,\mathrm{},\omega _n)`$ would have an analytic continuation, for each variable $`\omega _i`$ ($`1in)`$ to a region that include the singularity $`z_0.`$ This space $`\mathrm{\Phi }_+`$ would be our space of “regular,” “admissible” or “physical” states. Precisely, generalizing what we have done in paper , we will chose $`\mathrm{\Phi }_+`$ such that its states would satisfy Eq. (B5) of Appendix B. Analogously, the space $`\mathrm{\Phi }_{}`$ of the “unphysical” time-inverted states would satisfy Eq. (B6) of that appendix. With this choice the analytic continuation that we will perform has a rigorous meaning, since the operators act in a space $`\mathrm{\Phi }_+`$which endowed with adequate analytic properties. The demonstration of this fact is a mathematical problem, which is considered in Appendices B and C, where we generalize previous results. Then if $`|\phi \mathrm{\Phi }_+`$ and $`|\psi \mathrm{\Phi }_{}`$, from Eqs. (10) and (11) using appendix B and following the similar demonstration of paper , it can be proved that
$`\psi |\phi `$ $`=`$ $`\psi |(\stackrel{~}{a}^{()}\stackrel{~}{a}^{()}+{\displaystyle _\mathrm{\Gamma }}𝑑z\stackrel{~}{b}_z^{()}\stackrel{~}{b}_z^{()})|\phi ,`$ (21)
$$\psi |H|\phi =\psi |(z_0\stackrel{~}{a}^{()}\stackrel{~}{a}^{()}+_\mathrm{\Gamma }𝑑zz\stackrel{~}{b}_z^{()}\stackrel{~}{b}_z^{()})|\phi .$$
(23)
The residue at $`z_0`$ contributes to the first terms of the r.h.s. of these generalized partition of the identity and spectral decomposition of $`H,`$ as in paper , and, in a weak sense, the two previous equations can be written as
$$I=\stackrel{~}{a}^{()}\stackrel{~}{a}^{()}+_\mathrm{\Gamma }𝑑z\stackrel{~}{b}_z^{()}\stackrel{~}{b}_z^{()}$$
(24)
$$H=z_0\stackrel{~}{a}^{()}\stackrel{~}{a}^{()}+_\mathrm{\Gamma }𝑑zz\stackrel{~}{b}_z^{()}\stackrel{~}{b}_z^{()},$$
(25)
The creation and annihilation operators in all these equations reads
$`\stackrel{~}{a}^{()}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\alpha _+^{}(z_0)}}}\left[a+\lambda {\displaystyle _0^{\mathrm{}}}𝑑z{\displaystyle \frac{g(z)}{[z_0z]_+}}b_z\right],`$ (26)
$`\stackrel{~}{a}^{()}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\alpha _+^{}(z_0)}}}\left[a^{}+\lambda {\displaystyle _0^{\mathrm{}}}𝑑z{\displaystyle \frac{g(z)}{[z_0z]_+}}b_z^{}\right],`$ (28)
and
$`\stackrel{~}{b}_\omega ^{()}`$ $`=`$ $`b_\omega +{\displaystyle \frac{\lambda g(\omega )}{\eta _+(\omega )}}\left[a+\lambda {\displaystyle _0^{\mathrm{}}}𝑑\omega ^{}{\displaystyle \frac{g(\omega ^{})}{\omega \omega ^{}+i\epsilon }}b_\omega ^{}\right],`$ (29)
$`\stackrel{~}{b}_\omega ^{()}`$ $`=`$ $`b_\omega ^{}+{\displaystyle \frac{\lambda g(\omega )}{\alpha _{}(\omega )}}\left[a^{}+\lambda {\displaystyle _0^{\mathrm{}}}𝑑\omega ^{}{\displaystyle \frac{g(\omega ^{})}{\omega \omega ^{}i\epsilon }}b_\omega ^{}^{}\right].`$ (31)
The distribution $`\frac{1}{[z_0z]_+}`$ means
$$_0^{\mathrm{}}\frac{f(\omega )}{[z_0\omega ]_+}𝑑\omega =_\mathrm{\Gamma }\frac{f(z)}{z_0z}𝑑z=_0^{\mathrm{}}\frac{f(\omega )}{z_0\omega }𝑑\omega +2\pi if(z_0),$$
(32)
for every well-behaved analytical function $`f(z).`$ Observe that $`\stackrel{~}{b}_\omega ^{()}`$ does not change if we replace $`_\mathrm{\Gamma }`$ by $`_0^{\mathrm{}}`$ because, in this case, no pole is crossed, since $`\alpha _{}(z)`$ has not poles (in $`S_I`$). Nevertheless $`\stackrel{~}{b}_\omega ^{()}`$ does change because the pole is crossed when we modify the integration contour. We have shown this fact explicitly putting $`\eta _+(\omega )`$ in place of $`\alpha _+(\omega ),`$ where
$$\frac{1}{\eta _+(\omega )}=\frac{1}{\alpha _+(\omega )}+2\pi i\frac{\delta (zz_0)}{\alpha _+^{}(z_0)},$$
(33)
being $`\delta (zz_0)`$ the extension of the Dirac delta defined à la Gel’fand and Shilov . As a consequence of these facts $`\stackrel{~}{a}^{()}\stackrel{~}{a}^{()}`$ and $`\stackrel{~}{b}_\omega ^{()}\stackrel{~}{b}_\omega ^{()}`$. The star operation corresponds to the analytic generalization of the complex conjugation, which, acting on an analytic function $`f(z)`$, is defined byIt corresponds to the symbol # of paper . Here we follow the notation $``$ of paper in which we have also studied this model.
$$f^{}(z)=\left[f(z^{})\right]^{}.$$
(34)
From Eqs. (LABEL:aaa) and (LABEL:be) we see that we have four annihilation operators due to the presence of complex eigenvalues of $`H`$ with the corresponding doubling of solutions, since we have a pair of complex conjugate values. They are generalized eigenvalues of the two analytic continuations of $`H`$ into the lower ($``$) \[upper ($`+`$)\] -complex plane. These operators are
$`\stackrel{~}{a}^{()};\stackrel{~}{a}^{(+)};\stackrel{~}{b}_\omega ^{()};\stackrel{~}{b}_\omega ^{(+)}`$
The vacuum is the state annihilated by any annihilation operator. The Bogolubov transformation of Eqs. (LABEL:aaa) and (LABEL:be) does not mix creation and annihilation operators, therefore the vacuum just defined is actually the same state defined as the vacuum of the noninteracting system+bath. So from Eq. (9) the vacuum is the state $`|0|v|0,v`$, where $`|v`$ is the vacuum of the bath.
The corresponding creation operators are
$`\stackrel{~}{a}^{()};\stackrel{~}{a}^{(+)};\stackrel{~}{b}_\omega ^{()};\stackrel{~}{b}_\omega ^{(+)}`$
Starting from the common vacuum, by applying successively the operators $`\stackrel{~}{a}^{()}`$ and $`\stackrel{~}{b}_\omega ^{()}`$, the Fock basis $`\left\{|\stackrel{~}{}\right\}`$ is built, and with $`\stackrel{~}{a}^{(+)}`$ and $`\stackrel{~}{b}_\omega ^{(+)}`$ we build up the Fock basis $`\{|\overline{}`$ $`\}`$. In the case of $`\stackrel{~}{a}^{()}`$ and $`\stackrel{~}{a}^{(+)}`$ the corresponding vectors in the Fock bases of the one-particle sector are generalized eigenvectors of $`H`$ with purely complex eigenvalues. They represent unstable states, i.e. $`\stackrel{~}{a}^{()}|\stackrel{~}{0,v}=|\stackrel{~}{1,v}`$ is a one-particle generalized eigenvector of $`H`$ corresponding to a complex eigenvalue $`z_0`$ and $`\stackrel{~}{a}^{(+)}|\overline{0,v}=|\overline{1,v}`$ is a one-particle generalized eigenvector of $`H`$ corresponding to a complex eigenvalue $`z_0^{}.`$ In this way we are able to develop a second quantized version of the theory of unstable states .
Now we have two different number of quanta operators, $`\stackrel{~}{N}^{()}=_0^{\mathrm{}}𝑑\omega \stackrel{~}{b}_\omega ^{()}\stackrel{~}{b}_\omega ^{()}`$ and $`\stackrel{~}{N}^{(+)}=_0^{\mathrm{}}𝑑\omega \stackrel{~}{b}_\omega ^{(+)}\stackrel{~}{b}_\omega ^{(+)},`$ which are not Hermitian. So two different Fock bases can be built satisfying
$`\stackrel{~}{N}^{()}|\stackrel{~}{n,\omega _1\mathrm{}\omega _m}`$ $`=`$ $`m|\stackrel{~}{n,\omega _1\mathrm{}\omega _m},`$ (35)
$`\stackrel{~}{N}^{(+)}|\overline{n,\omega _1\mathrm{}\omega _m}`$ $`=`$ $`m|\overline{n,\omega _1\mathrm{}\omega _m}.`$ (37)
The spectral decomposition of the Hamiltonian reads
$`H^{()}`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}\mathrm{}{\displaystyle _0^{\mathrm{}}}𝑑\omega _1\mathrm{}𝑑\omega _m\left(z_0n+\omega _1+\mathrm{}+\omega _m\right)`$ (39)
$`\times |\stackrel{~}{n,\omega _1\mathrm{}\omega _m}\overline{n,\omega _1\mathrm{}\omega _m}|,`$
which acts on the right of the Fock space generated by basis $`\left\{|\stackrel{~}{}\right\}`$. But the “same” Hamiltonian can also be written in the following way (using the other analytical continuation, in which case it is evident that the next equation is only weak, and it has a precise meaning operating between $`|\phi \mathrm{\Phi }_{}`$ and $`|\psi \mathrm{\Phi }_+`$)
$`H^{(+)}`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}\mathrm{}{\displaystyle _0^{\mathrm{}}}𝑑\omega _1\mathrm{}𝑑\omega _m\left(z_0^{}n+\omega _1+\mathrm{}+\omega _m\right)`$ (41)
$`\times |\overline{n,\omega _1\mathrm{}\omega _m}\stackrel{~}{n,\omega _1\mathrm{}\omega _m}|,`$
which acts on the right of the Fock space generated by basis $`\left\{|\overline{}\right\}`$. In the same way, the identity reads
$`I^{()}`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}\mathrm{}{\displaystyle _0^{\mathrm{}}}𝑑\omega _1\mathrm{}𝑑\omega _m|\stackrel{~}{n,\omega _1\mathrm{}\omega _m}\overline{n,\omega _1\mathrm{}\omega _m}|,`$ (42)
$`I^{(+)}`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}\mathrm{}{\displaystyle _0^{\mathrm{}}}𝑑\omega _1\mathrm{}𝑑\omega _m|\overline{n,\omega _1\mathrm{}\omega _m}\stackrel{~}{n,\omega _1\mathrm{}\omega _m}|.`$ (44)
As the eigenvalues of Eqs. (39) and (41) are complex, in order to deal with unstable states we must find an adequate mathematical structure beyond the Hilbert space. In fact, these states are generalized states. Thus, in the Appendix B we see that kets $`|\stackrel{~}{}`$ and $`|\overline{}`$ are well defined in a rigged Hilbert space formalism, i.e. they must be thought as antilinear functionals acting on test spaces $`\mathrm{\Phi }_\pm `$ and, as elements of a vector space, they belong to the duals of $`\mathrm{\Phi }_\pm ,`$ symbolized by $`\mathrm{\Phi }_\pm ^\times .`$ They define a double Gel’fand triplet structure $`\mathrm{\Phi }_\pm _\pm \mathrm{\Phi }_\pm ^\times `$ .
What about energy conservation? The trouble emerges because the eigenvalues of the Hamiltonian are now complex; thus some states decay in time (e.g. vectors of $`\mathrm{\Phi }_+^\times `$ which vanish for long times, see Sec. V). Since some states vanish we may ask ourselves how the conservation of energy can be possible. The answer is that energy is conserved anyhow. In order to demonstrate this fact we will calculate the mean value of the Hamiltonian in a state $`|\phi (t)=\underset{nm}{}\underset{0}{\overset{\mathrm{}}{}}\mathrm{}\underset{0}{\overset{\mathrm{}}{}}𝑑\omega _1\mathrm{}𝑑\omega _mc_n(\omega _1\mathrm{}\omega _m)|\stackrel{~}{n,\omega _1\mathrm{}\omega _m}`$, precisely
$`E`$ $`=`$ $`\phi (t)|H^{()}|\phi (t)`$ (45)
$`=`$ $`{\displaystyle \underset{nn^{}}{}}{\displaystyle \underset{mm^{}}{}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑\omega _1\mathrm{}𝑑\omega _m𝑑\omega _1^{}\mathrm{}𝑑\omega _m^{}^{}c_n^{}(\omega _1\mathrm{}\omega _m)c_n^{}(\omega _1^{}\mathrm{}\omega _m^{}^{})`$ (48)
$`\times e^{i(z_0n^{}+\omega _1^{}+\mathrm{}+\omega _m^{}^{})t}e^{i(z_0^{}n+\omega _1+\mathrm{}+\omega _m)t}(z_0n+\omega _1+\mathrm{}+\omega _m)`$
$`\times \stackrel{~}{n,\omega _1,\mathrm{}\omega _m}|\stackrel{~}{n^{},\omega _1^{}\mathrm{}\omega _m^{}^{}}.`$
Taking into account the orthogonality relations (LABEL:conmu1) and (LABEL:conmu4) shown in Appendix B, and that $`z_0=\omega _0i\frac{\gamma }{2}`$, $`\gamma >0`$, Eq. (48) reduces to
$$E=\underset{m}{}\underset{0}{\overset{\mathrm{}}{}}\mathrm{}\underset{0}{\overset{\mathrm{}}{}}𝑑\omega _1\mathrm{}𝑑\omega _m\left|c_0(\omega _1\mathrm{}\omega _m)\right|^2(\omega _1+\mathrm{}+\omega _m),$$
(49)
which is time independent. Thus energy is conserved. Conservation of the norm and the number of particles can also be demonstrated in an analogous way changing $`H`$ by $`I`$.
Finally we can observe that Eqs. (LABEL:conmu1) and (LABEL:conmu4) show that the generalized eigenvectors have null norm and energy (with the exception of those with $`n=0`$) . In the literature they are called Gamow vectors, they are generalized states, and they represent just idealized mathematical states (see Appendix B), as is the case of the plane waves.
## IV Mixed states: Its evolution
A general pure state belonging to $`\mathrm{\Phi }_+`$ (see Appendix B) can be written as
$$|\mathrm{\Psi }=\underset{n}{}\underset{m}{}_0^{\mathrm{}}\mathrm{}_0^{\mathrm{}}𝑑\omega _1\mathrm{}𝑑\omega _mc_n(\omega _1,\mathrm{},\omega _m)\stackrel{~}{|n,\omega _1\mathrm{}\omega _m},$$
(50)
where $`\stackrel{~}{|n,\omega _1\mathrm{}\omega _n}\mathrm{\Phi }_+^\times `$, and the most general density<sup>\**</sup><sup>\**</sup>\**More accurately, we would say, that it is the most general possible decaying density matrix, as we will see. matrix can be written as
$`\rho `$ $`=`$ $`{\displaystyle \underset{nn^{}}{}}{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}\mathrm{}{\displaystyle _0^{\mathrm{}}}𝑑\omega _1\mathrm{}𝑑\omega _m𝑑\omega _1^{}\mathrm{}𝑑\omega _m^{}^{}c_{nn^{}}(\omega _1,\mathrm{},\omega _m,\omega _1^{},\mathrm{},\omega _m^{}^{})`$ (52)
$`\times \stackrel{~}{|n,\omega _1\mathrm{}\omega _m}\stackrel{~}{n,\omega _1^{},\mathrm{},\omega _m^{}^{}|}.`$
If this $`\rho `$ is the initial state $`\rho =\rho (0)`$, the evolution law of $`\rho (t)`$ reads
$$\rho (t)=e^{iH^{()}t}\rho (0)e^{iH^{(+)}t}.$$
(53)
As $`H`$ is only self-adjoint in a generalized way <sup>††</sup><sup>††</sup>††Recall that $`H`$ is self-adjoint in the Hilbert space, where $`=^\times `$, but in the generalized Hilbert space this property essentially becomes Eq. (LABEL:footnote). $`H^{()}`$ acts in a different way than $`H^{(+)}=H^{()}`$ and there are right and left eigenvalues,
$`H^{()}\stackrel{~}{|n,\omega _1\mathrm{}\omega _m}`$ $`=`$ $`(z_0n+\omega _1+\mathrm{}+\omega _m)\stackrel{~}{|n,\omega _1\mathrm{}\omega _m},`$ (54)
$`\stackrel{~}{n,\omega _1\mathrm{}\omega _m}|H^{(+)}`$ $`=`$ $`(z_0^{}n+\omega _1+\mathrm{}+\omega _m)\stackrel{~}{n,\omega _1\mathrm{}\omega _m}|.`$ (56)
Then we have
$`\rho (t)`$ $`=`$ $`{\displaystyle \underset{nn^{}}{}}e^{\frac{\gamma }{2}(n+n^{})t}e^{i\omega _0(nn^{})t}{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}\mathrm{}{\displaystyle _0^{\mathrm{}}}𝑑\omega _1\mathrm{}𝑑\omega _m𝑑\omega _1^{}\mathrm{}𝑑\omega _m^{}^{}`$ (59)
$`\times e^{i(\omega _1+\mathrm{}+\omega _m)t}e^{i(\omega _1^{}+\mathrm{}+\omega _m^{}^{})t}c_{nn^{}}(\omega _1,\mathrm{},\omega _m,\omega _1^{},\mathrm{},\omega _m^{}^{})`$
$`\times \stackrel{~}{|n,\omega _1\mathrm{}\omega _m}\stackrel{~}{n^{},\omega _1^{}\mathrm{}\omega _m^{}^{}|}.`$
For an arbitrary initial state $`\rho (0)`$ a time dependent asymptotic ($`t+\mathrm{}`$) state is reached. The explanation of this fact is simple. The modes of the bath are independent of each other \[see Eq. (7)\], and so we cannot expect that the bath reaches equilibrium (cf. Ref. ). Thus
$`\rho (t)`$ $``$ $`\rho _{}(t)={\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}\mathrm{}{\displaystyle _0^{\mathrm{}}}𝑑\omega _1\mathrm{}𝑑\omega _m𝑑\omega _1^{}\mathrm{}𝑑\omega _m^{}^{}e^{i(\omega _1+\mathrm{}+\omega _m)t}e^{i(\omega _1^{}+\mathrm{}+\omega _m^{}^{})t}`$ (61)
$`\times c_{00}(\omega _1,\mathrm{},\omega _m,\omega _1^{},\mathrm{},\omega _m^{}^{})\stackrel{~}{|0,\omega _1\mathrm{}\omega _m}\stackrel{~}{0,\omega _1^{}\mathrm{}\omega _m^{}^{}}|.`$
For completeness we also write the evolution equation for the density operator,
$`{\displaystyle \frac{d\rho (t)}{dt}}`$ $`=`$ $`i{\displaystyle \underset{nn^{}}{}}e^{\frac{\gamma }{2}(n+n^{})t}e^{i\omega _0(nn^{})t}{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}\mathrm{}{\displaystyle _0^{\mathrm{}}}𝑑\omega _1\mathrm{}𝑑\omega _m𝑑\omega _1^{}\mathrm{}𝑑\omega _m^{}^{}`$ (64)
$`\times e^{i(\omega _1+\mathrm{}+\omega _m)t}e^{i(\omega _1^{}+\mathrm{}+\omega _m^{}^{})t}(z_0n+\omega _1+\mathrm{}+\omega _mz_0^{}n^{}\omega _1^{}\mathrm{}\omega _m^{}^{})`$
$`\times c_{nn^{}}(\omega _1,\mathrm{},\omega _m,\omega _1,\mathrm{},\omega _m^{})\stackrel{~}{|n,\omega _1\mathrm{}\omega _m}\stackrel{~}{n^{},\omega _1^{}\mathrm{}\omega _m^{}^{}}|,`$
which is clearly equal to
$$\frac{d\rho (t)}{dt}=i\left(H^{()}\rho \rho H^{(+)}\right)=iL\rho ,$$
(65)
where $`L`$ is the generalized Liouvillian operator . So we see that the density operator follows an evolution described by a generalized Liouville-von Neumann equation.
In spite of the result obtained in Eq. (61) in Sec. V we show that the reduced density operator $`\rho _r,`$ which is obtained by taking the partial trace with respect to the environment modes, reaches equilibrium, namely a time independent state. An equivalent way to find the equilibrium state, closer to the spirit of our formalism, is to use a particular space of observables, as in paper .
## V Reduced density operator
As an illustration of the formalism we consider a simple example where the initial state is a very particular state of the composed system,
$$\rho (0)=\rho _S(0)\rho _E(0),$$
(66)
where
$$\rho _S(0)=c_{11}|11\left|+c_{10}\right|10\left|+c_{01}\right|01\left|+c_{00}\right|00|$$
(67)
is the initial state of the discrete oscillator, or the initial reduced density operator, (with $`c_{11},c_{00}0`$, $`c_{11}+c_{00}=1`$ and $`c_{10}=c_{01}^{}`$) and
$$\rho _E(0)=|vv|$$
(68)
is the initial state of the bath, which does not have any quantum, namely, it is in the ground state. This corresponds to a bath at zero temperature $`T=0`$ (in the zero and one-particle sector). The main features at any $`T`$ can be reproduced but we begin with this example because the mathematical computations are easier (recall that this model can be decomposed in sectors of constant number of quanta). Also our initial conditions are such that there is no correlation between the oscillator and the bath.
Our aim is to find the time dependence of the reduced matrix elements. It is derived from the time evolution of the density operator
$$\rho (t)=e^{iH^{()}t}\rho (0)e^{iH^{(+)}t}=e^{iH^{()}t}I^{()}\rho (0)I^{(+)}e^{iH^{(+)}t},$$
(69)
where $`I^{()}`$ and $`I^{(+)}`$ are the identities in spaces $`\mathrm{\Phi }_+`$ and $`\mathrm{\Phi }_{}`$ respectively<sup>‡‡</sup><sup>‡‡</sup>‡‡The difference in the conventions with respect to the use of $`+`$ and $``$ is the following. In operators $`+`$ and $``$ are related with the analytic continuations for $`\pm i\epsilon ,`$ while $`+`$ and $``$ in spaces are associated with the time evolution, which is only well defined for positive or negative times, respectively. \[see Eqs. (LABEL:Id)\].
$`\rho (t)`$ $`=`$ $`\left(e^{iz_0t}|\stackrel{~}{1,v}\overline{1,v}\left|+{\displaystyle _0^{\mathrm{}}}𝑑\omega e^{i\omega t}\right|\stackrel{~}{0,\omega }\overline{0,\omega }|\right)\rho (0)`$ (71)
$`\times \left(e^{iz_0^{}t}|\overline{1,v}\stackrel{~}{1,v}\left|+{\displaystyle _0^{\mathrm{}}}𝑑\omega ^{}e^{i\omega ^{}t}\right|\overline{0,\omega ^{}}\stackrel{~}{0,\omega ^{}}|\right).`$
We have only considered the terms of the identity that correspond to zero-particle and one-particle subspaces, since, from the conservation of the number of quanta, there is no contribution of other terms. We emphasize that no approximations were carried out up to here.
Once we have the time evolution of the density operator, the following step is to get the reduced density operator, tracing over the basis corresponding to the environment,
$`\rho _r(t)`$ $`=`$ $`tr_E\rho (t)={\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}\mathrm{}{\displaystyle _0^{\mathrm{}}}𝑑\omega _1\mathrm{}𝑑\omega _m𝑑\omega _1^{}\mathrm{}𝑑\omega _m^{}`$ (73)
$`\times \omega _1\mathrm{}\omega _m\left|\rho (t)\right|\omega _1^{}\mathrm{}\omega _m^{}\delta (\omega _1\omega _1^{})\mathrm{}\delta (\omega _m\omega _m^{}),`$
where the $`m=0`$ term means $`v\left|\rho (t)\right|v.`$
As we have said, the contribution of terms $`m=2,3,\mathrm{}`$ vanishes. Therefore, using the relations between “new” and “old” bases \[Eqs. (LABEL:aaa) and (LABEL:be)\] and the conservation of trace, we obtain a positive definite reduced density operator:
$`\rho _r(t)=c_{11}P(t)|11\left|+c_{10}\mathrm{\Delta }_0(t)\right|10|+`$
$$c_{01}\mathrm{\Delta }_0^{}(t)|01\left|+\left\{c_{00}+c_{11}\left[1P(t)\right]\right\}\right|00|,$$
(74)
where
$$\mathrm{\Delta }_0(t)=\frac{e^{iz_0t}}{\alpha _+^{}(z_0)}+_0^{\mathrm{}}𝑑\omega e^{i\omega t}\frac{\lambda ^2g^2(\omega )}{\eta _+(\omega )\alpha _{}(\omega )},$$
(75)
and $`P(t)=|\mathrm{\Delta }_0(t)|^2`$ is the survival probability of the state with only one quantum in the discrete part.
We can write $`P(t)`$ as the sum of four terms where the first one, $`\frac{e^{\gamma t}}{|\alpha _+^{}(z_0)|^2},`$ shows an exact exponential behavior. Expanding $`|\alpha _+^{}(z_0)|^2`$ as $`1+O(\lambda ^2),`$ we split the probability into two terms, one containing the purely exponential contribution and the other that we call “background,” giving rise to derivations from that purely exponential decay law, so that
$$P(t)=e^{\gamma t}+\mathrm{background}.$$
(76)
If we take a time neither very short nor very long, the background will be smaller than the purely exponential term (for $`\lambda 1`$) and it can be neglected, which leads to an exponential decay-law. This is not true for short times since $`\frac{dP}{dt}(0)=0`$, which leads to the so-called Zeno effect . For very long times the exponential term will decay faster than the background will do, which is known as Khalfin effect . We can force $`P(t)`$ to have an exponential appearance by defining the decay rate $`\mathrm{\Gamma }(t)`$ to be time dependent, namely
$$P(t)e^{\mathrm{\Gamma }(t)t},$$
(77)
with
$`\mathrm{\Gamma }(t)=\gamma {\displaystyle \frac{1}{t}}\mathrm{ln}\left(1+e^{\gamma t}\mathrm{background}\right).`$
Obviously for an intermediate time the background can be neglected and $`\mathrm{\Gamma }(t)\gamma .`$ The main restriction, imposed by the Zeno period, is
$$\frac{dP}{dt}(0)=\mathrm{\Gamma }(0)=0.$$
(78)
For a very long time the decay probability has also a non-exponential contribution as a consequence of the semi-finiteness of the energy spectrum. From papers we know that the survival amplitude goes to zero as $`t`$ goes to infinity as a consequence of the Riemann-Lebesgue theorem. Then the behavior of $`\mathrm{\Delta }_0`$ depends on the small-frequency behavior of $`g^2(\omega ).`$ For small $`\omega ,`$ $`\alpha _+(\omega )\omega _0,`$ where condition (8) was also used. Then the form of $`\mathrm{\Delta }_0`$ depends essentially on $`g(\omega )`$ for large $`t.`$ As an example we consider the case where $`g^2(\omega )\omega ^n\mathrm{exp}\left(\frac{\omega ^2}{\mathrm{\Lambda }^2}\right),`$ where $`\mathrm{\Lambda }`$ is a cutoff \[see paragraph before Eq. (19)\]. By evaluating the survival amplitude we have
$`\mathrm{\Delta }_0(t)=\lambda ^2{\displaystyle _0^{\mathrm{}}}𝑑\omega {\displaystyle \frac{g^2(\omega )}{\left|\alpha _+(\omega )\right|^2}}e^{i\omega t}{\displaystyle _0^{1/t}}𝑑\omega \omega ^ne^{i\omega t}\mathrm{exp}\left({\displaystyle \frac{\omega ^2}{\mathrm{\Lambda }^2}}\right),`$
where the contribution of high-frequency terms is negligible. Performing the change of variables $`\omega t=x,`$ we obtain
$$\mathrm{\Delta }_0(t)t^{(n+1)}_0^1𝑑xx^ne^{ix}\left(1\frac{x^2}{\mathrm{\Lambda }^2t^2}+\frac{x^4}{\mathrm{\Lambda }^4t^4}+\mathrm{}\right).$$
(79)
We can see that the survival amplitude merges into an algebraic long-time tail. The first relevant contribution behaves as $`t^{n1}`$. As a consequence the decay rate for long times must behave as
$$\mathrm{\Gamma }(t)\frac{\mathrm{ln}t}{t}.$$
(80)
The behavior at short times and intermediate times coincides with those obtained in Ref. . In Fig. 2 we show the qualitative behavior of $`P(t).`$ Zeno’s time, $`t_Z,`$ and Khalfin’s time, $`t_K,`$ are not in scale in the picture in order to show the three different contributions to the decay probability.
Eq. (74) is the exact solution to the proposed problem, without taking any approximation. The first, second, and third terms will vanish for $`t\mathrm{}`$; in fact $`P(t\mathrm{})=0`$ and the same happens for $`\mathrm{\Delta }_0(t)`$ (recall that $`P=|\mathrm{\Delta }_0|^2`$). The first term of (75) has the factor $`e^{\frac{\gamma }{2}t}`$ and the second one will tend to zero because of the Riemann-Lebesgue theorem. The probability of having the vacuum will grow in time. This means that all quanta in the discrete spectrum, except the ground state, decay into the continuum. So we find the equilibrium reduced density operator
$$\rho _{}=(c_{00}+c_{11})|00|=tr\rho |00|=|00|.$$
(81)
As expected, the equilibrium state is the vacuum, namely for $`t\mathrm{}`$ there are no quanta in the discrete spectrum, because the initial quantum has decayed into the bath (the discrete oscillator has spread its energy over the infinite oscillators of the bath with a distribution centered at the shifted frequency $`\omega _0`$) . This means that the discrete harmonic oscillator has thermalized at $`T=0`$. A similar result was recently obtained in Ref. .
In order to check the compatibility of the solution found in Eq. (74) we first briefly sketch the main points of the derivation we have done. We have obtained the exact solution of the Liouville equation. As a particular case, we have considered an initial condition restricted to the zero- and one-particle sectors and we have traced this solution over the environment modes. In that case, the survival probability $`P(t)`$ in Eq. (74) can be approximated by an exponential behavior, when the background contribution is neglected. Now the solution for $`\rho _r(t)`$ obtained through this approximation can be derived from a master equation of a Lindblad’s form, where the Lindblad generator is proportional to the destruction operator $`a,`$ since in our case we are only considering a zero-temperature bath in the case of a damping motion caused by friction :
$`\dot{\rho }_r(t)=i\mathrm{\Omega }_0[a^{}a,\rho _r]+{\displaystyle \frac{\gamma }{2}}\left(2a\rho _ra^{}a^{}a\rho _r\rho _ra^{}a\right).`$
The probability of finding $`n`$ quanta follows a Pauli master equation:
$$\frac{}{t}n\left|\rho _r\right|n=\gamma \left[\left(n+1\right)n+1\left|\rho _r\right|n+1nn\left|\rho _r\right|n\right].$$
(82)
It is easy to see that $`\rho _r`$ of Eq. (74), when the background is neglected, is solution of the Pauli equation (82).<sup>\**</sup><sup>\**</sup>\**Moreover $`\rho _r(t)`$ is possitive definites since this condition is equivalent to $`\rho _{11}\rho _{00}|\rho _{10}|^2`$ ($`\rho _{ij}=i\left|\rho _r\right|j`$) which is obviously satisfied by $`\rho _r(t)`$ provided it is satisfied by $`\rho _r(0)`$.
The results listed above are well known, but shown that our formalism works as the usual well established theories on the subject.
## VI Semigroups, Wigner time-inversion and irreversibility
One of the features of the system we are studying is its irreversible evolution; the system reaches the equilibrium at the far future, and of course the inverse evolution is not possible anymore. These properties were already found in Refs. and here they are briefly reviewed.
The presence of a time-asymmetric behavior can be shown in two different ways: As the splitting of the usual evolution group in two semigroups or as the impossibility to make a time inversion. We consider that the second one is the most eloquent.
### A Semigroups
Physical states can be chosen to be in test space $`\mathrm{\Phi }_+,`$ as we have done (or, with a simple and physically irrelevant change of convention in space $`\mathrm{\Phi }_{}`$), while corresponding generalized eigenvectors are in its dual $`\mathrm{\Phi }_+^\times `$ (or $`\mathrm{\Phi }_{}^\times `$). The proof is simple. From the Paley-Wiener theorem the following lemma can be deduced:
If $`f(\omega )H_+^2`$, $`e^{i\omega t}f(\omega )H_+^2`$ only for positive times. Similarly if $`f(\omega )H_{}^2`$, $`e^{i\omega t}f(\omega )H_{}^2`$ only for negative times.
The asymmetry in Hardy spaces can be immediately seen from the Fourier transform representation. From the Paley-Wiener theorem it is known that Hardy class functions from above can be represented as
$`f(\omega )={\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle _0^{\mathrm{}}}𝑑se^{i\omega s}\widehat{f}(s),`$
where
$`\widehat{f}(s)={\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega e^{i\omega s}f(\omega ).`$
The Fourier transform of Hardy class functions $`H_+^2`$ is in the space of square integrable functions with support on the positive real axis; and as the Fourier transform of $`e^{i\omega t}f(\omega )`$ is given by
$`\{e^{i\omega t}f(\omega )\}=\widehat{f}(st),`$
i.e. a function with support on $`[0,\mathrm{})`$ is transformed into a function with support on $`[t,\mathrm{})`$. For a negative time this last function has no longer support on $`[0,\mathrm{})`$ and therefore $`e^{i\omega t}f(\omega )`$ does not belong to $`H_+^2`$. Analogously it can be proved the same property for the Hardy class $`H_{}^2`$.
To simplify, we analyze the one-particle case; generalization to $`n`$-particle states is straightforward. Let $`\varphi (\omega )`$ be a function in $`\theta \left[𝒮H_\pm ^2\right]`$ such that, as a consequence of previous lemma, $`e^{i\omega t}\varphi (\omega )`$ will also be in the same space for $`t>0`$ only. As $`\varphi (\omega )=\omega |\varphi `$ then we can write in Dirac notation
$$e^{i\omega t}\varphi (\omega )=e^{i\omega t}\omega |\varphi =\omega |e^{iH^{()}t}|\varphi ,$$
(83)
and taking into account Eq. (B7) we can state that if $`|\varphi \mathrm{\Phi }_+`$ then $`e^{iH^{()}t}|\varphi \mathrm{\Phi }_+`$ only for $`t>0`$. In the same way, if $`|\varphi \mathrm{\Phi }_{}`$ then $`e^{iH^{(+)}t}|\varphi \mathrm{\Phi }_{}`$ for $`t<0`$. Then if we postulate that $`\mathrm{\Phi }_+`$ is the space of physical states and that the physical evolution brings physical states into physical states, it turns out that this evolution only take place in the period $`t>0,`$ so irreversibility naturally appears <sup>\*†</sup><sup>\*†</sup>\*†Moreover, states in $`\mathrm{\Phi }_+`$ are linear combination of generalized vectors $`|\stackrel{~}{n,\omega _1\mathrm{}\omega _m}`$, and these vectors evolve as
$$e^{iH^{()}t}|\stackrel{~}{n,\omega _1\mathrm{}\omega _m}=e^{i(\omega _0n+\omega _1+\mathrm{}+\omega _m)t}e^{\frac{\gamma }{2}nt}|\stackrel{~}{n,\omega _1\mathrm{}\omega _m}.$$
(84) Therefore, with the exception of $`n=0`$, they decay towards increasing $`t`$. Then we say that physical vectors in $`\mathrm{\Phi }_+`$ decay towards positive time (except the states belonging to $`\mathrm{\Phi }_+\mathrm{\Phi }_{},`$ as the vacuum state $`|0|v`$, which does not decay). Analogously
$$e^{iH^{(+)}t}|\overline{n,\omega _1\mathrm{}\omega _m}=e^{i(\omega _0n+\omega _1+\mathrm{}+\omega _m)t}e^{\frac{\gamma }{2}nt}|\overline{n,\omega _1\mathrm{}\omega _m},$$
(85) and we can say that the unphysical states in $`\mathrm{\Phi }_{}`$ decay towards negative times..
### B The Wigner time-inversion
The Wigner time-inversion operator acts in a real representation as
$$K\phi (x)=\phi ^{}(x).$$
(86)
\[cf. Eq. (3)\]. Then, as the complex conjugate of the functions of $`H_+^2`$ are the functions of $`H_{}^2`$ we have that
$`K`$ $`:`$ $`\mathrm{\Phi }_+\mathrm{\Phi }_{}\mathrm{\Phi }_+,`$ (87)
$`K`$ $`:`$ $`\mathrm{\Phi }_{}\mathrm{\Phi }_+\mathrm{\Phi }_{}.`$ (89)
Then the Wigner operator is not well defined either within $`\mathrm{\Phi }_+`$ or within $`\mathrm{\Phi }_{}`$, so those states in $`\mathrm{\Phi }_+`$ or $`\mathrm{\Phi }_{}`$ are not, in general, $`t`$symmetric. Therefore if we consider that only the states of $`\mathrm{\Phi }_+`$ are “physical” or “admissible” the Wigner time inversion transforms these states in “unphysical” ones and therefore it turns out to be impossible since unphysical states simply do not exist in nature. Then, through this mathematical structure, irreversibility is incorporated in our theory.
Nevertheless $`\left(\mathrm{\Phi }_+\mathrm{\Phi }_{}\right)`$ is not an empty set , so the time-inversion operator will be well defined there
$$K:(\mathrm{\Phi }_+\mathrm{\Phi }_{})(\mathrm{\Phi }_+\mathrm{\Phi }_{}),$$
(90)
and these states will describe reversible processes, which will be $`t`$symmetric.
All these things lead us to the postulate of the introduction: “Physical states are in $`\mathrm{\Phi }_+`$ ( or $`\mathrm{\Phi }_{}`$).” In fact, this postulate provides a mathematical structure to deal with irreversible process, as it was also shown in papers . The choice between $`\mathrm{\Phi }_+`$ and $`\mathrm{\Phi }_{}`$ is conventional and does not lead to physical consequences, but once we choose one of these spaces the distinction between past and future becomes substantial. Moreover, if we take into account the global structure of the universe, this choice can be motivated from the asymmetry of this structure. (, )
## VII Concluding remarks
We outline the main results of this work.
We have diagonalized the full Hamiltonian of our model and extended it in such a way that the solution is analytic in the interaction parameter. A rigorous mathematical formalism can be introduced in order to deal with unstable quantum systems (see the appendices).
Using this formalism, we have obtained a second quantized version of the decay of unstable systems and we have found the corresponding creation and annihilation operators of unstable states.
By means of a simple example, the exact time evolution of the reduced density matrix at zero temperature has been studied. We have obtained an exponential decay approach of $`P(t)`$ to the asymptotic value $`P(t\mathrm{})=0`$ which is expected when the particle is in thermal equilibrium with a zero temperature bath. For short times we have found a quadratic behavior for the decay probability $`P(t)`$ (Zeno effect). This short time deviation from the exponential decay law was recently measured by first time (see Ref. ). Other deviations from the exponential decay law, in this case for long times, naturally arise in our framework: Khalfin effect, which unfortunately are in practice far enough of any observational time scale.
## VIII Acknowledgments
We acknowledge the very helpful discussions with Edgardo García Alvarez and Roberto Laura. This work was partially supported by grants CI1\*-CT94-0004 and PSS 0992 of European Union, PID 3183/93 of CONICET, EX053 of Universidad de Buenos Aires, and 12217/1 of Fundación Antorchas.
## IX Appendices
## A Creation and annihilation operators
Consider the annihilation and creation “unsmeared operators,” $`b_\omega `$ and $`b_\omega ^{}`$, respectively, that we have used in our calculations. Commonly they are introduced in the mathematical framework of quantum field theory by virtue of expressions like
$$b\left(\varphi \right)=\varphi \left(\omega \right)b_\omega 𝑑\omega $$
(A1)
or
$$b^{}\left(\varphi \right)=\varphi ^{}\left(\omega \right)b_\omega ^{}𝑑\omega ,$$
(A2)
where $`b\left(\varphi \right)`$ and $`b^{}\left(\varphi \right)`$ are the (well-defined or “smeared”) annihilation and creation operators of the one-particle state $`\varphi `$ (being $``$ a Hilbert space) and their action is interpreted as the annihilation, respectively the creation, of a spectrum localized quanta, represented by a Dirac delta centered in the real value $`\omega `$. In this context we find that
$`b^{}\left(\varphi \right):{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}sym\left(^n\right){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}sym\left(^n\right),`$
$$b^{}\left(\varphi \right)\mathrm{\Phi }=b^{}\left(\varphi \right)(\varphi _0,\varphi _1,\mathrm{},\varphi _n,\mathrm{})=(0,b_0^{}\left(\varphi \right)\varphi _0,b_1^{}\left(\varphi \right)\varphi _1,\mathrm{},b_n^{}\left(\varphi \right)\varphi _n,\mathrm{}),$$
(A3)
where
$`b_n^{}\left(\varphi \right)\varphi _n`$ $`=`$ $`b_n^{}\left(\varphi \right)\left[sym\left(\varphi _n^{\left(1\right)}\varphi _n^{\left(2\right)}\mathrm{}\varphi _n^{\left(n\right)}\right)\right]`$ (A4)
$`=`$ $`sym\left(\varphi _n^{\left(1\right)}\varphi _n^{\left(2\right)}\mathrm{}\varphi _n^{\left(n\right)}\varphi \right)`$ (A5)
and $`b\left(\varphi \right)=\left[b^{}\left(\varphi \right)\right]^{}`$.
Of course, in the framework of the Hilbert space foundation of quantum mechanics “definitions” (A1) and (A2) are strictly formal.
It is evident that these equations are analogous to the formal definition of the Dirac delta
$$\varphi \left(x\right)=\varphi \left(\omega \right)\delta \left(x\omega \right)𝑑\omega ,$$
(A6)
so, as in the case of the Dirac delta, the rigorous meaning of the “unsmeared operators” $`b_\omega `$ and $`b_\omega ^{}`$ must be found in the rigged Hilbert space formulation of quantum mechanics.
One way to do this is considering the explicit definitions of $`b_\omega `$ and $`b_\omega ^{}`$ as distribution valued operators. Being $`𝒮`$ $`𝒮^\times `$ a rigged Hilbert space, we have that
$`b_\omega :{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}sym\left(𝒮^n\right)\left[{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}sym\left(𝒮^n\right)\right]^\times ,`$
$$b_\omega \mathrm{\Phi }=b_\omega (\varphi _0,\varphi _1,\mathrm{},\varphi _n,\mathrm{})=(b_{1\omega }\varphi _1,b_{2\omega }\varphi _2,\mathrm{},b_{n\omega }\varphi _n,\mathrm{}),$$
(A7)
where
$$b_{n\omega }\varphi _n=b_{n\omega }\left[sym\left(\varphi _n^{\left(1\right)}\varphi _n^{\left(2\right)}\mathrm{}\varphi _n^{\left(n\right)}\right)\right]$$
(A8)
and $`b_{n\omega }`$ acts as in its definition given in Sec. II. Also
$`b_\omega ^{}:{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}sym\left(𝒮^n\right)\left[{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}sym\left(𝒮^n\right)\right]^\times ,`$
$$b_\omega ^{}\mathrm{\Phi }=b_\omega ^{}(\varphi _0,\varphi _1,\mathrm{},\varphi _n,\mathrm{})=(0,b_{0\omega }^{}\varphi _0,b_{1\omega }^{}\varphi _1,\mathrm{},b_{n\omega }^{}\varphi _n,\mathrm{}),$$
(A9)
where
$`b_{n\omega }^{}\varphi _n`$ $`=`$ $`b_{n\omega }^{}\left[sym\left(\varphi _n^{\left(1\right)}\varphi _n^{\left(2\right)}\mathrm{}\varphi _n^{\left(n\right)}\right)\right]`$ (A10)
$`=`$ $`sym\left(\varphi _n^{\left(1\right)}\varphi _n^{\left(2\right)}\mathrm{}\varphi _n^{\left(n\right)}\delta _\omega \right).`$ (A11)
Observe that, with these definitions, only the normal product is well defined and that the “” symbol is merely a convenient notation, that is, there is not a rigorous Hermitian conjugation involved. But if we want to define the canonical commutation relations we will be in troubles, because the product of functionals is not uniquely defined. So this is not the right way either.
Fortunately, there is another way to define the annihilation and creation “operators,” $`b_\omega `$ and $`b_\omega ^{}`$. This point of view is more abstract than the previous one. Remember that the Dirac delta can be considered as a tempered distribution, i.e. a continuous linear functional on the one-particle regular states space $`𝒮`$. In an analogous way the annihilation and creation “operators,” $`b_\omega `$ and $`b_\omega ^{}`$, and all the respective “products” can be considered as continuous linear functionals on the Canonical Commutation Relations Algebra = $`CCR\left(𝒮\right)`$.
We summarize some properties which characterize this kind of algebras . Remember that, being $`𝒮`$ a nuclear metrizable space, there exists a non-decreasing basis $`\left\{p_\alpha \right\}_{\alpha I}`$ of continuous seminorms such that each seminorm is Hilbertian. Let us denote by $`_\alpha `$ the Hilbert space which is the completion of the quotient space $`𝒮/Ker\left(p_\alpha \right)`$ with respect to the quotient norm $`\widehat{p}_\alpha =p_\alpha /Ker\left(p_\alpha \right),`$ i.e. the space of equivalence classes defined by
$`\varphi _\alpha =\{\chi 𝒮:p_\alpha \left(\varphi \chi \right)=0\},`$
where
$`\widehat{p}_\alpha :𝒮/Ker\left(p_\alpha \right)𝐑_+,`$
$`\widehat{p}_\alpha \left(\varphi _\alpha \right)=p_\alpha \left(\varphi \right).`$
The -algebra $`CCR\left(𝒮\right)`$ is defined as the Hausdorff projective limit of the collection $`\left\{CCR\left(_\alpha \right)\right\}_{\alpha I}`$, where $`CCR\left(_\alpha \right)`$ is the C-algebra generated by the family of operators $`\{b\left(\varphi \right):\varphi _\alpha \},`$ with respect to the mappings that inject each $`CCR\left(_\alpha \right)`$ into $`CCR\left(_\beta \right)`$ if $`\alpha \beta `$, where the order in $`I`$ is the induced by the ordering of the basis of seminorms. We can characterize the $`CCR\left(𝒮\right)`$ as follows. Since every projective limit of a collection of C-algebras is a b-algebra, i.e. a complete symmetric -algebra whose topology is defined by a basis of continuous submultiplicative seminorms, $`CCR\left(𝒮\right)`$ is also a b-algebra (see Ref. ). Moreover, we have that $`CCR\left(𝒮\right)`$ is the strict inductive limit of the collection $`\left\{_{j=0}^nsym\left(𝒮^j\right)\right\}_{n=0}^{\mathrm{}}`$, so $`CCR\left(𝒮\right)`$ is a nuclear strict inductive limit of a collection of Fréchet spaces or $`^{}`$-algebra . With this we have that the algebra is complete, barreled, and nuclear.
Finally we have that $`CCR\left(𝒮\right)CCR\left(\right)\left[CCR\left(𝒮\right)\right]^\times ,`$ which represents a generalized Gel’fand triplet. So, viewing the relations (A1) and (A2) as generalized expansions in the sense of the well known Maurin’s theorem, one can identify $`b_\omega `$ and $`b_\omega ^{}`$ as continuous linear functionals on the algebra $`CCR\left(𝒮\right),`$ and we can say that they belong to $`\left[CCR\left(𝒮\right)\right]^\times `$.
## B Rigged extension
In this appendix we find a state space $`\mathrm{\Phi }`$ with the required properties to implement our formalism of unstable states. In order to adequately define the vectors obtained in Sec. III we must restrict the Hilbert space, which is the basic mathematical structure of ordinary quantum mechanics. Recall that Dirac’s bras are defined as linear functionals on kets of $`.`$<sup>\*‡</sup><sup>\*‡</sup>\*‡$``$ is the Hilbert space of the states we are considering. It can be the whole space of states or some subspace with precise physical properties, as the incoming or outgoing spaces. These functionals belong to $`^^\times ,`$ the topological dual of $`.`$ But in this case $`^^\times `$ is isomorphic to $`,`$ then one works indistinguishably with kets and bras. However, if we restrict the topology in order to take a dense subset $`\mathrm{\Phi }`$ of $`,`$ we break the one-to-one correspondence between elements $`\varphi `$ of $`\mathrm{\Phi }`$ and continuous linear or antilinear functionals $`F`$ over them. We will call $`\mathrm{\Phi }^{}`$ the dual space of linear functions and $`\mathrm{\Phi }^\times `$ the dual space of antilinear functional. We usually use the latter one. It leads to a triplet structure symbolized as $`\mathrm{\Phi }\mathrm{\Phi }^\times ,`$ where, to assure the convergence in the norm which defines the topology of $`\mathrm{\Phi },`$ we must require that $`\varphi |F`$ would be finite . This space $`\mathrm{\Phi }`$ will be the space of “regular” states $`\mathrm{\Phi }_+,`$ as we have explained above. Changing the our convention it can be $`\mathrm{\Phi }_{}`$
In our case, a necessary condition for $`|\varphi \mathrm{\Phi }`$ is that, the following expression, a generalization of Eq. (LABEL:aaa), would have a rigorous meaning,
$`\varphi |\stackrel{~}{n,v}`$ $`=`$ $`\varphi \left|\left[\stackrel{~}{a}^{()}\right]^n\right|0,v`$ (B1)
$`=`$ $`{\displaystyle \frac{1}{[\alpha ^{}(z_0)]^{\frac{n}{2}}}}[\varphi |nv+\lambda {\displaystyle _0^{\mathrm{}}}d\omega _1{\displaystyle \frac{g(\omega _1)}{[z_0\omega _1]_+}}\varphi |n1,\omega _1+\mathrm{}`$ (B3)
$`+\lambda ^n{\displaystyle _0^{\mathrm{}}}\mathrm{}{\displaystyle _0^{\mathrm{}}}d\omega _1\mathrm{}d\omega _n{\displaystyle \frac{g(\omega _1)}{[z_0\omega _1]_+}}\mathrm{}{\displaystyle \frac{g(\omega _n)}{[z_0\omega _n]_+}}\varphi |0\omega _1\mathrm{}\omega _n].`$
The last term of the second member of Eq. (B3) must be well defined, so the function $`\varphi |0,\omega _1\mathrm{}\omega _n=\varphi _0^{}(\omega _1\mathrm{}\omega _n)`$ has to have an analytic continuation in each variable $`\omega _i`$ ($`0in)`$ to a region which includes the singularity $`z_0`$, so that the integral defines an analytic $`n`$-dimensional function evaluated at $`z_0`$.
The simplest choice for $`\varphi |0,\omega _1\mathrm{}\omega _n`$ which does not depend on the localization of $`z_0`$ is that $`\varphi |0,\omega _1\mathrm{}\omega _n`$ would be a Hardy function from below $`H_{}^2`$ for each variable. It is equivalent to
$$\varphi |0,\omega _1\mathrm{}\omega _n\theta (𝒮H_{}^2)^n.$$
(B4)
This generalizes the criterium previously used for the one-particle sector . From this criterium it can be proved that all the mathematical expressions above are well defined (see Appendix C). Then in order to $`i,\omega _1\mathrm{}\omega _m|\varphi `$ $`(i+m=n,`$ $`n𝒩)`$ be well defined, it must belong to the following function space
$$i,\omega _1\mathrm{}\omega _m|\varphi \underset{m=0}{\overset{\mathrm{}}{}}\theta \left[𝒮H_+^2\right]^m,$$
(B5)
where $`𝒮`$ is the Schwartz space , and $`\theta `$ is the Heaviside step function, which gives the restriction to the positive real axis.
If we do the same in order to define $`\varphi |\overline{n,v}`$, we find another realization space
$$i,\omega _1\mathrm{}\omega _m|\varphi \underset{m=0}{\overset{\mathrm{}}{}}\theta \left[𝒮H_{}^2\right]^m,$$
(B6)
with $`(i+m=n,`$ $`n𝒩)`$. Therefore we define the following spaces
$$\mathrm{\Phi }_\pm =\left\{\varphi /i,\omega _1\mathrm{}\omega _m|\varphi \underset{m=0}{\overset{\mathrm{}}{}}\theta \left[𝒮H_\pm ^2\right]^m\right\}.$$
(B7)
The generalized eigenvectors belong to the dual spaces $`\mathrm{\Phi }_\pm ^\times `$, since they are antilinear functionals on spaces (B5) and (B6),
$`|i,\stackrel{~}{\omega _1\mathrm{}\omega _m}`$ $``$ $`\mathrm{\Phi }_+^\times ,`$ (B8)
$`|\overline{i,\omega _1\mathrm{}\omega _m}`$ $``$ $`\mathrm{\Phi }_{}^\times .`$ (B10)
These generalized eigenvectors fulfill the following relations (see )
$`i,\stackrel{~}{\omega _1\mathrm{}\omega _m}|\stackrel{~}{i^{},\omega _1^{}\mathrm{}\omega _m^{}^{}}`$ $`=`$ $`0,`$ (B11)
$`\overline{i,\omega _1\mathrm{}\omega _m}|\overline{i^{},\omega _1^{}\mathrm{}\omega _m^{}^{}}`$ $`=`$ $`0,`$ (B13)
$$\overline{i,\omega _1\mathrm{}\omega _m}|\stackrel{~}{i^{},\omega _1^{}\mathrm{}\omega _m^{}^{}}=\frac{\delta _{ii^{}}\delta _{mm^{}}}{(m!)^2}\underset{\sigma G_p}{}\underset{\tau G_p}{}\delta \left(\omega _{\sigma _1}^{}\omega _{\tau _1}\right)\mathrm{}\delta \left(\omega _{\sigma _m}^{}\omega _{\tau _m}\right).$$
(B14)
Eqs. (LABEL:conmu1) are valid except for $`i=i^{}=0,`$ and in this case we have
$`\stackrel{~}{0,\omega _1\mathrm{}\omega _n}|\stackrel{~}{0,\omega _1^{}\mathrm{}\omega _n^{}^{}}`$ $`=`$ $`{\displaystyle \frac{\delta _{nn^{}}}{(n!)^2}}{\displaystyle \underset{\sigma ,\tau G_p}{}}\delta \left(\omega _{\sigma _1}^{}\omega _{\tau _1}\right)\mathrm{}\delta \left(\omega _{\sigma _n}^{}\omega _{\tau _n}\right),`$ (B15)
$`\overline{0,\omega _1\mathrm{}\omega _n}|\overline{0,\omega _1^{}\mathrm{}\omega _n^{}^{}}`$ $`=`$ $`{\displaystyle \frac{\delta _{nn^{}}}{(n!)^2}}{\displaystyle \underset{\sigma ,\tau G_p}{}}\delta \left(\omega _{\sigma _1}^{}\omega _{\tau _1}\right)\mathrm{}\delta \left(\omega _{\sigma _n}^{}\omega _{\tau _n}\right).`$ (B17)
$`G_p`$ is the group of permutations. Eqs. (LABEL:conmu1) say that the norm of generalized eigenvectors is zero (except $`i=0`$), which is a necessary fact to conserve energy . It is not contradictory to have null-norm vectors because these are generalized vectors which are not in the usual Hilbert space and have an underlying indefinite metric structure . If we define the spaces $`_+`$ and $`_{}`$ as the spaces $`\mathrm{\Phi }_{}`$ and $`\mathrm{\Phi }_+`$ where the condition about $`𝒮`$ is not required, we have
$`_\pm =\left\{\varphi /n,\omega _1\mathrm{}\omega _m\mathrm{}|\varphi {\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}\theta \left[H_\pm ^2\right]^m\right\},`$
then we arrive to the triplets under Eq. (LABEL:Id) and in Eq. (LABEL:RHS).In this way $`\mathrm{\Phi }_\pm `$ is dense in $`_\pm `$ which is the outgoing (incoming) space . The $``$ cited in paper is actually the outgoing space $`_+.`$
Using Eq. (B7 ), we find a double structure of rigged Hilbert spaces for our model,
$`\mathrm{\Phi }_+`$ $``$ $`_+\mathrm{\Phi }_+^\times ,`$ (B18)
$`\mathrm{\Phi }_{}`$ $``$ $`_{}\mathrm{\Phi }_{}^\times .`$ (B20)
## C The double integral theorem
In Ref. it was demonstrated that if we want that the integral
$$_{𝐑^+}𝑑\omega \frac{g(\omega )\varphi |\omega }{z_0\omega }$$
(C1)
would be well defined, it is sufficient that
$$\varphi |\omega \theta (𝒮H_{}^2).$$
(C2)
In the two variables case it is the integral
$$_{𝐑^+}𝑑\omega \frac{g(\omega )}{z_0\omega }_{𝐑^+}𝑑\omega ^{}\frac{g(\omega ^{})\varphi |\omega ,\omega ^{}}{z_0\omega }$$
(C3)
the one that must be well defined. In this case we prove the following theorem.
Theorem.
If
$$\varphi (\omega ,\omega ^{})=\varphi |\omega ,\omega ^{}\theta (𝒮H_{}^2)^2,$$
(C4)
then integral (C3) is well defined.
Proof.
If condition (C4) is fulfilled, as $`𝒮`$ is a Fréchet spaceA Fréchet space is a metrizable complete space. we have (, page 459)
$$\varphi (\omega ,\omega ^{})=\underset{i=0}{\overset{\mathrm{}}{}}\lambda _i\varphi _1^i(\omega )\varphi _2^i(\omega ^{}),$$
(C5)
where $`_{i=0}^{\mathrm{}}|\lambda _i|<1,\varphi _1^i(\omega ),`$ $`\varphi _2^i(\omega )\theta (𝒮H_{}^2)`$ $`(i=1,2,\mathrm{})`$, $`\varphi _1^i,`$ $`\varphi _2^i0`$ when $`i\mathrm{},`$ and the r.h.s. of Eq. (C5) is absolutely convergent, namely the series
$`{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}p(\lambda _i\varphi _1^i\varphi _2^i)`$
is convergent for any continuous seminorm $`p`$ over $`\theta (𝒮H_{}^2)^2.`$ Let us now define the seminorm $`p_{z_0}`$ as
$$p_{z_0}(\varphi )=D^2_{𝐑^+}𝑑\omega \frac{|g(\omega )|}{|z_0\omega |}_{𝐑^+}𝑑\omega ^{}\frac{|g(\omega ^{})||\varphi |\omega ,\omega ^{}|}{|z_0\omega ^{}|},$$
(C6)
where $`D`$ is the distance from $`𝐑_+`$ to $`z_0`$ ($`D=\gamma /2`$). We must demonstrate that $`p_{z_0}`$ is a continuous seminorm. We use the Hölder inequality
$$\frac{g(\omega )}{z_0\omega }\frac{g(\omega ^{})\varphi |\omega ,\omega ^{}}{z_0\omega }_1=D^2p_{z_0}(\varphi )g(\omega )g(\omega ^{})\varphi (\omega ,\omega ^{})_1\frac{1}{(z_0\omega )(z_0\omega ^{})}_{\mathrm{}}$$
(C7)
for any $`z_0𝐂_{}`$ (the lower half-plane).<sup>\*∥</sup><sup>\*∥</sup>\*∥We remember that
$`f_1={\displaystyle _{𝐑_+}}|f(\omega )|𝑑\omega ,`$
$`f_{\mathrm{}}=sup\{|f(\omega )|;\omega 𝐑_+\},`$
and also that $`𝒮L_1`$ , i.e. any Schwartz function is integrable. Since
$$\frac{1}{(z_0\omega )(z_0\omega ^{})}_{\mathrm{}}=sup\left\{\left|\frac{1}{(z_0\omega )(z_0\omega ^{})}\right|;\omega ,\omega ^{}𝐑_+\right\}=D^2$$
(C8)
Equation (C7) reads
$$p_{z_0}(\varphi )g(\omega )g(\omega ^{})\varphi (\omega ,\omega ^{})_1$$
(C9)
for any $`z_0𝐂_{}.`$ Then $`p_{z_0}(\varphi )`$ is not only a continuous seminorm but also a continuous norm over $`\theta (𝒮H_{}^2)^2.`$ Then
$`{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}p_{z_0}[\lambda _i\varphi _1^i(\omega )\varphi _2^i(\omega ^{})]=D^2{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}|\lambda _i|{\displaystyle _{𝐑^+}}𝑑\omega {\displaystyle \frac{|g(\omega )|}{|z_0\omega |}}{\displaystyle _{𝐑^+}}𝑑\omega ^{}{\displaystyle \frac{|g(\omega ^{})||\varphi _1^i(\omega )||\varphi _2^i(\omega ^{})|}{|z_0\omega ^{}|}}`$
$$=D^2\underset{i=0}{\overset{\mathrm{}}{}}_{𝐑^+}𝑑\omega _{𝐑^+}𝑑\omega ^{}\left|\lambda _i\frac{g(\omega )}{z_0\omega }\frac{g(\omega ^{})}{z_0\omega ^{}}\varphi _1^i(\omega )\varphi _2^i(\omega ^{})\right|<\mathrm{}.$$
(C10)
So, calling
$$f_i^{z_0}(\omega ,\omega ^{})=\lambda _i\frac{g(\omega )}{z_0\omega }\frac{g(\omega ^{})}{z_0\omega ^{}}\varphi _1^i(\omega )\varphi _2^i(\omega ^{}),$$
(C11)
from the corollary of the Lebesgue theorem (, page 33) we know that, if the series
$`{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}f_i^{z_0}(\omega ,\omega ^{})<\mathrm{}`$
converge a.e. in $`𝐑_+\times 𝐑_+`$, then the series, considered as a function of $`(\omega ,\omega ^{}),`$ belongs to $`L_1`$ and
$$_{𝐑^+}𝑑\omega _{𝐑^+}𝑑\omega ^{}\underset{i=0}{\overset{\mathrm{}}{}}f_i^{z_0}(\omega ,\omega ^{})=\underset{i=0}{\overset{\mathrm{}}{}}_{𝐑^+}𝑑\omega _{𝐑^+}𝑑\omega ^{}f_i^{z_0}(\omega ,\omega ^{}).$$
(C12)
Then, going back to Eq. (C3) we have
$`{\displaystyle _{𝐑^+}}𝑑\omega {\displaystyle \frac{g(\omega )}{z_0\omega }}{\displaystyle _{𝐑^+}}𝑑\omega ^{}{\displaystyle \frac{g(\omega ^{})\varphi |\omega ,\omega ^{}}{z_0\omega }}={\displaystyle _{𝐑^+}}𝑑\omega {\displaystyle \frac{g(\omega )}{z_0\omega }}{\displaystyle _{𝐑^+}}{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}\lambda _i\varphi _1^i(\omega )\varphi _2^i(\omega ^{})`$
$$=\underset{i=0}{\overset{\mathrm{}}{}}\lambda _i_{𝐑^+}𝑑\omega \frac{g(\omega )\varphi _1^i(\omega )}{z_0\omega }_{𝐑^+}𝑑\omega ^{}\frac{g(\omega ^{})\varphi _2^i(\omega ^{})}{z_0\omega }.$$
(C13)
The l.h.s. of Eq. (C13) is well defined since it is an integral of a $`L_1`$ function. Moreover it is the sum of products of two well defined integrals, like the one of Eq. (C1), since from hypothesis $`\varphi _1^i(\omega ),\varphi _2^i(\omega ^{})\theta (𝒮H_{}^2).`$ Thus, the proof is complete.
Of course, this theorem can be generalized from the case of two factors, to the case of $`n`$ factors and, taking into account Eq. (32), it proves that Eq. (B3) is well defined if condition (B4) is fulfilled.
## D Figure captions
Fig. 1: Deformation of the contour of integration taking into account the presence of the complex pole $`z_0.`$
Fig. 2: Behavior of the decay probability showing the Zeno, exponential, and Khalfin phases.
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# Zinc as a tracer of metallicity evolution of Damped Ly 𝛼 systems
## 1 Introduction
The redshift-metallicity relation of QSO absorption systems is a fundamental probe of the chemical evolution of the universe. The study of such relation in Damped Lyman $`\alpha `$ (DLA) systems<sup>1</sup><sup>1</sup>1The QSO absorption systems with neutral hydrogen column density $`N`$(HI) $`>2\times 10^{20}`$ atoms cm<sup>-2</sup> are called Damped Lyman $`\alpha `$ systems owing to the presence of radiative damping wings in their Lyman $`\alpha `$ absorption profiles. probes, in particular, the metal enrichment of the associated galaxies located at cosmological distances along the QSO line of sight (Lu et al. 1996). Abundance studies of DLA systems can be used to trace the chemical evolution of galaxies in the early universe, starting from the redshifts of the most distant QSOs.
Abundance determinations for different elements are presently available for about 60 DLA systems, but the study of their redshift evolution is hampered by two main difficulties. One is the uncertainty of the abundance measurements owing to the possible effects of dust depletion. The other is the limited redshift coverage of the sample, which is not always adequate for probing the presence of evolution. One approach to tackle the first difficulty is to correct the observed abundances for dust depletion effects (Vladilo 1998, Savaglio et al. 2000). Another approach is to use an element unaffected by dust depletion as a tracer of metallicity evolution (Pettini et al. 1997, 1999). Zinc is known to have little affinity with dust since it is essentially undepleted in the interstellar medium (Roth & Blades 1995). Abundance measurements in Galactic metal-poor stars yield \[Zn/Fe\]<sup>2</sup><sup>2</sup>2We adopt the usual convention \[X/Y\] = log $`N`$(X)/$`N`$(Y) $``$ log (X/Y) $`0`$ (Sneden, Gratton & Crocker 1991), suggesting that zinc is a good tracer of iron. The quality of available stellar data leaves open the possibility that new, more precise measurements may reveal small deviations of \[Zn/Fe\] from the solar ratio such as those found for the iron-peak elements Cr, Mn and Co (Ryan, Norris & Beers 1996). However, even in this case, zinc would still be a good indicator of metallicity, even if not a perfect tracer of Fe.
Studies of zinc abundances in DLA systems have not revealed evidence for evolution of the column-density weighted \[Zn/H\] metallicity (Pettini et al. 1997, 1999). One difficulty in detecting evolution is the lack of measurements at $`z_{\mathrm{abs}}>3`$, when the redshifted Zn II resonant doublet $`\lambda _{\mathrm{rest}}2025,2062`$ Å falls in the reddest part of the visible spectrum. In addition, the column-density weighted metallicity is more prone to be affected by low number statistics than the unweighted metallicity, as we discuss in Section 3. Here we present the results of a search for redshift evolution performed by considering both the unweighted and the weighted metallicities (Sections 2 and 3, respectively) and based on the sample of \[Zn/H\] literature data, which includes now our recent measurement at $`z_{\mathrm{abs}}3.4`$ (Molaro et al. 2000).
## 2 The anti-correlation between \[Zn/H\] and redshift
In Fig. 1 we show the metallicity \[Zn/H\] in DLA systems derived from the Zn II and H I column density determinations currently available in literature. References to the original works are indicated in the figure caption. The UVES measurement of the $`z_{\mathrm{abs}}=3.3901`$ system in QSO 0000-26 is the big circle at the right bottom of the figure. When necessary, the zinc column densities were redetermined by using the oscillator strengths by Bergeson & Lawler (1993) which are adopted for the full sample. The meteoritic abundance by number $`\mathrm{log}(\mathrm{Zn}/\mathrm{H})_{}=7.33\pm 0.04`$ (Grevesse, Noels & Sauval 1996) was adopted in all cases as a solar reference value. The error bars shown in the figure were derived by error propagation of the Zn II and H I column density errors quoted by the authors and of the meteoritic abundance error.
No systematic differences appear to be present between the ZnII column densities obtained from 4-m class telescopes (e.g. Pettini et al. 1994, 1997, 1999) and those measured with the Keck telescope (e.g. Lu et al. 1996, Prochaska & Wolfe 1996, 1997, 1999). In the few cases when a re-determination of the same Zn II column density with higher resolution and signal-to-noise ratio is available the difference is well within the observational errors quoted by the authors. The reason for this is probably the weak degree of saturation of the Zn II doublet which allows accurate column densities to be determined even with spectra of modest resolution and signal-to-noise ratio. We believe therefore that the data sample shown in Fig. 1 is sufficiently homogeneous to perform a statistical analysis.
Since the dependence of \[Zn/H\] as a function of $`z`$ is not known a priori, we first performed a non parametric correlation analysis by means of Kendall’s $`\tau `$ using routine kendl1 of Press et al. (1992). We obtain evidence for rank reversal ($`\tau =0.29`$), — i.e. evidence for anti-correlation — at 97.2% confidence level. We then performed a least square linear regression through the {$`z_{\mathrm{abs}}`$, \[Zn/H\]} data points and obtained an anti-correlation with slope $`m=0.33\pm 0.12`$ and intercept $`q=0.43\pm 0.24`$. Student’s t tests indicate that the slope and the intercept differ from zero at 98.7% and 92.0% confidence level, respectively. The linear regression results are confirmed by a bootstrap analysis with 10,000 samples, which yields a mean slope $`<m>=0.32\pm 0.13`$ and a mean intercept $`<q>=0.44\pm 0.25`$.
While the evidence for anti-correlation is quite robust, the data show a scatter larger than the typical measurement errors. The scatter is still above the individual errors even considering only the Keck sub-sample. Therefore the dispersion is genuine and not due to inhomogeneity of the data. An intrinsic scatter is expected if DLA systems are associated with galaxies of different types and different chemical evolution histories. An additional source of scatter is the random galactocentric distances of the gas intercepted, in the presence of metallicity gradients.
The influence of intrinsic scatter is a matter of concern in performing the ordinary least squares regression on the data. Therefore we re-analysed the data sample with the BCES method of Akritas & Bershady (1996) which is ideally suited to cope with this situation. We find a slope $`m_{\mathrm{BCES}}=0.27\pm 0.11`$ and an intercept $`q_{\mathrm{BCES}}=0.54\pm 0.20`$. The mean slope and intercept of 10,000 bootstrap samples analysed with the BCES method are $`<m_{\mathrm{BCES}}>=0.32\pm 0.13`$ and $`<q_{\mathrm{BCES}}>=0.45\pm 0.25`$. The BCES results confirm, within the errors, the parameters derived from the ordinary linear regression analysis.
Albeit small, zinc shows some dust depletion which may affect the above analysis. In the Galactic ISM, the typical zinc depletion is $``$ $`0.19`$ dex (Roth & Blades 1995), while in DLA systems is $`0.12`$ dex owing to the modest dust-to-gas ratios typical of these absorbers (Vladilo 1998). The intercept of the linear regressions that we find may be therefore underestimated by $`\stackrel{<}{}`$ 0.2 dex (probably $``$ 0.1 dex), a systematic error which is still within the statistical errors. As far as the slope is concerned, we may expect an increase of its absolute value (i.e. a steepening of the anti-correlation) as a consequence of correcting for dust depletion. The reason for this is that dust-to-gas ratio and metallicity are correlated in DLA systems (Vladilo 1998) and we may expect a negligible correction at high redshift, where the metallicity is low, but a non negligible correction at low redshift, where the metallicity rises. In the most extreme case (i.e. when we consider the depletion completely negligible at $`z=3.5`$ and we apply a ISM correction of $`+0.19`$ dex at $`z=0.5`$) we estimate that $`|m|`$ would increase by about 0.06 dex. The real systematic error is probably smaller and hence well within the statistical errors of the slope.
## 3 The column-density weighted metallicity
The analysis performed in the previous section is appropriate to study the evolution of metals in individual systems. However, in order to estimate the mean cosmic metallicity of the neutral gas in the universe, the abundance of each DLA system has to be weighted by its HI column density (Pei & Fall 1995, Pettini et al. 1997). It is easy to show that this is equivalent to compute the expression
$$<Z>=\mathrm{log}\frac{\underset{i}{}N(\mathrm{Zn}\mathrm{II})_i}{_iN(\mathrm{H}\mathrm{I})_i}\mathrm{log}(\mathrm{Zn}/\mathrm{H})_{},$$
(1)
where the sums are extended to the DLA systems in a given redshift bin. The mean cosmic metallicity $`<Z>`$ is useful for probing models of cosmic chemical evolution (Pei & Fall 1995), but it is particularly affected by low-number statistics. In fact, not only the number of systems available in each redshift bin is limited but, in addition, only the few of them with the highest $`N`$(HI) give a significant contribution to $`<Z>`$.
We estimated $`<Z>`$ in 4 redshift bins. Rather than adopting bins of constant redshift width, we binned 4 groups of DLA systems with adjacent redshifts, each group with comparable value of $`_iN(\mathrm{H}\mathrm{I})_i`$ and hence comparable statistical significance. We then determined the effective redshift of each bin with the expression $`<z_{\mathrm{abs}}>=_i\left[z_{\mathrm{abs},i}N(\mathrm{H}\mathrm{I})_i\right]/_iN(\mathrm{H}\mathrm{I})_i`$. The resulting 4 data points are shown as filled diamonds in Fig. 2. The error bars of $`<Z>`$ are estimated according to Eq. (3) by Pettini et al. (1997). A quick look at the figure does not show evidence for metallicity evolution. This result is confirmed by a regression analysis of the {$`<z_{\mathrm{abs}}>,<Z>`$} data points, which yields $`m=0.13\pm 0.07`$. This slope is lower than the one derived from the analysis of individual DLA systems and differs from zero only at 79% confidence level.
## 4 Summary and conclusions
We find evidence for an anti-correlation between the absolute zinc abundance \[Zn/H\] and the absorption redshift $`z_{\mathrm{abs}}`$ of DLA systems, with a slope $`0.3\pm 0.1`$ in the range $`0.5\stackrel{<}{}z_{\mathrm{abs}}\stackrel{<}{}3.5`$. The zinc metallicity increases from $`3\%`$ up to $`25\%`$ of the solar value from $`z_{\mathrm{abs}}3.5`$ to $`z_{\mathrm{abs}}0.5`$. Should DLA absorbers continue the same trend also from $`z_{\mathrm{abs}}0.5`$ to $`z_{\mathrm{abs}}0`$, the typical present-day metallicity would be $`35\%`$ solar, even though a value as high as $`60\%`$ is still within the errors of the intercept. Correcting for dust depletion effects would slightly steepen the anti-correlation, but well within the statistical error of the slope; the characteristic present-day metallicity would rise up to $`50\%`$ solar, with values as high as $`100\%`$ solar still within the errors.
The slope of the metallicity redshift relation that we derive is in good agreement with the value recently derived by Savaglio et al. (2000). The literature data base considered by these authors is larger than ours and includes different elements in addition to zinc. However, most of these elements are known to be severely depleted into dust in the ISM and the results by Savaglio et al. are based on an algorithm that corrects the abundances for depletion effects. The present results do not require a modeling of the elemental depletion patterns nor assumptions on the intrinsic abundance patterns in DLA systems.
In spite of the correlation with redshift, the zinc metallicities show evidence for intrinsic scatter. Models of galactic chemical evolution have already been able to explain such scatter by considering the surface brightness and the formation redshift of the galaxies, as well as the galactocentric distance of the gas intercepted (Jiménez, Bowen & Matteucci 1999). While such analysis has shown a general consistency between the zinc observations and the predicted evolutionary tracks, the present results demonstrate for the first time the evolution of zinc metallicity on pure observational grounds.
While we find evolution of the zinc metallicity of individual systems, we do not find evolution of the column-density weighted metallicity $`<Z>`$. This is consistent with the results of previous studies by Pettini et al. (1997, 1999). The lack of evolution of the zinc mean cosmic metallicity might be due, at least in part, to the lack of a sufficiently large data base, since the measurement of $`<Z>`$ is extremely sensitive to low number statistics. However, there are also reasons to believe that $`<Z>`$ is affected by some selection bias. In Fig. 2 we use different symbols for the systems with $`N`$(HI) $`>10^{21}`$ cm<sup>-2</sup> (empty squares) and $`N`$(HI) $`10^{21}`$ cm<sup>-2</sup> (empty circles). One can see from the figure that the systems with high column density have, in general, low metallicity — an effect originally pointed out by Boissé et al. (1998). High column density systems are the main contributors to $`<Z>`$ and the lack of such absorbers with \[Zn/H\] $`>1`$ at low $`z`$ tends to hide the global rise of metallicity with cosmic time. The lack of DLA systems of high column density and metallicity at low redshift is somewhat surprising because (i) clouds with $`N`$(HI) $`>10^{21}`$ atoms cm<sup>-2</sup> and high metallicity do exist in the disk of our Galaxy and in low redshift spirals; (ii) study of the HI content of the local universe suggest that spirals should be the main contributors to the DLA population at $`z0`$ (Rao & Briggs 1993). Nevertheless, spirals are a small fraction of the intervening DLA galaxies observed in low-$`z`$ imaging studies (Le Brun et al. 1997; Rao & Turnshek 1998; see also refs. in Table 1 by Vladilo 1999). This deficiency of spirals suggests the presence of some selection effect. Selection effects that can bias the observed population of DLA absorbers include QSO obscuration by DLA dust (Fall & Pei 1993) and gravitational lensing (Smette, Claeskens & Surdej 1997). However, also the surface brightness of the intervening galaxies and the galactocentric distances of the clouds intercepted can play a role in affecting the observed population. As discussed in Vladilo (1999), these effects generally conspire to decrease the fraction of chemical enriched regions in the sample population, dust obscuration alone yielding a QSO visual extinction of $`1`$ magnitude when $`N`$(HI) $`>10^{20.7}`$ cm<sup>-2</sup> at solar metallicity. Considering the likely presence of this bias and the severe dependence of $`<Z>`$ on low number statistics, the lack of evolution of $`<Z>`$ should not be used to conclude that the mean cosmic metallicity of DLA absorbers does not evolve. Comparison between empirical $`<Z>`$ determinations and model predictions of global enrichment of the universe should await a better understanding of the role played by any selection bias and a significant enlargement of the observational data base.
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# Two ways of biasing galaxy formation
## 1 Introduction
Galaxy clustering in the nearby Universe has been mapped through a variety of surveys, including different populations of luminous objects. These run from optical galaxies in the APM, CfA and LCRS surveys to the sources of the IRAS catalogue at 60 $`\mu `$m. Reconstructing the overall mass power spectrum from these data represents one of the main goals of modern cosmology. It is already known, however, that different tracer populations show different clustering amplitudes even after redshift-space and small-scale corrections are applied. Thus, their clustering patterns are not unambiguosly related to any one given mass power spectrum (see Peacock 1999 for a review).
The simplest and most common description of biasing adopted in the literature is that, at any spatial position $`𝐱`$, the fluctuation in the number density of galaxies $`\delta _g(𝐱)`$ responds linearly and locally to the underlying mass fluctuation $`\delta (𝐱)`$, namely $`\delta _g(𝐱)=b^E\delta (𝐱)`$, where $`b^E`$ is a space-independent bias factor (e.g. Dekel & Rees 1987). As discussed below, higher-order bias factors can be introduced, but the point is that such a bias prescription is inherently Eulerian: it relates the present-day galaxy and mass clustering properties, ignoring their past evolution. However, if gravity is the main force acting in the Universe, there is no doubt that galaxy biasing evolves in time, as collapsing mass fluctuations keep accreting luminous matter onto them, the galaxy distribution eventually relaxing to the mass one (Fry 1996; Tegmark and Peebles 1998). So, the biasing in the present-day galaxy distribution might well be rooted into the deep past of the history of the Universe: the strong Lyman break galaxy clustering seems to suggest that this might be the case (Steidel et al. 1997). Any primordial biasing, arising at the epoch of galaxy formation, cannot be described by an Eulerian model. Instead, a Lagrangian one has to be adopted: it is the primordial fluctuation in galaxies that is proportional to the mass fluctuation, $`\delta _g(𝐪)=b^L\delta (𝐪)`$, where $`𝐪`$ denotes the Lagrangian position; in general $`b^E`$ differs from $`b^L`$ and, in principle, higher order factors can be defined.
In this Letter we show that the local Eulerian and Lagrangian bias models are inconsistent. In fact, the clustering patterns predicted by the two bias models are different. Specifically we study the galaxy bispectrum and skewness, both in real and redshift space, on scales where the mildly non-linear approximation suffices, starting from Gaussian initial conditions. In Section 2 we review the general Eulerian and Lagrangian bias models in terms of infinite hierarchies of bias factors $`\{b_j^E\}`$ and $`\{b_j^L\}`$. In the whole discussion, these must be considered as free, position-independent, parameters. In Section 3 we discuss the galaxy bispectrum and skewness in real space, for both bias models. In Section 4 we carry out the same analysis, but taking into account the effect of redshift distortions. Section 5 contains our conclusions.
## 2 The twofold biasing prescription
Let us start by fixing the notation of basic quantities. If $`\phi _o(𝐪)`$ is the primordial gravitational potential (growing mode only, smoothed on some scale $`R_o`$ and linearly extrapolated to the present time), then $`\delta ^{(1)}(𝐪,z)=D(z)_q^2\phi _o(𝐪)`$ is the linear density field, and $`D(z)`$ its growth factor with $`z`$ the cosmological redshift \[we put $`D(0)=1`$\]. The linear peculiar velocity is given by $`𝐮^{(1)}(𝐪)=_q\phi _o(𝐪)`$, and it is constant in time. The Eulerian density field will be indicated by $`\delta (𝐱,z)`$ and the $`n`$-th order perturbative solutions $`\delta ^{(n)}`$ are such that $`\delta =_n\delta ^{(n)}`$ (Goroff et al. 1986). The Fourier transform is, e.g., $`\stackrel{~}{\delta }(𝐤)=𝑑𝐱\delta (𝐱)\mathrm{exp}i𝐤𝐱`$.
### 2.1 Local Eulerian bias
In this approach, the galaxy number density field at a given position $`𝐱`$ and time $`z`$ (e.g. ‘here’ and ‘now’) is assumed to be a local function of the underlying mass density field at the same location and instant, $`\delta _g(𝐱,z;R)[\delta (𝐱,z;R)],`$ where the smoothing scale $`R`$ is much larger than the typical size of the selected objects. Usually, assuming that $`[\delta ]`$ can be expanded about $`\delta =0`$ as a power series, an infinite set of “Eulerian bias factors” $`b_j^E`$ can be defined (Fry & Gaztañaga 1993):
$$\delta _g=\underset{j=0}{\overset{\mathrm{}}{}}\frac{b_j^E}{j!}\delta ^j.$$
(1)
This series is such that $`\delta _g=0`$ and $`\delta _g(\delta =1)=1`$. The linear coefficient $`b_1^E`$ corresponds to the usual bias factor. The origin of this local Eulerian prescription is essentially phenomenological, and it is a priori devoid of any insight about the dynamics of the clustering. Galaxy clustering is analyzed for instance in terms of $`N`$-point correlation functions $`_{n=2}^N\delta _g(𝐱_n,z)`$, and the bias factors are tuned to fit the observational data. This is the approach that has been implicitly adopted in most of the published literature on biasing, at least in its leading approximation.
### 2.2 Local Lagrangian bias
According to this alternative prescription, the sites of galaxy formation are identified with specific regions of the primordial density field. It is then appropriate to define a “primordial” galaxy density field, $`\delta _g(𝐪)`$, measuring the (smoothed) overdensity of galaxies in fieri at the Lagrangian position $`𝐪`$ at a given time $`z`$ (formally $`z=\mathrm{}`$, i.e. ‘there’ and ‘then’) which is biased with respect to the primordial (linear) density field at the same location and instant, namely $`\delta _g(𝐪)[ϵ_o(𝐪)]=_{j=0}^{\mathrm{}}(b_{oj}^L/j!)ϵ_o(𝐪)^j;`$ here $`ϵ_o(𝐪)=^2\phi _o(𝐪)`$ is the linear density field extrapolated to the present time. We can equivalently write, for similarity with the Eulerian case,
$$\delta _g(𝐪)=\underset{j=0}{\overset{\mathrm{}}{}}\frac{b_j^L}{j!}\delta ^{(1)j},$$
(2)
with $`\delta ^{(1)}(𝐪,z)=D(z)ϵ_o(𝐪)`$ and the Lagrangian factors $`b_j^L`$ are defined accordingly in terms of the original $`b_{oj}^L`$.
Both galaxies and dark matter flow through Eulerian space towards mass concentrations. Therefore, even assuming that the spots of galaxy formation can be identified in Lagrangian space, one has to consider large-scale motions in order to compute the statistics of present-day structures. One therefore needs to assign a dynamical prescription, because evolution changes the original galaxy distribution: this is the main difference with respect to the Eulerian bias scheme, where no dynamics is taken into account.
One way is to use galaxies as test particles of the underlying gravitational field (Fry 1996), then the evolved galaxy density field at the Eulerian position $`𝐱`$ and instant $`z`$ is related to the primordial galaxy field and evolved density field by the relation (Catelan et al. 1998)
$$1+\delta _g(𝐱,z)=[1+\delta _g(𝐪)][1+\delta (𝐱,z)].$$
(3)
We stress the fact that eq.(3) is inherently non-local. Smoothed regions in Lagrangian space can be mapped to Eulerian space through the transformation $`𝐱=𝐪+𝐒(𝐪,z)`$, where $`𝐒`$ is the displacement vector. In the Zel’dovich (1970) approximation, the simplest transformation, $`𝐒(𝐪,z)=D(z)𝐮^{(1)}(𝐪)`$. Thus, the resulting $`\delta _g(𝐱,z)`$ and $`\delta (𝐱,z)`$ are not deterministically related. In fact, for any given $`\delta `$ the galaxy field $`\delta _g`$ can assume different values (see Dekel & Lahav 1999). This stochastic behaviour is inherent to the gravitational instability dynamics.
The question now is the following: are the predictions about the clustering (in terms of standard statistics as the correlation functions, for example) as deduced from the local Eulerian and Lagrangian bias equivalent? In other terms, do there exist two sets of non-trivial Eulerian and Lagrangian bias factors, $`\{b_j^E\}`$ and $`\{b_j^L\}`$, such that the predictions for galaxy clustering are identical? In order to answer to these questions, let us analyze the galaxy bispectrum from Gaussian initial conditions as induced by mildly non-linear density evolution.
## 3 Galaxy bispectrum
### 3.1 Eulerian bias case
The lowest order contribution to the galaxy bispectrum $`(2\pi )^3\delta _D(𝐤_1+𝐤_2+𝐤_3)B_g(𝐤_1,𝐤_2,𝐤_3;z)=\stackrel{~}{\delta }_g(𝐤_1,z)\stackrel{~}{\delta }_g(𝐤_2,z)\stackrel{~}{\delta }_g(𝐤_3,z)`$ comes from the appearance of non-negligible second-order fluctuations $`\delta _g^{(2)}`$. From eq. (1), this is $`\delta _g^{(2)}(𝐱,z)=b_1^E\delta ^{(2)}(𝐱,z)+\frac{1}{2}b_2^E\delta ^{(1)2}(𝐱,z).`$ (Note that this expression does not have zero mean, so an offset term should be introduced; however, since we are interested in the spectral properties of the galaxy clustering, we will ignore it since it contributes only to $`𝐤=\mathrm{𝟎}`$.) Defining $`\nu _{12}𝐤_1𝐤_2/k_1k_2`$, and the second-order growth factor $`E\frac{3}{7}\mathrm{\Omega }^{2/63}D^2`$ either in an open Universe with no cosmological constant (Bouchet et al. 1992) or $`E\frac{3}{7}\mathrm{\Omega }^{1/140}D^2`$ in a Universe with a cosmological constant or quintessence (Kamionkowski & Buchalter 1999), we introduce the symmetric kernel
$$J_S^{(2)}\frac{1}{2}\left(1\frac{E}{D^2}\right)+\frac{1}{2}\left(\frac{k_1}{k_2}+\frac{k_2}{k_1}\right)\nu _{12}+\frac{1}{2}\left(1+\frac{E}{D^2}\right)\nu _{12}^2,$$
(4)
and the second-order convolution integral operator $`^{(2)}\stackrel{~}{\delta }^{(1)}\stackrel{~}{\delta }^{(1)}`$ (Fry 1984). Finally, we can simply write $`\stackrel{~}{\delta }^{(2)}=^{(2)}J_S^{(2)}`$ and
$$\stackrel{~}{\delta }_g^{(2)}=^{(2)}\left(b_1^EJ_S^{(2)}+\frac{1}{2}b_2^E\right).$$
(5)
Thus, the galaxy bispectrum is (Matarrese, Verde & Heavens 1997)
$$B_g^E=2D^4b_1^{E\mathrm{\hspace{0.17em}2}}[(b_1^EJ_S^{(2)}+\frac{1}{2}b_2^E)P(k_1)P(k_2)+\mathrm{c}.\mathrm{t}.],$$
(6)
where $`P(k,z)=D(z)^2P(k)`$ is the mass linear power spectrum. The skewness of the galaxy density field smoothed on scale $`R`$ is therefore (Fry 1994),
$$S_g^E(R)=\frac{b_1^E(42\frac{E}{D^2})+3b_2^Eb_1^E\gamma (R)}{b_1^{E\mathrm{\hspace{0.17em}2}}},$$
(7)
where $`\gamma =d\mathrm{ln}\sigma _R^2/d\mathrm{ln}R`$ and $`\sigma _R^2`$ is the rms density on scale $`R`$. For a scale-free mass power spectrum $`P(k)k^n`$ and a top-hat smoothing function, one obtains $`\gamma =n+3`$ (Bernardeau 1994). We remind the reader that in the Einstein–de-Sitter Universe, $`42E/D^2=34/7`$. Mass bispectrum can be recovered by setting in these formulae $`b_1^E=1`$ and $`b_2^E=0`$, $`B_m^E2D^4[J_S^{(2)}P(k_1)P(k_2)+\mathrm{c}.\mathrm{t}.]B_m`$ (Fry 1984). The growth of $`B_m`$ is self-similar, i.e. mass particles do not move from their initial positions, and the wavectors $`𝐤`$ actually correspond to those positions.
### 3.2 Lagrangian bias case
Let us now repeat the previous calculations assuming the Lagrangian biasing scheme in eq.(2). In this case, expanding eq.(3) up to second-order, we obtain, after Zel’dovich transforming the Lagrangian coordinate $`𝐪`$ to the Eulerian one $`𝐱`$ at $`z`$, $`\delta _g=b_0^L+(1+b_0^L+b_1^L)\delta ^{(1)}+\delta _g^{(2)}`$, where,
$$\delta _g^{(2)}=(1+b_0^L)\delta ^{(2)}Db_1^L𝐮^{(1)}\delta ^{(1)}+(b_1^L+\frac{1}{2}b_2^L)\delta ^{(1)2}.$$
(8)
This expression generalizes the analogous one in Catelan et al. (1998) for Press-Schechter dark matter halos, for which $`b_0^L=0`$. The Zel’dovich approximation, adopted here, suffices to transform from $`𝐪`$ to $`𝐱`$: this explains the presence of the inertia term, proportional to the velocity. (We assume the scale $`R_o`$ large enough so that shell-crossing is absent on scale $`R`$.) The Fourier transform of eq.(8) is
$$\stackrel{~}{\delta }_g^{(2)}=^{(2)}\left((1+b_0^L)J_S^{(2)}+_S^{(2)}\right),$$
(9)
where
$$_S^{(2)}b_1^L+\frac{1}{2}b_2^L+\frac{1}{2}b_1^L\left(\frac{k_1}{k_2}+\frac{k_2}{k_1}\right)\nu _{12}$$
(10)
describes the effects of Lagrangian biasing during the mildly nonlinear regime. So, the galaxy bispectrum is
$`B_g^L`$ $`=`$ $`2D^4(1+b_0^L+b_1^L)^2[((1+b_0^L)J_S^{(2)}+_S^{(2)})`$ (11)
$`\times `$ $`P(k_1)P(k_2)+\mathrm{c}.\mathrm{t}.],`$
and the galaxy skewness turns out to be
$$S_g^L(R)=\frac{(1+b_0^L)(42\frac{E}{D^2})+6b_1^L+3b_2^L(1+b_0^L+b_1^L)\gamma (R)}{(1+b_0^L+b_1^L)^2}.$$
(12)
We emphasize the fact that it is the term $`\frac{1}{2}b_1^L(k_1/k_2+k_2/k_1)\nu _{12}`$ in $`_S^{(2)}`$, Fourier transform of the inertia term $`b_1^L𝐮^{(1)}\delta ^{(1)}`$ in eq.(8), which carries the signature of the gravitational dynamics, inherently absent in the Eulerian description. Clearly, such a signature would reflect into a distinctive shape dependence of the galaxy bispectrum, best quantified by the ‘effective’ amplitude $`Q`$ (see below).
### 3.3 Disentangling the two biasing schemes
#### 3.3.1 Galaxy Bispectrum and Skewness
The different clustering predictions of the two biasing models may be emphasized simply by calculating the difference $`\mathrm{\Delta }B_g`$ of the bispectra in eq.(6) and eq.(11) or of the skewnesses $`\mathrm{\Delta }S_g`$ in eq.(7) and eq.(12). It can be easily verified that no two sets of nontrivial independent bias factors $`\{b_j^E\}`$ and $`\{b_j^L\}`$ can be found such that $`\mathrm{\Delta }B_g=0=\mathrm{\Delta }S_g`$. We can explicitly write the final expressions assuming that at least the lowest-order bias factors are related, namely $`b_1^E=1+b_0^L+b_1^L`$. This last relation may be easily derived in linear regime, but it shows to be preserved even during the mildly non-linear regime (Mo & White 1996; Mo, Jing & White 1997 for the case $`b_0^L=0`$). Thus $`\mathrm{\Delta }B_g=B_g^EB_g^L`$ is
$`\mathrm{\Delta }B_g`$ $`=`$ $`D^4(1+b_0^L+b_1^L)^2\{[(b_2^Eb_2^L)b_1^L(1+{\displaystyle \frac{E}{D^2}})`$ (13)
$`\times `$ $`(1\nu _{12}^2)]P(k_1)P(k_2)+\mathrm{c}.\mathrm{t}.\}`$
$``$ $`\mathrm{\Delta }B_g(1,2)+\mathrm{\Delta }B_g(1,3)+\mathrm{\Delta }B_g(2,3).`$
Correspondingly, the skewness difference is
$$\mathrm{\Delta }S_g=\frac{3(b_2^Eb_2^L)2b_1^L(1+\frac{E}{D^2})}{(1+b_0^L+b_1^L)^2}.$$
(14)
Intriguingly, the dependence on the filtering scale cancels out, and we are left with a residual difference in the skewness of the two bias models that is scale-independent. Expressions for the Einstein-de-Sitter Universe can be recovered by setting $`1+E/D^2=4/7`$.
The two biasing schemes cannot be distinguished by measuring the skewness alone, but they can be distinguished from the shape dependence of the bispectrum. The important point is that $`_S^{(2)}`$, the new term that arises in the Lagrangian-biasing bispectrum \[cf. eq. (11)\] is linearly independent of the mass kernel $`J_S^{(2)}`$ and the constant $`b_2^E/2`$, the two terms that make up the Eulerian-biasing bispectrum \[cf. eq. (6)\]. Thus, no combination of parameters $`b_1^E`$ and $`b_2^E`$ can allow the Eulerian-biasing bispectrum to mimic the Lagrangian-biasing bispectrum.
Since the bias of high-redshift populations is $``$ 3 while that of local populations are closer to unity, we heuristically expect bias evolution of order unity, and thus the Lagrangian-bias parameters $`b_0^L`$ and $`b_1^L`$ to be quantities of order unity (further modeling is required to give a more precise answer). It is likely that surveys such as the SDSS and 2dF will be able to measure the bispectrum with enough precision to distinguish the predictions of Eulerian and Lagrangian biasing if the bias parameters are of order unity. Matarrese, Verde & Heavens (1997) have determined that SDSS/2dF data should be able to determine $`b_1^E`$ and $`b_2^E`$ to roughly a few percent within the context of Eulerian-bias models; that is, the coefficients of $`J_S^{(2)}`$ and the constant in eq. (6) can be determined to a few percent. To distinguish between Eulerian and Lagrangian biasing requires that the data be fit to an additional term, $`_S^{(2)}`$, as well. Although we have not revisited the calculation in detail, it seems reasonable that if the coefficients of $`J_S^{(2)}`$ and the constant can be fit to a few percent, then the coefficients of $`_S^{(2)}`$ can be fit with a precision not much poorer. In this case, the data can discriminate between Eulerian- and Lagrangian-biasing schemes.
#### 3.3.2 $`Q`$-Amplitudes
Eulerian and Lagrangian biasings predict for the galaxy bispectrum two different shape dependences, which cannot be obtained from one another by simply tuning the bias parameters $`b_1^L,b_2^L`$ and $`b_1^E,b_2^E`$. An efficacious way of emphasizing such a shape dependence is through the bispectrum ‘amplitudes’ $`Q_g^E`$ and $`Q_g^L`$, where, for example,
$$Q_g^L\frac{B_g^L(𝐤_1,𝐤_2,𝐤_3,z)}{[P_g^L(k_1,z)P_g^L(k_2,z)+\mathrm{c}.\mathrm{t}.]},$$
(15)
and similarly for $`Q_g^E`$ (Fry 1984). $`Q`$-amplitudes are essentially insensitive to the scale and the overall geometry. We obtain
$$Q_g^E=Q_g^L+\frac{\mathrm{\Delta }B_g}{P_gP_g}.$$
(16)
It is very useful to express $`Q_g^E`$ and $`Q_g^L`$ in terms of the bispectrum amplitude of the underlying mass density distribution $`Q_mB_m/PP`$. We have, respectively,
$$Q_g^E=\frac{Q_m}{b_1^E}+\frac{b_2^E}{b_1^{E\mathrm{\hspace{0.17em}2}}},$$
(17)
for the Eulerian amplitude (Fry 1994), and, for the Lagrangian amplitude,
$`Q_g^L`$ $`=`$ $`{\displaystyle \frac{Q_m}{1+b_0^L+b_1^L}}+{\displaystyle \frac{b_2^L}{(1+b_0^L+b_1^L)^2}}`$ (18)
$`+`$ $`{\displaystyle \frac{(1+E/D^2)b_1^L}{(1+b_0^L+b_1^L)^2}}{\displaystyle \frac{[(1\nu _{12}^2)P(k_1)P(k_2)+\mathrm{c}.\mathrm{t}.]}{[P(k_1)P(k_2)+\mathrm{c}.\mathrm{t}.]}}.`$
The novelty contained in eq.(18) is the angular dependence appearing in the right hand side: it is not related to the $`Q_m`$–shape dependence, and it is independent from the values of $`\{b_i^L\}`$. Measurements of $`Q_g`$ from galaxy catalogs for two different shapes of the triangle $`𝐤_1+𝐤_2+𝐤_3=\mathrm{𝟎}`$ can in principle disentangle the two biasing factors $`b_1`$ and $`b_2`$ and the two biasing schemes as well. In Figure 1 we plot $`Q_g^E`$ and $`Q_g^L`$ for a $`\mathrm{\Lambda }CDM`$ model; the values of the biasing factors $`b_{1,2}^E`$ and the choice of the scales are based on Scoccimarro et al. (2000).
More sophisticated and predictive relations may be proposed if one assumes that the set of Lagrangian bias factors $`\{b_j^L\}`$ are not free parameters, as in the present discussion, but rigourously computed within the framework of a given theoretical model. The ‘excursion set’ formalism (Peacock & Heavens 1990; Bond et al. 1991), for example, where dark-matter halos are identified by first-upcrossings of a collapse thereshold, predicts that $`b_1^L`$ and $`b_2^L`$ are functions of both halo size and redshift (Mo & White 1996; Porciani et al. 1998).
## 4 Redshift distortion effects
Given that the two biasing schemes are in principle distinguishable, we proceed to calculate the Eulerian and Lagrangian bispectra in redshift space, which is where they are most likely to be measured. Peculiar motions associated with structures on any scale distort the clustering pattern in redshift space (Kaiser 1987). So, in order to reconstruct the actual distribution of galaxies from redshift catalogues, we must be able to invert the distortion process. This can be easily done if we consider a distant region of the Universe so that the distortions essentially occur along the line-of-sight, and we restrict to large scales for which the mildly non-linear approximation suffices. If $`𝐫`$ is the physical coordinate, and $`u=𝐯𝐫/r`$ is the line-of-sight component of the peculiar velocity $`𝐯`$, assuming that the observer’s peculiar velocity is zero, the apparent galaxy fluctuation $`\delta _g^s(𝐬)`$ at the apparent position $`𝐬=(1+u/r)𝐫`$ is related to the actual one $`\delta _g`$ computed at the same apparent position by the relation
$$\delta _g^s(𝐬)=\delta _g(𝐬)u^{}(𝐬)\left[u(𝐬)\left(\delta _g(𝐬)u^{}(𝐬)\right)\right]^{}.$$
(19)
Here $`u^{}`$ indicates the first radial derivative of $`u`$. Since in this section we will compute the effects of redshift distortions on the galaxy bispectrum, both for a local Eulerian and Lagrangian bias, only corrections up-to second order are considered. We remind the reader that in the distant-observer limit, the Fourier transform of $`d/drik\mu `$ where $`\mu =𝐤𝐫/kr`$ and $`\stackrel{~}{u}=i\mu f(\mathrm{\Omega })\stackrel{~}{\eta }/k`$, where $`f(\mathrm{\Omega })\mathrm{\Omega }^{0.6}`$, $`a`$ is the universal scale factor, and $`\eta `$ is the divergence of the velocity field.
### 4.1 EB galaxy bispectrum in redshift space
In this case, inserting eq.(5) into eq.(19), the Fourier transform of $`\delta _g^s(𝐬)`$ is given by
$$\stackrel{~}{\delta }_g^s=(b_1^E+\mu ^2f)\stackrel{~}{\delta }^{(1)}+^{(2)}𝒮_E^{(2)},$$
(20)
where the redshift-distorted symmetric Eulerian-bias kernel is
$`𝒮_E^{(2)}`$ $``$ $`b_1^EJ_S^{(2)}+\mu ^2fK_S^{(2)}+{\displaystyle \frac{1}{2}}b_2^E`$ (21)
$`+`$ $`{\displaystyle \frac{1}{2}}b_1^Ef\left[\mu _1^2+\mu _2^2+\mu _1\mu _2\left({\displaystyle \frac{k_1}{k_2}}+{\displaystyle \frac{k_2}{k_1}}\right)\right]`$
$`+`$ $`f^2\left[\mu _1^2\mu _2^2+{\displaystyle \frac{1}{2}}\mu _1\mu _2\left(\mu _1^2{\displaystyle \frac{k_1}{k_2}}+\mu _2^2{\displaystyle \frac{k_2}{k_1}}\right)\right].`$
The quantity $`K_S^{(2)}`$ describes the second-order contribution to $`\eta `$ (Goroff et al. 1986). The distorted galaxy bispectrum is (Heavens, Matarrese & Verde 1998)
$$B_g^s(E)=2D^4(b_1^E+\mu _1^2f)(b_1^E+\mu _2^2f)𝒮_E^{(2)}P(k_1)P(k_2)+\mathrm{c}.\mathrm{t}.$$
(22)
### 4.2 LB galaxy bispectrum in redshift space
We adopt in this case the expression in eq.(9), obtaining, after analogous calculations,
$$\stackrel{~}{\delta }_g^s=(1+b_0^L+b_1^L+\mu ^2f)\stackrel{~}{\delta }^{(1)}+^{(2)}𝒮_L^{(2)}.$$
(23)
Thus, the galaxy bispectrum is
$`B_g^s(L)`$ $`=`$ $`2D^4(1+b_0^L+b_1^L+\mu _1^2f)(1+b_0^L+b_1^L+\mu _2^2f)`$ (24)
$`\times `$ $`𝒮_L^{(2)}P(k_1)P(k_2)+\mathrm{c}.\mathrm{t}.,`$
where, in this bias prescription, the redshift-distorted second-order Lagrangian-bias kernel $`𝒮_L^{(2)}`$ is
$`𝒮_L^{(2)}`$ $``$ $`(1+b_0^L)J_S^{(2)}+_S^{(2)}+\mu ^2fK_S^{(2)}`$ (25)
$`+`$ $`{\displaystyle \frac{1}{2}}(1+b_0^L+b_1^L)f\left[\mu _1^2+\mu _2^2+\mu _1\mu _2\left({\displaystyle \frac{k_1}{k_2}}+{\displaystyle \frac{k_2}{k_1}}\right)\right]`$
$`+`$ $`f^2\left[\mu _1^2\mu _2^2+{\displaystyle \frac{1}{2}}\mu _1\mu _2\left(\mu _1^2{\displaystyle \frac{k_1}{k_2}}+\mu _2^2{\displaystyle \frac{k_2}{k_1}}\right)\right].`$
### 4.3 Comparing bias in redshift space
If we assume, once again, the validity of the algebric relation $`b_1^E=1+b_0^L+b_1^L`$, it follows that between the redshift-distorted kernels holds the relation $`𝒮_E^{(2)}=𝒮_L^{(2)}+b_1^LJ_S^{(2)}+\frac{1}{2}b_2^E_S^{(2)}.`$ This relation should be immediately compared with the one in eq.(13), to understand that we obtain the concise expression between the quantities $`\mathrm{\Delta }B_g^s`$ and $`\mathrm{\Delta }B_g`$ which emphasize the inconsistency between the two biasing prescriptions,
$$\mathrm{\Delta }B_g^s=(1+\mu _1^2\beta )(1+\mu _2^2\beta )\mathrm{\Delta }B_g(1,2)+\mathrm{c}.\mathrm{t}.,$$
(26)
where $`\beta f/(1+b_0^L+b_1^L)`$. Thus, the only redshift effect on the quantity $`\mathrm{\Delta }B_g(i,j)`$ comes from the first-order distortion of the galaxy number density field, $`\stackrel{~}{\delta }_g^{s(1)}=(1+\mu ^2\beta )\stackrel{~}{\delta }_g^{(1)}`$. It has to be like that, if one thinks that the distortion effects due to peculiar motions are either independent of the bias factors or proportional to the first-order bias factors, then they cancel out. In redshift space, eq.(16) becomes
$$Q_g^s(E)=Q_g^s(L)+\frac{\mathrm{\Delta }B_g^s}{P_g^sP_g^s}.$$
(27)
Though the structure of the expression (16) is preserved in redshift space, $`Q_g(Q_m)`$-relations like those in eqs.(17) and (18) are not. A comprehensive investigation of the effects of redshift distortions on $`B_g^E`$ and $`B_g^L`$ is in progress; see also Scoccimarro et al. (1999) for an analysis of $`B_g^s(E)`$.
## 5 Discussion and conclusions
We compared the galaxy clustering predictions of the local Eulerian bias scheme versus those of the Lagrangian one. We showed that the two bias models are inconsistent, since the predicted three-point galaxy correlations are different. A similar inconsistency certainly characterizes correlations of higher order, or of lower order but higher perturbative corrections. Qualitatively, these results are independent on whether the Lagrangian zero-order bias factor $`b_0^L`$ is zero, as for Press-Schechter dark matter halos, or not, as in the most general case we have considered here. The galaxy bispectrum is much better suited to distinguish between the two bias models than the corresponding skewness, since the latter is spatially averaged: the bispectrum depends on the shape of the triangle $`𝐤_1+𝐤_2+𝐤_3=\mathrm{𝟎}`$, thus two shapes can disentagle the two bias factors $`b_1`$ and $`b_2`$ (Matarrese, Verde & Heavens 1997; Scoccimarro 2000) and the two bias models. The shape dependence is best quantified by the $`Q`$-amplitudes discussed in Subsection 3.3.: the reduced amplitude of the bispectrum from Lagrangian bias in eq.(18) displays a dependence on the triangle configuration which is not contained in eq.(17), and, since it is dynamically induced, which is independent from the bias factors. The next generation redshift catalogues, as the ongoing Two Degree Field Survey and the Sloan Digital Sky Survey, will contain enough galaxies to establish whether $`B_g^s(E)`$ in eq.(22) or $`B_g^s(L)`$ in eq.(24) better fits the observational data, but they cannot be both correct, whatever the assumed cosmology.
Both bias schemes represent rather extreme and idealized approaches. Lagrangian models imply a sort of infinite-memory process, since the sites for galaxy formation are known from the beginning, and dynamical evolution changes their spatial distribution. On the other hand, in local Eulerian schemes galaxies are simply ‘painted’ on a snapshot of the density field, without a record of the past. However, even though real galaxy formation is probably a process with intermediate characteristics with respect to the biasing schemes discussed here, recent models based on a Lagrangian selection of the sites for object formation were shown to be very successful in reproducing the clustering of dark-matter halos found in numerical simulations (e.g. Catelan, Matarrese & Porciani 1998; Porciani, Catelan & Lacey 1999). The issue discussed in this Letter surely deserves further investigation, both in real and redshift space. It would be of interest to test which biasing scheme better describes galaxy power spectrum and bispectrum from a combination of numerical simulations and semianalytic models (Porciani et al., in preparation).
Acknowledgments We thank Sabino Matarrese for discussions. The anonymous referee improved the presentation of these results. CP is supported by a Golda Meir Fellowship. This work was supported at Caltech in part by the DoE, NSF, and NASA.
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# The Pauli Equation for Probability Distributions
## I Introduction
Since the early days of quantum mechanics, we have been forced to coexist with complex probability amplitudes without worrying about their lack of any reasonable physical meaning. One should not ignore, however, that the wave-like properties of quantum objects still raise conceptual problems on whose solutions, a general consensus is far from having been reached .
A possible way out of this difficulty has been implicitly suggested by Feynmann , who has shown that, by dropping the assumption that the probability for an event must always be nonnegative, one can avoid the use of probability amplitudes in quantum mechanics. This proposal goes back to the work by Wigner , who first introduced nonpositive pseudoprobabilities to represent quantum mechanics in phase space, and to the Moyal approach to quantum mechanics .
From a conceptual point of view, the elimination of the waves from quantum theory is in line with the procedure inaugurated by Einstein with the elimination of the aether in the theory of electromagnetism.
The phase-space formulation of quantum mechanics provides a means of analyzing quantum-mechanical systems while still employing a classical framework. Moreover, a quantum mechanics without wave functions has been discussed in .
Recentely, the problem of quantum state measurement, initially posed by Pauli , received a lot of attention . The tomographic approach to the quantum state of a system has allowed to establish a map between the density operator (or any its representation) and a set of probability distributions, often called “marginals”. The latter have all the characteristics of classical probabilities; they are nonnegative, measurable, and normalized.
Based on this connection, a classical-like description of quantum dynamics by means of “symplectic tomography” has been formulated , providing a bridge between classical and quantum worlds. That is, the evolution of a quantum system with continuous observables (namely, quadrature components of a field mode) was described in terms of a classical-like equation for a marginal distribution. Different aspects of this classical-like description using tomographic probabilities were recently analyzed .
On the other hand, discrete observables, like spin or angular momentum, are so important in quantum mechanics as the continuous ones are. Hence, the tomography scheme for discrete variables was introduced , and the marginal distribution for rotated spin variables has been constructed , deriving an evolution equation for this function.
Here, we would extend the approach by considering a spin-$`\frac{1}{2}`$ particle moving in a potential, then constructing the marginal distributions for space coordinates and spin projections and finally deriving the evolution equation for such probabilities, which would be the analog of the Pauli equation. It would also be a generalization of approaches attempted by us in previous papers .
Essentially, our aim is to eliminate the hybrid procedure of describing the dynamical evolution of a system, which consists of a first stage, where the theory provides a deterministic evolution of the wave function, followed by a hand-made construction of the physically meaningful probability distributions. If the probabilistic nature of the microscopic phenomena is fundamental, and not simply due to our ignorance, as in classical statistical mechanics, why should it be impossible to describe them in probabilistic terms from the very beginning? On the other hand, the language of probability, suitably adapted to take into account all the relevant constraints, seems to be the only language capable of expressing the fundamental role of “chance” in nature .
The paper is organized as follow:
In Section 2, we review the general approach to construct known tomography schemes using density matrix in the specifically transformed reference frames. In Section 3, we derive the general evolution equation for tomographic probabilities (marginal distributions) which describe the quantum state instead of density matrix. In Section 4, the general scheme of tomography construction is used to re-derive the particular example of symplectic tomography, which is applied for measuring states depending on continuous quadrature. In Section 5, the general scheme is used to re-derive the construction of spin-state tomography. In Section 6, the general scheme of Section 2 is then applied to obtain tomographic probabilities in the combined situation described by spatial (multidimensional too) and spin variables. In Section 7, some examples are studied in the context of the probability representation of quantum mechanics. Section 8 concludes.
We are using the natural unit ($`\mathrm{}=c=1`$).
## II General approach to quantum tomography
In this section, we give a short review of the general principles used to construct a tomography scheme for measuring quantum states. Recently, we established a quite general principle of constructing measurable probabilities, which determine completely the quantum state in the tomographic approach; more refined treatments then followed . Here, we apply our general approach to derive the evolution equation for the tomographic probabilities that is alternative in some sense to the Schrödinger equation for the wave function (or the quantum Liouville equation for the density matrix).
Let us consider a quantum state described by the density operator $`\widehat{\rho }`$, which is a nonnegative hermitian operator, i.e.,
$$\widehat{\rho }^{}=\widehat{\rho },\text{Tr}\widehat{\rho }=1,$$
(1)
and
$$v\widehat{\rho }v=\rho _{v,v}0.$$
(2)
We label the vector basis $`v`$ in the space of pure quantum states by the multidimensional index $`v=(v_1,v_2,\mathrm{},v_N),`$ where the number $`N`$ shows the number of degrees of freedom of the system under consideration. Among indexes $`v_k`$, $`k=1,\mathrm{},N`$, there are continuous ones like position (or momentum) and discrete ones like spin projections. In this sense, the wave function $`\psi (v)=v\psi `$ of a pure state $`\psi `$ depends both on continuous and discrete observables. Formula (2) can be rewritten by using the hermitian projection operator
$$\widehat{\mathrm{\Pi }}_v=vv,$$
(3)
in the following form
$$\rho _{v,v}=\text{Tr}\left\{\widehat{\mathrm{\Pi }}_v\widehat{\rho }\right\}.$$
(4)
The physical meaning of the projector $`\widehat{\mathrm{\Pi }}_v`$ is that it extracts the state $`v`$ with given $`v`$ (for example, with given position and spin projection), which is an eigenstate of the-commuting-hermitian operators $`\widehat{V}=(\widehat{V}_1,\widehat{V}_2,\mathrm{},\widehat{V}_N)`$
$$\widehat{V}_kv=v_kv.$$
(5)
In the space of states, there is a family of unitary transformation operators $`\widehat{U}(\sigma )`$ depending on the parameters $`\sigma =(\sigma _1,\mathrm{},\sigma _k\mathrm{})`$, that can be sometimes identified with a group-representation operators. In these cases, the parameters $`\sigma `$ describe the group element. It was shown that known tomography schemes can be considered from the viewpoint of the group theory by using appropriate groups. More recentely this concept has been developed obtaining an elegant group theoretical approach to quantum state measurement . Here, we formulate the tomographic approach in the following way. Let us introduce a “transformed density operator”
$$\widehat{\rho }_\sigma =\widehat{U}^1(\sigma )\widehat{\rho }\widehat{U}(\sigma ).$$
(6)
Its diagonal elements are still nonnegative probabilities
$$z\widehat{\rho }_\sigma z=z\widehat{\rho }zw(z,\sigma ).$$
(7)
Here, $`z`$ is one of the possible vectors $`v`$, while the symbol $`z`$ denotes the transformed vectors
$$z=\widehat{U}(\sigma )z,$$
(8)
which in turn are eigenstates of the-transformed-operators
$$\widehat{Z}=\widehat{U}(\sigma )\widehat{V}\widehat{U}^1(\sigma ).$$
(9)
As a consequence of the unit trace of the density operator we also have the normalization condition
$$𝑑zw(z,\sigma )=1.$$
(10)
Of course, in case of discrete indices, the integral in Eq. (10) is replaced by a sum over discrete variables.
Formula (7) can be interpreted as the probability density for the measurement of the observable $`\widehat{V}`$ in an ensemble of transformed reference frames labeled by the index $`\sigma `$, if the state $`\widehat{\rho }`$ is given. Along with this interpretation, one can also consider the transformed projector
$$\widehat{\mathrm{\Pi }}_z(\sigma )=\widehat{U}(\sigma )\widehat{\mathrm{\Pi }}_z\widehat{U}^1(\sigma )=zz,$$
(11)
the explicit expression for the probability $`w(z,\sigma )`$ takes the form
$$w(z,\sigma )=\text{Tr}\left\{\widehat{\rho }\widehat{\mathrm{\Pi }}_z(\sigma )\right\}=\text{Tr}\{\widehat{\rho }zz\}.$$
(12)
These probability densities are also called “marginal” distributions as generalization of the concept introduced by Wigner . The tomography schemes are based on the possibility to find the inverse of Eq. (12). If it is possible to solve Eq. (12), considering the probability $`w(y,\sigma )`$ as known function and the density matrix as unknown operator, the quantum state can be described by the positive probability instead of the density matrix. This property is the essense of state reconstruction techniques. In such cases, the inverse of Eq. (12) takes the form
$$\widehat{\rho }=w(z,\sigma )\widehat{K}(z,\sigma )𝑑z𝑑\sigma .$$
(13)
Thus, there exist a family of operators $`\widehat{K}(z,\sigma )`$ depending on both the variables $`z`$ and parameters $`\sigma `$ such that the density operator is reconstructed, if the probability $`w(z,\sigma )`$ is known. It is worth remarking that transformations $`\widehat{U}(\sigma )`$ can form other algebraic constructions, which have no structure of groups . The only condition for the existence of a tomography scheme is the possibility to invert Eq. (12). In the cases of optical tomography , symplectic tromography , and spin tomography , the sets of transformations $`\widehat{U}(\sigma )`$ have the structure of corresponding Lie groups (i.e., rotation, symplectic and spin).
## III The time evolution equation
We are now interested in finding the evolution equation for the probability $`w(z,\sigma ,t)`$, in which $`t`$ is the time parameter. Using Eq. (12) one has
$$_tw(z,\sigma ,t)=\text{Tr}\left\{\left[_t\widehat{\rho }(t)\right]\widehat{\mathrm{\Pi }}_z(\sigma )\right\}.$$
(14)
On the other hand, the density operator satisfies the Liouville-Von Neumann equation
$$_t\widehat{\rho }(t)=i[\widehat{\rho }(t),\widehat{H}],$$
(15)
with $`\widehat{H}`$ the system Hamiltonian. By inserting Eq.(15) in (14), and with the aid of Eq.(13), we find the evolution equation for the probability $`w`$ in a closed form
$$_tw(z,\sigma ,t)=𝑑z^{}𝑑\sigma ^{}w(z^{},\sigma ^{},t)\text{Tr}\left\{i[\widehat{K}(z^{},\sigma ^{}),\widehat{H}]\widehat{\mathrm{\Pi }}_z(\sigma )\right\}.$$
(16)
Equation (16) represents the classical-like version of the Liouville-Von Neumann equation, thus, it would be the analog of the Pauli equation for a system with space and spin degrees of freedom.
## IV Quantum tomography with continuous variables
Let us consider, in a one-dimensional system, an operator $`\widehat{X}`$ as the linear combination of position $`\widehat{q}`$ and momentum $`\widehat{p}`$
$$\widehat{X}=\mu \widehat{q}+\nu \widehat{p},$$
(17)
which depends upon real parameters $`\mu `$, $`\nu `$ and, due to its hermiticity, is a measurable observable. Since the linear canonical transformation (17) belongs to the symplectic group $`Sp(2,R)`$, the tomography scheme under discussion was called “symplectic tomography” .
The probability (marginal) related to the observable (17) is given by
$$w(x,\mu ,\nu )=x\widehat{\rho }x,$$
(18)
where $`\widehat{\rho }`$ is the system’s density operator, while the eigenstates $`x`$ of the operator (17) can be written as
$$x=dqqxq,$$
(19)
with $`q`$ the position eigenkets. The wave function $`qx`$ can be easily calculated by using the following equality
$$q\widehat{X}x=q\mu \widehat{q}+\nu \widehat{p}x,$$
(20)
and then transforming it in a partial differential equation
$$xqx=\mu qqxi\nu \frac{}{q}qx.$$
(21)
The solution is
$$qx=\left(\frac{1}{|\nu |}\right)^{1/2}\mathrm{exp}\left[i\frac{x}{\nu }q\frac{i}{2}\frac{\mu }{\nu }q^2\right].$$
(22)
It is worth noting that as soon as $`\mu 1`$ and $`\nu 0`$, then $`xx`$ and the wavefunction (22) tends to $`\delta (qx)`$.
Furthermore, Eq.(18) can be formally rewritten as
$$w(x,\mu ,\nu )=\mathrm{Tr}\left\{\widehat{\rho }\widehat{\mathrm{\Pi }}_x(\mu ,\nu )\right\},$$
(23)
where the transformed projector is given by
$$\widehat{\mathrm{\Pi }}_x(\mu ,\nu )=\widehat{U}(\mu ,\nu )\widehat{\mathrm{\Pi }}_x\widehat{U}^1(\mu ,\nu ),\widehat{\mathrm{\Pi }}_x=xx.$$
(24)
Here, the transformation $`\widehat{U}(\sigma )`$ is chosen to be the symplectic group representation
$$\widehat{U}(\mu ,\nu )=\mathrm{exp}\left[i\varphi \left(\frac{\widehat{p}^2}{2}+\frac{\widehat{q}^2}{2}\right)\right]\mathrm{exp}\left[\frac{i\lambda }{2}\left(\widehat{q}\widehat{p}+\widehat{p}\widehat{q}\right)\right].$$
(25)
The rotation and scaling parameters $`\varphi `$ and $`\lambda `$ are related to $`\mu `$ and $`\nu `$ by the following formulas
$`\mu =\lambda \mathrm{cos}\varphi ,`$ $`\nu =\lambda ^1\mathrm{sin}\varphi ,`$ (26)
$`\varphi ={\displaystyle \frac{1}{2}}\text{arcsin}\left(2\mu \nu \right),`$ $`\lambda =\pm {\displaystyle \frac{1}{4}}\sqrt{{\displaystyle \frac{1\pm \sqrt{14\mu ^2\nu ^2}}{2}}}.`$ (28)
This means that the marginal distribution $`w(x,\mu ,\nu )`$ for this particular case of symplectic tomography is given by the relation
$`w(x,\mu ,\nu )`$ $`=`$ $`\text{Tr}\{xx\mathrm{exp}\left[i\varphi ({\displaystyle \frac{\widehat{p}^2}{2}}+{\displaystyle \frac{\widehat{q}^2}{2}})\right]\mathrm{exp}\left[{\displaystyle \frac{i\lambda }{2}}(\widehat{q}\widehat{p}+\widehat{p}\widehat{q})\right]`$ (30)
$`\times \widehat{\rho }\mathrm{exp}[i\varphi ({\displaystyle \frac{\widehat{p}^2}{2}}+{\displaystyle \frac{\widehat{q}^2}{2}})]\mathrm{exp}[{\displaystyle \frac{i\lambda }{2}}(\widehat{q}\widehat{p}+\widehat{p}\widehat{q})]\}.`$
Such measurable probability can be explicitly expressed as
$$w(x,\mu ,\nu )=𝑑y𝑑k\mathrm{exp}\left[ikx+\frac{i\mu \nu k^2}{2}+iky\mu \right]\rho (y+\nu k,y),$$
(31)
where $`\rho (y+\nu k,y)=y+\nu k|\widehat{\rho }|y`$ is the representation of the density matrix over the position eigenkets. The marginal satisfy the following homogeneous property
$`w(x,\mu ,\kappa \nu )`$ $`=`$ $`{\displaystyle \frac{1}{\kappa }}w(x/\kappa ,\mu /\kappa ,\nu ),`$ (32)
$`w(x,\kappa \mu ,\nu )`$ $`=`$ $`{\displaystyle \frac{1}{\kappa }}w(x/\kappa ,\mu ,\nu /\kappa ).`$ (34)
The above relation (31) can be inverted as
$$\widehat{\rho }=𝑑x𝑑\mu 𝑑\nu w(x,\mu ,\nu )\widehat{K}(x,\mu ,\nu ),$$
(35)
where the kernel operator takes the form
$$\widehat{K}(x,\mu ,\nu )=\frac{1}{2\pi }ϵ^2\mathrm{exp}\left[iϵX+\frac{iϵ^2\mu \nu }{2}\right]e^{iϵ\mu \widehat{q}}e^{iϵ\nu \widehat{p}}.$$
(36)
Here, $`ϵ`$ can be set equal 1; this freedom reflects the overcompleteness of information obtainable by means of all possible marginals (30.
The multi-mode generalization is straightforward, and the analog of formula (30) holds with the following replacement
$`x`$ $``$ $`\stackrel{}{x},\stackrel{}{x}=(x_1,x_2,\mathrm{}),`$ (37)
$`\varphi \left({\displaystyle \frac{\widehat{p}^2}{2}}+{\displaystyle \frac{\widehat{q}^2}{2}}\right)`$ $``$ $`\varphi _1\left({\displaystyle \frac{\widehat{p}_1^2}{2}}+{\displaystyle \frac{\widehat{q}_1^2}{2}}\right)+\varphi _2\left({\displaystyle \frac{\widehat{p}_2^2}{2}}+{\displaystyle \frac{\widehat{q}_2^2}{2}}\right)+\mathrm{},`$ (38)
$`\lambda \left(\widehat{q}\widehat{p}+\widehat{p}\widehat{q}\right)`$ $``$ $`\lambda _1\left(\widehat{q}_1\widehat{p}_1+\widehat{p}_1\widehat{q}_1\right)+\lambda _2\left(\widehat{q}_2\widehat{p}_2+\widehat{p}_2\widehat{q}_2\right)+\mathrm{}.`$ (39)
Relations of the parameters $`\lambda _k,\varphi _k`$ to the parameters $`\mu _k,\nu _k`$ are the same of Eq. (LABEL:param).
## V Quantum tomography with discrete variables
Here, we consider a spin-$`j`$ system. Following we will derive the expression for the density matrix of a spin state in terms of measurable probability distributions.
For arbitrary values of spin, let the spin state have the density matrix
$$\rho _{mm^{}}^{(j)}=jm\widehat{\rho }^{(j)}jm^{},m=j,j+1,\mathrm{},j1,j,$$
(40)
where
$$\widehat{j}_3jm=mjm,\widehat{j}^2jm=j(j+1)jm,$$
(41)
and
$$\widehat{\rho }^{(j)}=\underset{m=j}{\overset{j}{}}\underset{m^{}=j}{\overset{j}{}}\rho _{mm^{}}^{(j)}jmjm^{}.$$
(42)
The operator $`\widehat{\rho }^{(j)}`$ is the density operator of the state under consideration.
The general group construction of tomographic schemes was also used for spin tomography . The idea is to consider the diagonal elements of the density matrix $`\widehat{\rho }`$ in another reference frame, i.e. rotated one. To this end we introduce a rotated measureble spin projection
$$\widehat{J}_3(\alpha ,\beta ,\gamma )=\widehat{D}(\alpha ,\beta ,\gamma )\widehat{j}_3\widehat{D}^1(\alpha ,\beta ,\gamma ),$$
(43)
where the unitary rotation operator $`\widehat{D}`$ depends on the Euler angles $`\alpha ,\beta ,\gamma `$. The role of the observable $`\widehat{Z}`$ is now played by the spin projection $`\widehat{J}_3`$, while the rotation-transformation parameters are the Euler angles $`\sigma _1=\alpha ,`$ $`\sigma _2=\beta `$, $`\sigma _3=\gamma `$. The transformation $`\widehat{U}(\sigma )`$ is given by the matrix representation of the rotation group, i.e., the Wigner $`D`$-function .
The marginals are
$$w(s,\alpha ,\beta ,\gamma )=js\widehat{\rho }js,$$
(44)
where the rotated spin states becomes
$$js=\underset{m=j}{\overset{j}{}}D_{sm}^{(j)}(\alpha ,\beta ,\gamma )jm.$$
(45)
Here the matrix elements $`D_{m^{}m}^{(j)}(\alpha ,\beta ,\gamma )`$ (Wigner $`D`$-functions) are the matrix elements of the rotation-group representation
$$D_{m^{}m}^{(j)}(\alpha ,\beta ,\gamma )=e^{im^{}\gamma }d_{m^{}m}^{(j)}(\beta )e^{im\alpha },$$
(46)
where
$$d_{m^{}m}^{(j)}(\beta )=\left[\frac{(j+m^{})!(jm^{})!}{(j+m)!(jm)!}\right]^{1/2}\left(\mathrm{cos}\frac{\beta }{2}\right)^{m^{}+m}\left(\mathrm{sin}\frac{\beta }{2}\right)^{m^{}m}P_{jm^{}}^{(m^{}m,m^{}+m)}(\mathrm{cos}\beta ),$$
(47)
with $`P_n^{(a,b)}(x)`$ the Jacobi polynomials .
Moreover, the transformed spin projector will be
$$\widehat{\mathrm{\Pi }}_s(\alpha ,\beta ,\gamma )=\widehat{D}(\alpha ,\beta ,\gamma )jsjs\widehat{D}^1(\alpha ,\beta ,\gamma )=jsjs,$$
(48)
then, we have
$$w(s,\alpha ,\beta ,\gamma )=\underset{m_1=j}{\overset{j}{}}\underset{m_2=j}{\overset{j}{}}D_{sm_1}^{(j)}(\alpha ,\beta ,\gamma )\rho _{m_1m_2}^{(j)}D_{sm_2}^{(j)}(\alpha ,\beta ,\gamma ).$$
(49)
Since
$$D_{m^{}m}^{(j)}(\alpha ,\beta ,\gamma )=(1)^{m^{}m}D_{m^{}m}^{(j)}(\alpha ,\beta ,\gamma ),$$
(50)
the marginal distribution really depends only on two angles, $`\alpha `$ and $`\beta `$. Hence
$$w(s,\alpha ,\beta ,\gamma )w(s,\alpha ,\beta ),$$
(51)
which satisfies the normalization condition
$$\underset{s=j}{\overset{j}{}}w(s,\alpha ,\beta )=1.$$
(52)
As an example, for a spin-$`\frac{1}{2}`$ state with spin projection $`+1/2`$, we have
$$\widehat{\rho }=\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right),$$
(53)
and the marginal distributions will be
$$w\left(s=\frac{1}{2},\alpha ,\beta \right)=\mathrm{cos}^2\frac{\beta }{2},w\left(s=\frac{1}{2},\alpha ,\beta \right)=\mathrm{sin}^2\frac{\beta }{2}.$$
(54)
In Refs. , in view of the properties of the Wigner $`D`$-function and the Clebsch–Gordan coefficients, Eq. (49) was inverted and the density matrix was expressed in terms of the marginal distribution
$`\rho _{m_1m_2}^{(j)}`$ $`=`$ $`(1)^{m_2}{\displaystyle \underset{j_3=0}{\overset{2j}{}}}{\displaystyle \underset{m_3=j_3}{\overset{j_3}{}}}(2j_3+1)^2{\displaystyle \underset{s=j}{\overset{j}{}}}{\displaystyle (1)^sw(s,\alpha ,\beta )}`$ (55)
$`\times `$ $`D_{0m_3}^{(j_3)}(\alpha ,\beta ,\gamma )W_{ss\mathrm{\hspace{0.17em}0}}^{jjj_3}W_{m_1m_2m_3}^{jjj_3}{\displaystyle \frac{d\mathrm{\Omega }}{8\pi ^2}}`$ (56)
where $`m_1,m_2=j,j+1,\mathrm{},j`$ and $`W_{m_1m_2m_3}^{j_1j_2j_3}`$ are the Wigner-$`3j`$ symbols . The integration is performed over the rotation parameters, i.e.
$$𝑑\mathrm{\Omega }=_0^{2\pi }𝑑\alpha _0^\pi \mathrm{sin}\beta d\beta _0^{2\pi }𝑑\gamma .$$
(57)
Equation (55) can be presented in an invariant operator form . We systematically introduce the following notation, first for the function on the unit sphere
$$\mathrm{\Phi }_{jm_1m_2}^{(j_3)}(\alpha ,\beta )=(1)^{m_2}\underset{m_3=j_3}{\overset{j_3}{}}D_{0m_3}^{(j_3)}(\alpha ,\beta ,\gamma )W_{m_1m_2m_3}^{jjj_3},$$
(58)
and then for the operator on the unit sphere
$$\widehat{A}_j^{(j_3)}(\alpha ,\beta )=(2j_3+1)^2\underset{m_1=j}{\overset{j}{}}\underset{m_2=j}{\overset{j}{}}jm_1\mathrm{\Phi }_{jm_1m_2}^{(j_3)}(\alpha ,\beta )jm_2.$$
(59)
In order to write a final expression for the density operator, we introduce an operator on the unit sphere which contains a dependence on the measurable projection of the spin
$$\widehat{K}^{(j)}(s,\alpha ,\beta )=(1)^s\underset{j_3=0}{\overset{2j}{}}W_{ss\mathrm{\hspace{0.17em}0}}^{jjj_3}\widehat{A}_j^{(j_3)}(\alpha ,\beta ).$$
(60)
Finally, we obtain a compact expression for the density operator,
$$\widehat{\rho }^{(j)}=\underset{s=j}{\overset{j}{}}\frac{d\mathrm{\Omega }}{8\pi ^2}w(s,\alpha ,\beta )\widehat{K}^{(j)}(s,\alpha ,\beta ).$$
(61)
Formula (61) admits of the following interpretation. To determine the spin state for a spin $`j`$, one has to experimentally measure the projection $`s`$ of the spin for each direction specified by the angles $`\alpha `$ and $`\beta `$, obtaining a distribution function $`w(s,\alpha ,\beta )`$. The sum on the r.h.s. of Eq.(61) for a given point on the unit sphere represents the average operator $`\widehat{K}^{(j)}(s,\alpha ,\beta )`$. Then, the integral over the whole solid angle gives the desired density operator. Finally, we recognize that, for the spin case, the operator (60) plays the role of the operator $`\widehat{K}(z,\sigma )`$ of Eq. (13), employed in the general scheme of Section 2.
## VI The general case
We are now able to consider the case of a particle with $`N1`$ spatial degrees of freedom, plus one spin-$`\frac{1}{2}`$ degree. In this case, the state vector $`v`$ has the form
$$\stackrel{}{q},m=q_1,\mathrm{}q_{N1}\frac{1}{2},m,$$
(62)
where $`\stackrel{}{q}`$ is the eigenvalue of the position operator $`\widehat{\stackrel{}{q}}`$ and the spin projection $`m=(1/2,\mathrm{\hspace{0.17em}1}/2)`$ is the eigenvalue of the Pauli matrix $`\widehat{\sigma }_z`$.
The transformation operator $`\widehat{U}(\sigma )`$ used to construct the tomography scheme, for this case, depends on $`2(N1)`$ parameters determining the symplectic transform, and on three Euler angles determining the spin rotation.
The transformation operator $`\widehat{U}(\sigma )`$ of Eq. (6) becomes the product of operators
$$\widehat{U}(\sigma )=_{k=1}^{N1}\widehat{U}(\mu _k,\nu _k)\widehat{U}(\alpha ,\beta ,\gamma ).$$
(63)
For the case of spin-$`\frac{1}{2}`$, the representation of the rotation group is given by
$$D(\alpha ,\beta ,\gamma )=\left(\begin{array}{cc}e^{i\alpha /2}\mathrm{cos}\left(\beta /2\right)e^{i\gamma /2}& e^{i\alpha /2}\mathrm{sin}\left(\beta /2\right)e^{i\gamma /2}\\ e^{i\alpha /2}\mathrm{sin}\left(\beta /2\right)e^{i\gamma /2}& e^{i\alpha /2}\mathrm{cos}\left(\beta /2\right)e^{i\gamma /2}\end{array}\right),$$
(64)
which determines the operator
$$\widehat{U}(\alpha ,\beta ,\gamma )=\underset{m_1=1/2}{\overset{1/2}{}}\underset{m_2=1/2}{\overset{1/2}{}}D_{m_1m_2}^{(1/2)}(\alpha ,\beta ,\gamma )\frac{1}{2},m_1\frac{1}{2},m_2.$$
(65)
The marginal distribution $`w(z,\sigma )`$ (12) will depend on $`N1`$ continuous (noncompact) variables $`z_1=x_1,`$ $`\mathrm{},`$ $`z_{N1}=x_{N1},`$ and one discrete spin projection $`z_N=s,`$ as well as on parameters $`\mu _k,\nu _k`$ and on Euler angles $`\alpha ,\beta .`$ The dependence of the marginal distribution on the Euler angle $`\gamma `$ disappears, as it was shown in the previous section, due to the structure of Wigner $`D`$-functions.
In order to get an analog of the Pauli evolution equation for the marginal distribution, we consider the general equation (16) where the operator $`\widehat{K}(z^{},\sigma ^{})`$ has the form
$$\widehat{K}(z^{},\sigma ^{})=\frac{1}{8\pi ^2}_{k=1}^{N1}\widehat{K}(x_k,\mu _k\nu _k)\widehat{K}^{(1/2)}(s,\alpha ,\beta ).$$
(66)
Here, the operator $`\widehat{K}(x_k,\mu _k,\nu _k)`$ has the form of Eq.(36) with $`ϵ=1,`$ and the operator $`\widehat{K}^{(1/2)}(s,\alpha ,\beta )`$ is given by formula (60) with $`j=1/2`$. Moreover, we have to introduce the marginal distribution $`w(\stackrel{}{x},\stackrel{}{\mu },\stackrel{}{\nu },s,\alpha ,\beta ,t)`$ describing a state of spin-$`\frac{1}{2}`$ particle which depends on the continuous variables $`\stackrel{}{x}`$, discrete spin projection $`s`$, symplectic reference frame’s labels $`\stackrel{}{\mu }`$ and $`\stackrel{}{\nu }`$, and Euler angles $`\alpha `$ and $`\beta `$. Then, for a given Hamiltonian $`\widehat{H}`$ the general equation (16) takes the form of a Pauli-like equation equation
$`_tw(\stackrel{}{x},\stackrel{}{\mu },\stackrel{}{\nu },s,\alpha ,\beta ,t)`$ $`=`$ $`{\displaystyle \underset{s^{}=1/2}{\overset{1/2}{}}}{\displaystyle 𝑑\stackrel{}{X}^{}𝑑\stackrel{}{\mu }^{}𝑑\stackrel{}{\nu }^{}𝑑\mathrm{\Omega }^{}w(\stackrel{}{x}^{},\stackrel{}{\mu }^{},\stackrel{}{\nu }^{},s^{},\alpha ^{},\beta ^{},t)}`$ (67)
$`\times `$ $`\mathrm{\Theta }(\stackrel{}{x},\stackrel{}{\mu },\stackrel{}{\nu },s,\alpha ,\beta ;\stackrel{}{x}^{},\stackrel{}{\mu }^{},\stackrel{}{\nu }^{},s^{},\alpha ^{},\beta ^{}),`$ (68)
where
$$\mathrm{\Theta }=\frac{i}{8\pi ^2}\stackrel{}{x},s[_{k=1}^{N1}\widehat{K}(x_k^{},\mu _k^{}\nu _k^{})\widehat{K}^{(1/2)}(s^{},\alpha ^{},\beta ^{}),\widehat{H}]\stackrel{}{x},s.$$
(69)
The structure of the derived Pauli-like equation for probability distributions depends on the particular tomography schemes we have considered. Obviously, it would be useful to find the schemes which give the simplest form for such dynamical equation, nevertheless this is a nontrivial problem related to the possibility of finding properly transformed projector (11). The latter are investigated in Ref. , but for different purposes.
### A Limit cases
We want now to consider two limiting cases of the above equation (67).
First of all we consider the (one-dimensional) spatial case only with free motion
$$\widehat{H}=\frac{\widehat{p}^2}{2}.$$
(70)
The spin part does not contribute since $`\widehat{H}`$ does not contain the spin operators, that is
$`{\displaystyle \frac{d\mathrm{\Omega }^{}}{8\pi ^2}w(s^{},\alpha ^{},\beta ^{},)s\widehat{K}^{(j)}(s^{},\alpha ^{},\beta ^{})s}={\displaystyle \underset{m_1,m_2=j}{\overset{j}{}}}D_{sm_1}^{(j)}(\alpha ,\beta ,\gamma )D_{sm_2}^{(j)}(\alpha ,\beta ,\gamma )`$ (71)
$`\times {\displaystyle \underset{j_3=0}{\overset{2j}{}}}{\displaystyle \underset{m_3=j_3}{\overset{j_3}{}}}{\displaystyle \underset{s^{}=j}{\overset{j}{}}}()^{m_2s^{}}(2j_3+1)^2W_{s^{}s^{}\mathrm{\hspace{0.17em}0}}^{jjj_3}W_{m_1m_2m_3}^{jjj_3}`$ (72)
$`\times {\displaystyle }{\displaystyle \frac{d\mathrm{\Omega }^{}}{8\pi ^2}}w(s^{},\alpha ^{},\beta ^{},)D_{0m_3}^{(j_3)}(\alpha ^{},\beta ^{},\gamma ^{})=w(s,\alpha ,\beta ,),`$ (73)
where $``$ indicates other possible variables. Then, for what concerns the spatial part, it is important to calculate the commutator between the kernel and the Hamiltonian, given by
$$[e^{i\mu ^{}\widehat{q}}e^{i\nu ^{}\widehat{p}},\widehat{p}^2]=e^{i\mu ^{}\widehat{q}}e^{i\nu ^{}\widehat{p}}(2\mu ^{}\widehat{p}\mu ^2).$$
(74)
Now, one can write
$`_tw(x,\mu ,\nu ,t)`$ $`=`$ $`{\displaystyle \frac{i}{4\pi }}{\displaystyle 𝑑x^{}𝑑\mu ^{}𝑑\nu ^{}w(X^{},\mu ^{},\nu ^{},t)e^{iX^{}+i\mu ^{}\nu ^{}/2}}`$ (75)
$`\times `$ $`{\displaystyle 𝑑qxe^{i\mu ^{}\widehat{q}}e^{i\nu ^{}\widehat{p}}|qq|(2\mu ^{}\widehat{p}\mu ^2)x}`$ (76)
By using the explicit form for the wave functions $`qx`$ (22), toghether with the homogeneous property (LABEL:hom), it is possible to reduce the above equation to a very simple form
$$_tw=\mu _\nu w$$
(77)
which was derived in a different way in Ref. .
As a second case we study the dynamics of spin-$`\frac{1}{2}`$ degree only. The Hamiltonian we wish to consider is
$`\widehat{H}=\left(\begin{array}{cc}a& \hfill 0\\ 0& \hfill c\end{array}\right).`$ (80)
Of course, the spatial degree is not affected, so its variables can be disregarded; this also results from the fact that
$$\frac{1}{2\pi }𝑑x^{}𝑑\mu ^{}𝑑\nu ^{}w(x^{},\mu ^{},\nu ^{},)xe^{i\mu ^{}\widehat{q}}e^{i\nu ^{}\widehat{p}}xe^{ix^{}+i\mu ^{}\nu ^{}/2}=w(x,\mu ,\nu ,)$$
(81)
In this case the relation between the transformed spin state projection and the untransformed one is given by
$$s=\widehat{D}_{s,1/2}^{(1/2)}(\alpha ,\beta )|\frac{1}{2}+\widehat{D}_{s,1/2}^{(1/2)}(\alpha ,\beta )|\frac{1}{2}.$$
(82)
Again, the central task is the calculation of the commutator between the kernel and the Hamiltonian. It is easy to see that
$`s[\widehat{K}^{(1/2)}(s^{},\alpha ^{},\beta ^{}),\widehat{H}]s=(1)^s^{}{\displaystyle \underset{j_3=0}{\overset{1}{}}}W_{s^{}s^{}\mathrm{\hspace{0.17em}0}}^{1/\mathrm{2\hspace{0.17em}1}/2j_3}(2j_3+1)^2`$ (83)
$`\times {\displaystyle \underset{m_1m_2,1/2}{\overset{1/2}{}}}()^{m_2}{\displaystyle \underset{m_3=j_3}{\overset{j_3}{}}}D_{0m_3}^{(j_3)}(\alpha ^{},\beta ^{},\gamma ^{})W_{m_1m_2m_3}^{1/\mathrm{2\hspace{0.17em}1}/2j_3}`$ (84)
$`\times ()^{1/2m_2}(ac)D_{sm_1}^{(1/2)}(\alpha ,\beta ,\gamma )D_{sm_2}^{(1/2)}(\alpha ,\beta ,\gamma ).`$ (85)
Due to the properties of the Wigner-$`3j`$ symbols we may see that the terms with $`j_3=0,1`$, and $`m_3=0`$ do not give contributions; moreover, changing the value of $`s^{}`$, it changes only the sign. Thus, we will get
$`_tw(\frac{1}{2},\alpha ,\beta ,t)`$ $`=`$ $`{\displaystyle \frac{d\mathrm{\Omega }^{}}{8\pi ^2}\left[w(\frac{1}{2},\alpha ^{},\beta ^{},t)w(\frac{1}{2},\alpha ^{},\beta ^{},t)\right]}`$ (86)
$`\times `$ $`{\displaystyle \frac{3}{2}}(ac)\mathrm{sin}\beta ^{}\mathrm{sin}\beta \mathrm{sin}(\alpha \alpha ^{})`$ (87)
and by using the normalization condition it can be rewritten as
$$_tw(s,\alpha ,\beta ,t)=3(ac)\mathrm{sin}\beta \frac{d\mathrm{\Omega }^{}}{8\pi ^2}w(s,\alpha ^{},\beta ^{},t)\mathrm{sin}\beta ^{}\mathrm{sin}(\alpha \alpha ^{})$$
(88)
which is similar to that derived in Ref. (the differencies are due to the degeneracy of the spin-$`\frac{1}{2}`$ systems). It should be noted in the above equation that the argument $`s`$ is the same in both sides; this is consistent with the fact that $`\widehat{H}`$ in Eq.(80) does not mix states with different $`s`$. On the other hand it can be easily checked that the sum over $`s`$ at r.h.s. of Eq.(88) causes the integral to become zero; this is consistent with the fact that at l.h.s. we will obtain the time derivative of a constant. Also, if $`a=c`$, the r.h.s. of Eq.(88) will be zero since the Hamiltonian (80) will be proportional to the identity and will not produce any evolution.
## VII Examples
In the previous section, we discussed the probability for the joint measurement of the spin and spatial variables. Therefore, here we would like to consider some examples involving both variables.
At first we consider a system with the following hamiltonian
$$\widehat{H}=\frac{1}{2}\left(\widehat{p}^2+\widehat{q}^2\right)+\left(|\frac{1}{2}\frac{1}{2}||\frac{1}{2}\frac{1}{2}|\right).$$
(89)
It could describe e.g. one vibrational degree of a trapped electron plus its spin . The measurability of marginals in this system is investigated in Ref. . Here, as a straigthforward extension of the arguments of Sec. VI.1 we obtain
$`_tw(x,\mu ,\nu ,s,\alpha ,\beta )`$ $`=`$ $`\left(\mu _\nu \nu _\mu \right)w(x,\mu ,\nu ,s,\alpha ,\beta )`$ (90)
$`+`$ $`6\mathrm{sin}\beta {\displaystyle \frac{d\mathrm{\Omega }^{}}{8\pi ^2}w(x,\mu ,\nu ,s,\alpha ^{},\beta ^{},t)\mathrm{sin}\beta ^{}\mathrm{sin}(\alpha \alpha ^{})}.`$ (91)
Let us now consider an initial entangled state like
$$\mathrm{\Psi }(0)=\frac{1}{\sqrt{2}}\left(0\frac{1}{2}+1\frac{1}{2}\right),$$
(92)
where $`n`$ represents the number eigenstate of a harmonic oscillator. At Eq.(92) corresponds the following marginal
$$w(x,\mu ,\nu ,s,\alpha ,\beta ,t=0)=\frac{1}{2}\left[w_{00}+w_{11}+w_{01}+w_{10}\right],$$
(93)
where
$`w_{00}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\pi (\mu ^2+\nu ^2)}}}\mathrm{exp}\left[{\displaystyle \frac{x^2}{\mu ^2+\nu ^2}}\right]D_{s{\scriptscriptstyle \frac{1}{2}}}^{(1/2)}(\alpha ,\beta ,\gamma )D_{s{\scriptscriptstyle \frac{1}{2}}}^{(1/2)}(\alpha ,\beta ,\gamma ),`$ (94)
$`w_{11}`$ $`=`$ $`{\displaystyle \frac{2x^2}{\sqrt{\pi (\mu ^2+\nu ^2)^3}}}\mathrm{exp}\left[{\displaystyle \frac{x^2}{\mu ^2+\nu ^2}}\right]D_{s{\scriptscriptstyle \frac{1}{2}}}^{(1/2)}(\alpha ,\beta ,\gamma )D_{s{\scriptscriptstyle \frac{1}{2}}}^{(1/2)}(\alpha ,\beta ,\gamma ),`$ (95)
$`w_{01}`$ $`=`$ $`i{\displaystyle \frac{\sqrt{2}x(\nu i\mu )}{\sqrt{\pi (\mu ^2+\nu ^2)^3}}}\mathrm{exp}\left[{\displaystyle \frac{x^2}{\mu ^2+\nu ^2}}\right]D_{s{\scriptscriptstyle \frac{1}{2}}}^{(1/2)}(\alpha ,\beta ,\gamma )D_{s{\scriptscriptstyle \frac{1}{2}}}^{(1/2)}(\alpha ,\beta ,\gamma ),`$ (96)
$`w_{10}`$ $`=`$ $`w_{01}^{}.`$ (97)
Then, the solution of the Pauli equation (90) is
$$w(x,\mu ,\nu ,s,\alpha ,\beta ,t)=\frac{1}{2}\left[w_{00}+w_{11}+w_{01}e^{3it}+w_{10}e^{3it}\right].$$
(98)
As a second example, we want to consider the case of Landau levels , i.e. a charged particle moving in a classical magnetic field $`\stackrel{}{B}`$ being time-independent and axial symmetric. The particle’s movement along the axis being free, instead the Hamiltonian of the transverse motion reads
$$\widehat{H}=\frac{1}{2}\left[\left(\widehat{p}_1\widehat{A}_1\right)^2+\left(\widehat{p}_2\widehat{A}_2\right)^2\right],\widehat{\stackrel{}{A}}=\left[\stackrel{}{B}\times \frac{\widehat{\stackrel{}{r}}}{2}\right],$$
(99)
where $`\widehat{\stackrel{}{r}}=(\widehat{q}_1,\widehat{q}_2)`$ is the radius-vector of the particle’s center, $`\widehat{p}_1`$ and $`\widehat{p}_2`$ are the particle’s momentum components in the transverse plane. Having $`\stackrel{}{B}`$ along the third axis and choosing $`\stackrel{}{B}=2`$, we get
$$\widehat{H}=\frac{1}{2}\left(\widehat{p}_1^2+\widehat{p}_2^2+\widehat{q}_1^2+\widehat{q}_2^2\right)+\left(\widehat{p}_1\widehat{q}_2\widehat{p}_2\widehat{q}_1\right)+\left(|\frac{1}{2}\frac{1}{2}||\frac{1}{2}\frac{1}{2}|\right)$$
(100)
In this case the kernel $`\mathrm{\Theta }`$ of Eq.(69) is given by
$`\mathrm{\Theta }`$ $`=`$ $`{\displaystyle \frac{i}{4\pi ^2}}{\displaystyle \frac{dq_1}{\nu _1}\frac{dq_2}{\nu _2}e^{i_{l=1}^2[(\mu _l^{}\nu _l^{}/2x_l^{})+\mu _l(q_l\nu _l^{})+(x_l\nu _l^{}+\mu _l\nu _{l}^{}{}_{}{}^{2}\mu _l\nu _l^{}q_l)/\nu _l]}}`$ (101)
$`\times `$ $`\{{\displaystyle \underset{l=1}{\overset{2}{}}}(\mu _l^{}\mu _lq_l/\nu _l\mu _l^{}x_l/\nu _l+\nu _l^{}q_l\mu _{l}^{}{}_{}{}^{2}/2\nu _{l}^{}{}_{}{}^{2}/2)(x_2\mu _2q_2)\nu _1^{}/\nu _2`$ (102)
$`+`$ $`(\mu _2^{}q_1\mu _2^{}\nu _1^{})+(x_1\mu _1q_1)\nu _2^{}/\nu _1(\mu _1^{}q_2\mu _1^{}\nu _2^{})\}\delta _{s,s^{}}\delta (\mathrm{\Omega }\mathrm{\Omega }^{})`$ (103)
$`+`$ $`{\displaystyle \frac{6}{8\pi ^2}}\mathrm{sin}\beta \mathrm{sin}\beta ^{}\mathrm{sin}(\alpha \alpha ^{})\delta (\stackrel{}{x}\stackrel{}{x}^{})\delta (\stackrel{}{\mu }\stackrel{}{\mu }^{})\delta (\stackrel{}{\nu }\stackrel{}{\nu }^{}).`$ (104)
As nontrivial example we also consider here an initial state which is the entangled superposition
$$\mathrm{\Psi }(0)=\frac{1}{\sqrt{2}}\left[\mathrm{0\hspace{0.17em}0}\frac{1}{2}+\mathrm{1\hspace{0.17em}0}\frac{1}{2}\right].$$
(105)
It leads to nonfactorisable marginal
$$w\left(\stackrel{}{x},\stackrel{}{\mu },\stackrel{}{\nu },s,\alpha ,\beta ,t=0\right)=\frac{1}{2}\left[w_{0000}+w_{1010}+w_{0010}+w_{1000}\right],$$
(106)
where
$$w_{n_1n_2n_1^{}n_2^{},m_1m_2}=w_{n_1n_2n_1^{}n_2^{}}\times D_{sm_1}^{(1/2)}(\alpha ,\beta ,\gamma )D_{sm_2}^{(1/2)}(\alpha ,\beta ,\gamma ).$$
(107)
Here, $`m_1=\frac{1}{2}`$ ($`m_1=\frac{1}{2}`$) replaces downarrow (uparrow) while the spatial part $`w_{n_1n_2n_1^{}n_2^{}}`$ is explicitely calculated in the Appendix.
It is now easy to see that the solution of the Eq.(67) with the kernel (101), subject to the above initial condition, is
$$w(\stackrel{}{x},\stackrel{}{\mu },\stackrel{}{\nu },s,\alpha ,\beta ,t)=\frac{1}{2}\left[w_{0000}+w_{1010}+w_{0010}e^{3it}+w_{1000}e^{3it}\right].$$
(108)
## VIII Conclusion
We conclude that it is possible to obtain an evolution equation for the tomographic probabilities (marginal distributions) of an arbitrary tomography scheme. The main result of our paper is the analog of the Pauli equation for spin-$`\frac{1}{2}`$ particle.
The explicit expression for the marginal distribution for a trapped particle as well as for Landau levels has been studied. It results that in the nonstationary case they obey the analog of the Pauli equation.
The examples considered demonstrate that the usual problems of conventional quantum mechanics can be cast into the form in which only positive probabilities are used to describe quantum states and their evolution. A possible disadvantage of the approach proposed is a complicated evolution equation (67), but, perhaps, this is the price one ought to pay for the possibility of describing quantum objects in terms of classical probabilities.
Anyway, our argumentations can constitute a step further from the Bohr position about the inapplicability of classical modes of description in the quantum domain. In fact, while we belive that quantum mechanics is not classical physics in disguise, we retain (some) classical concepts still applicable against counterintuitive notions like complex statefunctions.
We also belive that the developed classical-like formalism could be applied to describe quantum mechanical paradoxes, because usually, if there is a paradox in quantum mechanics, there should also be a classical one, perhaps worse . These aspects will be investigated in a forthcoming paper as well as the extension of the presented approach to the relativistic domain , in order to find an analog of the Dirac equation.
## Acknowledgments
S.M. would like to thank the Lebedev Physical Institute for the kind hospitality during the first stage of this work. O.V.M. thanks the Department of Physics and Mathematics of the University of Camerino for partial support. O.V.M. and V.I.M. are grateful to the Russian Foundation for Basic Research for the partial support under Project No. 99-02-17753. S.M. and P.T. are grateful to the M.U.R.S.T. for the partial support under the Project “Cofinanaziamento”.
## Appendix
The wave function of the particle’s coherent state in a magnetic field $`\stackrel{}{B}`$ is
$$\mathrm{\Psi }_{\alpha ,\beta }(q_1,q_2)=\frac{1}{\sqrt{\pi }}\mathrm{exp}\left\{\frac{q_1^2+q_2^2}{2}\frac{|\alpha |^2}{2}\frac{|\beta |^2}{2}i\alpha \beta +\beta \left(q_1+iq_2\right)+i\alpha \left(q_1iq_2\right)\right\},$$
(109)
where $`q_1`$ and $`q_2`$ are the particle’s coordinates and $`\alpha `$ and $`\beta `$ complex numbers.
The coherent state (109) is the superposition of number states
$$\mathrm{\Psi }_{\alpha ,\beta }(q_1,q_2)=\mathrm{exp}\left(\frac{\alpha ^2}{2}\frac{\beta ^2}{2}\right)\underset{n=0}{\overset{\mathrm{}}{}}\underset{n^{}=0}{\overset{\mathrm{}}{}}\frac{\alpha ^n\beta ^n^{}\mathrm{\Psi }_{nn^{}}(q_1,q_2)}{\sqrt{n!n^{}!}}.$$
(110)
In view of the general relationship between the marginal distribution and wave function , we have
$`w(x_1,x_2,\mu _1,\nu _1,\mu _2,\nu _2)={\displaystyle \frac{1}{4\pi ^2\nu _1\nu _2}}`$ (111)
$`\times \left|{\displaystyle \mathrm{exp}\left(\frac{iy_1^2\mu _1}{2\nu _1}\frac{iy_1x_1}{\nu _1}+\frac{iy_2^2\mu _2}{2\nu _2}\frac{iy_2x_2}{\nu _2}\right)\mathrm{\Psi }_{\alpha \beta }(y_1,y_2)𝑑y_1𝑑y_2}\right|^2,`$ (112)
where parameters $`\mu _1,\nu _1,\mu _2,\nu _2`$, as usually, mark reference frames; then, one obtains for the marginal distribution of the particle’s coherent state without spin in a magnetic field the following expression
$`w_{\alpha \beta }(x_1,x_2,\mu _1,\nu _1,\mu _2,\nu _2)={\displaystyle \frac{\mathrm{exp}\left[|\alpha |^2|\beta |^2i\left(\alpha \beta \alpha ^{}\beta ^{}\right)\right]}{\pi \sqrt{\left(\nu _1^2+\mu _1^2\right)\left(\nu _2^2+\mu _2^2\right)}}}`$ (113)
$`\times \mathrm{exp}\{{\displaystyle \frac{\left(\nu _1+i\mu _1\right)\left(i\alpha \nu _1+\beta \nu _1ix_1\right)^2+\left(\nu _1i\mu _1\right)\left(i\alpha ^{}\nu _1\beta ^{}\nu _1+ix_1\right)^2}{2\nu _1\left(\mu _1^2+\nu _1^2\right)}}`$ (114)
$`+{\displaystyle \frac{\left(\nu _2+i\mu _2\right)\left(\alpha \nu _2+i\beta \nu _2ix_2\right)^2+\left(\nu _2i\mu _2\right)\left(\alpha ^{}\nu _2i\beta ^{}\nu _2+ix_2\right)^2}{2\nu _2\left(\mu _2^2+\nu _2^2\right)}}\}.`$ (115)
Multiplying (113) by $`\mathrm{exp}\left(|\alpha |^2+|\beta |^2\right)`$ and expanding the expression obtained into the power series, we arrive at
$$w_{\alpha \beta }(x_1,x_2,\mu _1,\nu _1,\mu _2,\nu _2)e^{|\alpha |^2}e^{|\beta |^2}=\underset{n_1,n_2,n_1^{},n_2^{}=0}{\overset{\mathrm{}}{}}\frac{\alpha ^{n_1}(\alpha ^{})^{n_2}\beta ^{n_1^{}}(\beta ^{})^{n_2^{}}w_{n_1n_2n_1^{}n_2^{}}}{\sqrt{n_1!n_2!n_1^{}!n_2^{}!}}.$$
(116)
Taking into account the property of the generating function for multivariate Hermite polynomials , namely,
$$\mathrm{exp}\left\{\frac{1}{2}\stackrel{}{u}M\stackrel{}{u}+\stackrel{}{u}M\stackrel{}{\zeta }\right\}=\underset{n_1,n_2,n_1^{},n_2^{}=0}{\overset{\mathrm{}}{}}\frac{\alpha ^{n_1}(\alpha ^{})^{n_2}\beta ^{n_1^{}}(\beta ^{})^{n_2^{}}}{n_1!n_2!n_1^{},n_2^{}!}H_{n_1n_2n_1^{},n_2^{}}^{\left\{M\right\}}\left(\stackrel{}{\zeta }\right),$$
(117)
where the vector $`\stackrel{}{u}`$ has components $`\stackrel{}{u}=(\alpha ,\alpha ^{},\beta ,\beta ^{})`$, and comparing (116) with (117), we obtain
$`w_{n_1n_2n_1^{}n_2^{}}(x_1,x_2,\mu _1,\nu _1,\mu _2,\nu _2)={\displaystyle \frac{1}{\pi \sqrt{\left(\nu _1^2+\mu _1^2\right)\left(\nu _2^2+\mu _2^2\right)}}}`$ (118)
$`\times \mathrm{exp}\left({\displaystyle \frac{x_1^2}{\mu _1^2+\nu _1^2}}{\displaystyle \frac{x_2^2}{\mu _2^2+\nu _2^2}}\right){\displaystyle \frac{H_{n_1n_2n_1^{},n_2^{}}^{\left\{M\right\}}\left(\stackrel{}{\zeta }\right)}{\sqrt{n_1!n_2!n_1^{}!n_2^{}!}}},`$ (119)
where the $`4\times 4`$ matrix $`M`$ reads
$$M=\left(\begin{array}{cc}M^{(1)}& \hfill M^{(2)}\\ M^{(4)}& \hfill M^{(3)}\end{array}\right).$$
(120)
The $`2\times 2`$ matrices $`M^{(r)}`$ are given by
$`M_{k,l}^{(r)}`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{2}{}}}{\displaystyle \frac{\nu _j}{\nu _j+i()^l\mu _j}}()^{j+(r+1)/2}\delta _{k+l,\mathrm{even}},r=1,3,`$ (121)
$`M_{k,l}^{(r)}`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{2}{}}}i\left[{\displaystyle \frac{\nu _j}{\nu _j+i()^l\mu _j}}()^{l+(j1)(r/21)}+()^{l1}\right]\delta _{k+l,\mathrm{even}},r=2,4,`$ (123)
The argument of the multivariate Hermite polynomials $`\stackrel{}{\zeta }=(\zeta _1,\zeta _1^{},\zeta _2,\zeta _2^{})`$ is expressed in terms of the parameters as follows
$`\zeta _1`$ $`=`$ $`{\displaystyle \frac{ix_1}{\sqrt{\mu _1^2+\nu _1^2}}}\mathrm{exp}\left(i\mathrm{tan}^1{\displaystyle \frac{\mu _2}{\nu _2}}\right){\displaystyle \frac{x_2}{\sqrt{\mu _2^2+\nu _2^2}}}\mathrm{exp}\left(i\mathrm{tan}^1{\displaystyle \frac{\mu _1}{\nu _1}}\right),`$ (124)
$`\zeta _2`$ $`=`$ $`{\displaystyle \frac{ix_2}{\sqrt{\mu _1^2+\nu _1^2}}}\mathrm{exp}\left(i\mathrm{tan}^1{\displaystyle \frac{\mu _1}{\nu _1}}\right){\displaystyle \frac{x_1}{\sqrt{\mu _2^2+\nu _2^2}}}\mathrm{exp}\left(i\mathrm{tan}^1{\displaystyle \frac{\mu _2}{\nu _2}}\right).`$ (126)
Taking $`n_1=n_2`$ and $`n_1^{}=n_2^{}`$ we obtain the marginal distribution $`w_{nn^{}}(x_1`$,$`x_2`$,$`\mu _1`$,$`\nu _1`$,$`\mu _2`$,$`\nu _2)`$ for the Landau level states $`nn^{}`$
$`w_{nn^{}}(x_1,x_2,\mu _1,\nu _1,\mu _2,\nu _2)w_{nnn^{}n^{}}(x_1,x_2,\mu _1,\nu _1,\mu _2,\nu _2)`$ (128)
$`={\displaystyle \frac{1}{\pi \sqrt{\left(\nu _1^2+\mu _1^2\right)\left(\nu _2^2+\mu _2^2\right)}n!n^{}!}}\mathrm{exp}\left({\displaystyle \frac{x_1^2}{\mu _1^2+\nu _1^2}}{\displaystyle \frac{x_2^2}{\mu _2^2+\nu _2^2}}\right)H_{nnn^{}n^{}}^{\left\{M\right\}}\left(\stackrel{}{\zeta }\right),`$
where $`n`$ is the main quantum number and $`n^{}n=l`$ is the angular momentum quantum number.
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# Exact Equations and Scaling Relations for 𝑓̄₀-avalanche in the Bak-Sneppen Evolution Model
## Abstract
Infinite hierarchy of exact equations are derived for the newly observed $`\overline{f}_0`$-avalanche in the Bak-Sneppen model. By solving the first order exact equation, we find that the critical exponent $`\gamma `$, governing the divergence of the average avalanche size, is exactly $`1`$ (for all dimensions), which has been confirmed by extensive simulations. Solution of the gap equation yields another universal result $`\rho =1`$ ($`\rho `$ is the exponent of relaxation to attractor). Scaling relations are established among the critical exponents ($`\gamma `$, $`\tau `$, $`D`$, $`\sigma `$ and $`\nu `$) for $`\overline{f}_0`$-avalanche.
PACS number(s): 87.10.+e, 05.40.+j, 64.60.Lx
In the Bak-Sneppen (BS) evolution model , random numbers, $`f_i`$, chosen from a flat distribution between $`0`$ and $`1`$, $`p(f)`$, are assigned independently to each species located on a $`d`$-dimensional lattice of linear size $`L`$. At each time step, the extremal site, i.e., the species with the smallest random number, and its $`2d`$ nearest neighboring sites, are assigned $`2d+1`$ new random numbers also chosen from $`p(f)`$. This updating continues indefinitely. After a long transient process the system reaches a statistically stationary state where the density of random numbers in the system is uniform above $`f_c`$ (the self-organized threshold) and vanishes for $`f<f_c`$ .
Despite the fact that it is an oversimplification of real biological process, BS model exhibits such common interesting features observed by paleontologists as punctuated equilibria, power-law probability distributions of lifetimes of species and of the sizes of extinction events. These behaviors suggest that the ecology of interacting species might have evolved to a self-organized critical state.
BS model displays spatial-temporal complexity, which also emerges from many natural phenomena, such as fractals , $`1/f`$ noise , etc. This strongly suggests that various complex behaviors may be attributed to a common underlying mechanism. Authors in Ref. suggest that the relation of these different phenomena can be established on the basis of their unique models. It is even proposed by them that spatial-temporal complexity comes out as the direct results of avalanche dynamics in driven systems, and different complex phenomena are related via scaling relations to the fractal properties of the avalanches. It hence can be inferred that avalanche dynamics plays a key role in dealing with complex systems, especially when one needs to know the macroscopic features of the systems, since lingering on the inner structure of individuals will not be helpful .
Avalanche is a kind of macroscopic phenomenon driven by local interactions. The size of an avalanche may be extremely sensitive to the initial configuration of the system, while the distribution of the sizes (spatial and temporal) of avalanches, i.e. the ”fingerprint” should be robust with respect to the modifications, due to the universality of complexity and the definition of self-organized criticality (SOC) . In this sense, the extent that we know about avalanche will determine to what extent we do a complex system. Avalanche dynamics provides insight into complexity and enables one to further investigate the system studied.
Though avalanche dynamics may be a possible underlying mechanism of complexity, the definitions of avalanches can be vastly different for various complex systems, or for same sorts of systems, even for the same one. In BTW model , an avalanche is intrigued by the adding of a grain or several grains of sand into the system. The avalanche is considered over when the heights of all the sites are less than the critical value, say, $`4`$. In BS model , several types of avalanches, for instance, $`f_0`$-avalanche, $`G(s)`$-avalanche, forward avalanche and backward avalanche, etc, are presented. These different definitions of avalanches may show their unique hierarchal structures, while they manifest the common fractal feature of the complex system, that is, SOC. It can be inferred that various types of avalanches are equivalent in the sense that they imply complexity.
Since similar structures and common features evidently arise in different types of avalanches, it is straightforward that various avalanches differ each other only in the contexts from which one comprehends them. As known, the major aim of avalanche study is to investigate the universal rules possibly hidden behind the evolution of the systems or the models. Hence, the means of understanding the avalanches appear crucial. Better ways may enable one to know more about the system or the model and hence to have better comprehension of the features corresponding to complexity. From this point of view, when studying avalanches one should try to choose the easier ways instead of the more difficult ones.
The evolution of the highly sophiscated BS model shows a hierarchal structure specified by avalanches, which correspond to sequential mutations below certain threshold. It has been noted that in BS model an avalanche is initiated when the fitness of the globally extremal site (the species with the least random number) is larger than the self-organized threshold. That is, the triggering event of an avalanche is directly related to the fitness, the feature of individuals. In other words, the avalanche is directly associated with the feature of individuals instead of general features of the ecosystem as a whole. Is it feasible that the avalanches are directly intrigued by the global feature of the whole system? Can such global feature be expressed in terms of the corresponding quantity? If such quantity found and such avalanches observed, may the new avalanches provide a new and easier way in investigating properties of the model?
One of our previous works presents such a different hierarchy of avalanches ($`\overline{f}_0`$-avalanche) for BS model. We defined a global quantity, $`\overline{f}`$, which denotes the average fitness of the system. The new type of avalanches are directly related to $`\overline{f}`$. In this paper, we present a master equation for the hierarchal structure of $`\overline{f}_0`$-avalanches. It prescribes the cascade process of smaller avalanches merging into bigger avalanches when the critical parameter $`\overline{f}_0`$ is changed. An infinite series of exact equations can be derived from this master equation . The first order exact equation, together with an scaling ansatz of the average sizes of avalanches, shows the exact result of $`\gamma `$, the critical exponent governing self-organization, to be universally $`1`$ for all dimensional BS models, which has been confirmed by extensive simulations of the model. We also establish scaling relations related to some critical exponents for $`\overline{f}_0`$-avalanche and make predictions on the values of some exponents.
The quantity $`\overline{f}`$ is a global one of the ecosystem and can be expected to involve some general information about the whole system. It may represent the average population or living capability of the whole species system. Larger $`\overline{f}`$ shows that the average population is immense or the average living capability is great, and vice versa. $`\overline{f}`$ is defined as
$$\overline{f}=\frac{1}{L^d}\underset{i=1}{\overset{L^d}{}}f_i,$$
(1)
where $`f_i`$ is the fitness of the $`i`$th species of a system consists of $`L^d`$ species. Let BS model start to evolve. At each time step of the evolution, apart from the random numbers of the globally extremal site and its $`2d`$ nearest neighboring sites, the signal $`\overline{f}`$ is also tracked. Initially, $`\overline{f}`$ tends to increase step-wisely. As the evolution continues further, $`\overline{f}`$ approaches a critical value $`\overline{f}_c`$ and remain statistically stable around $`\overline{f}_c`$. The plot of $`\overline{f}`$ versus time step $`s`$ shows that the increasing signals of $`\overline{f}`$ follow a Devil’s staircase , which implies that punctuated equilibrium emerges. Denote $`F(s)`$ the gap of the punctuated equilibrium. Actually, $`F(s)`$ tracks the peaks in $`\overline{f}`$. After some careful derivation one can write down an exact gap equation
$$\frac{dF(s)}{ds}=\frac{\overline{f}_cF(s)}{L^dS_{F(s)}},$$
(2)
where $`S_{F(s)}`$ denotes the average size of avalanches occurred during the gap $`F(s)`$ when $`\overline{f}<F(s)`$. This exact gap equation will be exactly solved in this paper.
Signals $`\overline{f}(s)`$ play important roles in defining $`\overline{f}_0`$-avalanche. For any value of the auxiliary parameter $`\overline{f}_0`$ ($`0.5<\overline{f}_0<1.0`$), an $`\overline{f}_0`$-avalanche of size $`S`$ is defined as a sequence of $`S1`$ successive events when $`\overline{f}(s)<\overline{f}_0`$ confined between two events when $`\overline{f}(s)>\overline{f}_0`$. This definition ensures that the mutation events during an avalanche are spatially and temporally correlated. It can also guarantee the hierarchal structure of the avalanches: larger avalanches consists of smaller ones. As $`\overline{f}_0`$ is raised, smaller avalanches gather together and form bigger ones. The statistics of $`\overline{f}_0`$ will inevitably have a cutoff if $`\overline{f}_0`$ is not chosen to be $`\overline{f}_c`$. This will not affect the size distribution provided that $`\overline{f}_0`$ approaches $`\overline{f}_c`$. Extensive simulations show that exponents $`\tau `$ of $`\overline{f}_0`$-avalanche size distribution are 1.800 and 1.725 for 1D and 2D BS models respectively, amazingly different from the counterparts of $`f_0`$-avalanche, 1.07 and 1.245 . This strengthens the speculation that $`\overline{f}_0`$-avalanche is a different type of avalanche, distinguished from any types of avalanches found previously.
Denote $`P(S,\overline{f}_0)`$ the probability of acquiring a $`\overline{f}_0`$-avalanche of size $`S`$. The signals $`\overline{f}(s)`$ ($`\overline{f}_0<\overline{f}(s)<\overline{f}_0+d\overline{f}_0`$) will stop the $`\overline{f}_0`$-avalanches and not ($`\overline{f}_0+d\overline{f}_0`$)-avalanches. That is , as $`\overline{f}_0`$ is raised by an infinitesimal amount $`d\overline{f}_0`$ some of $`\overline{f}_0`$-avalanches merge together to form bigger ($`\overline{f}_0+d\overline{f}_0`$)-avalanches. This exhibits a hierarchal structure of $`\overline{f}_0`$-avalanches and will be prescribed by the below exact master equation. In some sense, the master equation reflects the ”flow” of probability of avalanche size distribution with respect to the change in $`\overline{f}_0`$.
Simulations show that $`\overline{f}`$ approaches $`\overline{f}_c`$ and remain statistically stable in the critical state. This feature is greatly different from the feature of $`f_{\mathrm{min}}`$ (fitness of globally extremal site), which can vary between 0 and 1. While $`\overline{f}`$ in the critical state fluctuates slightly around $`\overline{f}_c`$. Therefore, the $`\overline{f}_0`$-avalanches will have no good statistics if $`\overline{f}_0`$ is chosen as the value far less than $`\overline{f}_c`$, since there only exists smaller avalanches in the model. To acquire a better and reasonable distribution of $`\overline{f}_0`$-avalanches sizes, one should choose the value of $`\overline{f}_0`$ under the condition $`\overline{f}_0\overline{f}_c`$. It should be emphasized that the master equation listed below is valid also for $`\overline{f}_0\overline{f}_c`$.
Both theoretical analysis and extensive simulations suggest that the signals $`\overline{f}(s)`$ which terminate $`\overline{f}_0`$-avalanches are uncorrelated and evenly distributed between $`(\overline{f}_0,\overline{f}_c)`$ provided that $`\overline{f}_0\overline{f}_c`$. The direct consequence of this observation is that the probability of an $`\overline{f}_0`$-avalanche merging to $`\overline{f}_0+d\overline{f}_0`$-avalanche is prescribed by $`\frac{d\overline{f}_0}{\overline{f}_c\overline{f}_0}`$. It is important to note that any two subsequent avalanches are mutually independent for the following arguments to be true. In other words, the probability distribution of $`\overline{f}_0`$-avalanches, initiated immediately after the termination of an $`\overline{f}_0`$-avalanche of size $`S`$ is independent of $`S`$. This is true because in BS model the dynamics within an $`\overline{f}_0`$-avalanche is completely independent of the particular value of the signals $`\overline{f}(s)>\overline{f}_0`$ in the background that were left by the previous avalanches.
Here present the master equation. As $`\overline{f}_0`$ is raised by an infinitesimal amount $`d\overline{f}_0`$, the probability ”flowing” out of the size distribution of $`\overline{f}_0`$-avalanches is given by $`P(S,\overline{f}_0)\frac{d\overline{f}_0}{\overline{f}_c\overline{f}_0}`$, while the probability ”flowing” into is given by $`_{S_1=1}^{S1}\frac{P(S_1,\overline{f}_0)}{\overline{f}_c\overline{f}_0}P(SS_1,\overline{f}_0)`$. Let $`\overline{f}_0\overline{f}_c`$ and $`d\overline{f}_00`$, one can write down the master equation as
$`(\overline{f}_c\overline{f}_0){\displaystyle \frac{P(S,\overline{f}_0)}{\overline{f}_0}}`$ $`=`$ $`P(S,\overline{f}_0)`$ (4)
$`+{\displaystyle \underset{S_1=1}{\overset{S1}{}}}P(S_1,\overline{f}_0)P(SS_1,\overline{f}_0).`$
The first term on the right hand of the equation expresses the loss of avalanches of size $`S`$ due to the merging with the subsequent one, while the second one describes the gain in $`P(S,\overline{f}_0)`$ due to merging of avalanches of size $`S_1`$ with avalanches of size $`SS_1`$.
In order to investigate the exact master equation it is convenient to make some variable changes. Define $`h=\mathrm{ln}(\overline{f}_c\overline{f}_0)`$. Therefore, $`\overline{f}_0=\overline{f}_c`$ corresponds to $`h=+\mathrm{}`$. Since in the master equation $`\overline{f}_0`$ is chosen to be close to $`\overline{f}_c`$, $`h`$ varies from a very large number to $`+\mathrm{}`$. Due to the variable change the variable $`h`$ is chosen from the distribution $`P(h)=e^h`$, which seems to be more ”natural”. In the following part we will use the new variable $`h`$ instead of $`\overline{f}_0`$. The master equation can be rewritten, in terms of $`h`$, as
$$\frac{P(S,h)}{h}=P(S,h)+\underset{S_1=1}{\overset{S1}{}}P(S,h)P(SS_1,h).$$
(5)
Making Laplace transformation of Eq.(4), after some calculation, one obtains
$$\frac{\mathrm{ln}(1p(\beta ,h))}{h}=p(\beta ,h),$$
(6)
where $`p(\beta ,h)=\underset{S=1}{\overset{\mathrm{}}{}}P(S,h)e^{\beta S}`$. This exact equation is the key one of this work. Many interesting physical features can be derived from it. As $`h<+\mathrm{}`$ avalanches size will have a cutoff. The normalization of $`P(S,h)`$ can be expressed as $`p(0,h)=\underset{S=1}{\overset{\mathrm{}}{}}P(S,h)=1`$. Expanding both sides of Eq. (5) as Tylor series throughout a neighborhood of the point $`\beta =0`$, one can immediately obtain
$`{\displaystyle \frac{}{h}}[1S_h\beta +{\displaystyle \frac{1}{2!}}S^2_h\beta ^2{\displaystyle \frac{1}{3!}}S^3_h\beta ^3+\mathrm{}]=`$ (7)
$`[S_h\beta {\displaystyle \frac{1}{2!}}S^2_h\beta ^2+{\displaystyle \frac{1}{3!}}S^3_h\beta ^3+\mathrm{}]\times `$ (8)
$`[1+S_h\beta {\displaystyle \frac{1}{2!}}S^2_h\beta ^2+{\displaystyle \frac{1}{3!}}S^3_h\beta ^3+\mathrm{}]`$ . (9)
Since the equation (6) holds for arbitrary $`\beta `$, comparing the coefficients of different powers of $`\beta `$ in the above Taylor series gives an infinite series of exact equations. Comparison of the coefficients of $`\beta ^1`$ results in
$$\frac{\mathrm{ln}S_h}{h}=1.$$
(10)
Eq. (7) is extremely interesting. Changing variable $`h`$ back into $`\overline{f}_0`$, one can obtain the ”gamma” equation
$$\frac{d\mathrm{ln}S_{\overline{f}_0}}{d\overline{f}_0}=\frac{1}{\overline{f}_c\overline{f}_0}.$$
(11)
Inserting the scaling ansatz $`S_{\overline{f}_0}(\overline{f}_c\overline{f}_0)^\gamma `$ into Eq. (8), one immediately obtain an interesting result
$$\gamma =1$$
(12)
.
It should be noted that $`\gamma =1`$ is universal, that is, independent of the dimension. The value of $`\gamma `$ for $`\overline{f}_0`$-avalanches is different from those for $`f_0`$-avalanche found in Ref. , which are 2.70 and 1.70 for 1D and 2D BS models respectively. Extensive simulations show $`\gamma =0.99\pm 0.01`$ and $`\gamma =0.98\pm 0.01`$ for 1D and 2D BS models respectively. Fig. (1) shows our simulation results, which confirms the universal result $`\gamma =1`$.
Higher powers of $`\beta `$ gives new exact equations. Here present the first two
$$\frac{}{h}(\frac{S^2_h}{S_h})=2S;$$
(13)
$$\frac{}{h}(\frac{S^3_h}{3S_h}\frac{S^2_h^2}{2S_h^2})=S^2_h.$$
(14)
Next present the solution of the exact gap equation for $`\overline{f}_0`$-avalanches. Inserting the scaling relation $`S_{F(s)}(\overline{f}_cF(s))^1`$ into the equation and integrating, one obtains
$$\mathrm{\Delta }\overline{f}(s)=\overline{f}_cF(s)(\frac{s}{L^d})^\rho =(\frac{s}{L^d})^1,$$
(15)
where $`\rho `$ is the exponent of relaxation to attractor . Thus, we obtain $`\rho =1`$. Interestingly, $`\rho `$ is also a universal exponent for all dimensional BS models. It shows that the critical point ($`\mathrm{\Delta }\overline{f}=0`$) is approached algebraically with an exponent $`1`$.
Up to now, we have obtained some exponents of corresponding physical properties of $`\overline{f}_0`$-avalanches: $`\tau `$, avalanche size distribution , $`D`$, avalanche dimension , $`\gamma `$, average avalanche size , and $`\rho `$, relaxation to attractor . Recall another two exponents : $`\nu `$, $`\sigma `$, which are defined as $`r_{co}(\overline{f}_c\overline{f}_0)^\nu `$ and $`S_{co}=(\overline{f}_c\overline{f}_0)^{\frac{1}{\sigma }}`$ respectively. Here $`r_{co}`$ and $`S_{co}`$ are referred to as the cut-off of the spatial extent of avalanches (due to the limit system size ) and that of the avalanche size (due to the fact that $`\overline{f}_0`$ is not chosen as $`\overline{f}_c`$) respectively. It is natural to establish some scaling relations of these exponents for $`\overline{f}_0`$-avalanches similar to those found in Ref. for $`f_0`$-avalanches. Nevertheless, these two types of avalanches manifest similar fractal properties. Hence some common features should be shared by them. Integrating of the equation $`S=_1^{(\overline{f}_c\overline{f}_0)^{\frac{1}{\sigma }}}SP(S,\overline{f}_0)𝑑S`$ and the scaling $`S(\overline{f}_c\overline{f}_0)^1`$ result in
$$\gamma =\frac{2\tau }{\sigma }=1.$$
(16)
Due to the compactness of avalanches, we have $`S_{co}r_{co}^D=(\mathrm{\Delta }f)^{\nu D}`$, thus
$$\nu =\frac{1}{\sigma D}=\frac{1}{(2\tau )D}.$$
(17)
Eqs. (11)-(12) establish scaling relations among the critical exponents, and they imply that the self-organization time to reach the critical state is independent of the initial configuration of the system. A system of size $`L`$ reaches the stationary state when $`[\mathrm{\Delta }f(s)]^\nu L`$. It can be inferred from Eqs. (11)-(12) that, if one chooses $`\tau `$ and $`D`$ as two independent exponents other exponents can be expressed in terms of them. Among the six exponents mentioned above, $`\tau `$ and $`D`$ can be numerically measured , $`\gamma `$ and $`\rho `$ can be analytically obtained, while $`\nu `$ and $`\sigma `$ are difficult to explore despite some methods measuring the corresponding exponents for $`f_0`$-avalanches were introduced in Ref. . Therefore, we can rely on the scaling relations and values of the exponents obtained to predict the values of $`\nu `$ and $`\sigma `$. We predict $`\sigma =0.2`$ (1D) and 0.275 (2D), $`\nu =2.04`$ (1D) and 1.17 (2D).
Comparing $`\overline{f}_0`$-avalanche with $`f_0`$-avalanche we find that the former is more readily to be treated. Two critical exponents can be analytically obtained and are found to be universal for all dimensional BS models. Furthermore, the infinite hierarchy of exact equations and the exact gap equations, together with their solutions, provide exclusive investigation of the new type of avalanches. Another asset of $`\overline{f}_0`$-avalanche is that it involves some information concerning the whole system. It can be concluded that $`\overline{f}_0`$-avalanche does enable us to comprehend the complex system from an effective and different context. The weak point of this avalanche is that it loses some knowledge of individuals. It is still unknown how these individual features will matter. It is worthwhile to investigate the avalanche dynamics further in the future.
This work was supported in part by NSFC in China and Hubei-NSF. X.C. thanks T. Meng for hospitalities during his visit in Berlin.
Figure Captions
Fig. 1: The average size of avalanches $`S`$ vs $`(\overline{f}_c\overline{f}_0)`$ for (a) 1D and (b) 2D Bak-Sneppen evolution models. The asymptotic slope yields $`\gamma =0.99\pm 0.01`$ and $`0.98\pm 0.01`$ respectively.
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# Alberta-Thy 07-00 May 2000 Holographic Stress Tensor for Kerr-AdS Black Holes and Local Failure on IR-UV Connection
## Abstract
We show that in general holographic stress tensor may contain a new term of divergence of a spacelike unit normal acceleration. Then, it is shown that in contrast to previous descriptions, a new stress tensor for Kerr-AdS solutions can be a traceless one. Interestingly, this prescription entails a local failure on the IR-UV connection.
A precise form of the AdS/CFT correspondence has been formulated by Gubser, Polyakov, Klebanov and Witten . The basic statement is given by the correspondence between the partition function of a string (or M) theory and the generating functional of correlation functions of a boundary conformal field theory (CFT). In the description, the two are functionals of boundary fields $`\varphi ^{(0)}`$ which have dual perspective as boundary values of bulk fields $`\varphi `$ and sources for the operators of CFT. According to this view, in the limit of large number of D-branes and small string coupling, the effective action of strong coupling large $`N`$ CFT is given by evaluating the action functional for solutions to classical supergravity equations of motion
$$S_{string}[\varphi ^{cl}(\varphi ^{(0)})]=W_{CFT}[\varphi ^{(0)}].$$
(1)
Concentrating on the stress tensor in CFT and the corresponding bulk gauge field, metric $`g_{\mu \nu }`$, it may be assumed that all other bulk fields vanish on the boundary. For AdS spaces, due to existence of a second order pole on the asymptotic boundary, the metric $`g_{\mu \nu }`$ does not induce a unique metric $`g^{(0)}`$ on the boundary. Instead, the boundary field that satisfies the boundary condition of the bulk metric is a conformal structure $`[g^{(0)}]`$. However, the conformal invariance is to be broken on the regularization of the bulk supergravity action as arbitrarily picking a particular representative $`g^{(0)}`$ of the conformal structure $`[g^{(0)}]`$. Taking this scheme, Henningson and Skenderis have shown that the conformal anomaly in CFT (ultraviolet (UV) effect) arises from infrared (IR) divergences in the bulk theory. This is an explicit example of the IR-UV connection , which applies to holographic theories , and becomes a non-trivial check for the conjectured AdS/CFT correspondence. Generalizations and applications of the investigation has been studied in .
According to the equation (1), divergences arising from the supergravity action is the usual UV divergences in CFT. Thus, the regularization for the supergravity action can be achieved by introducing local counterterms. Compared to the reference background subtraction method , the prescription, so call counterterm subtraction method, is a nice way to regularize a gravitational action apparently preserving general covariance. The counterterm subtraction method has been developed in its own interest and applications .
On the other hand, various black holes have been studied in the context of the AdS/CFT correspondence, and some interesting observations have been made, e.g., on electrically charged Reissner-Nordström-AdS , rotating Kerr-AdS , and Kerr-Newman-AdS black holes . In this paper, we are concerned about the IR-UV connection between the Kerr-AdS black holes and boundary CFT’s living on rotating Einstein universes. This subject has been served by a manuscript in which the correspondence was probed by calculating the Casimir energies and/or conformal anomalies from bulk theory using the counterterm subtraction method.
An interesting observation in has been made for the five-dimensional Kerr-AdS and the dual $`𝒩=4`$ super Yang-Mills (SYM) theory on a rotating Einstein universe in four dimensions. In usual, when a conformal anomaly is present, the classical bulk action contains a logarithmic divergence, which cannot be apparently canceled by a counterterm. For the five-dimensional Kerr-AdS, the stress tensor was not traceless, but logarithmic divergence did not appear in the on-shell action. Nevertheless, the trace of stress tensor precisely matched to that of the dual SYM. The authors argued that not only the integrated conformal anomaly vanishes, but also the anomaly is proportional to $`\mathrm{}R`$, where $`R`$ is the boundary scalar curvature. Therefore, supplementing ordinary counterterm action with an additional counterterm, then one obtains a new traceless stress tensor. They also proposed that the different choices for counterterms corresponds to the choice of different schemes for regularization in ordinary field theories in four dimensions that is due to the freedom of taking an undetermined coefficient of $`\mathrm{}R`$ in the anomaly.
However, this prescription for the paradox seems unreasonable. First of all, in the one loop effective action of the $`𝒩=4`$ SYM on four-dimensional rotating Einstein universe, the $`\mathrm{}R`$ term in present has to be distinguished from the $`\mathrm{}R`$ term with an undetermined coefficient, e.g., depending on the choices of minimally coupled and conformally coupled scalars . The former is just the usual logarithmic UV divergence in four dimensions $`R^{ab}R_{ab}R^2/3`$. The proportionality, $`\mathrm{}RR^{ab}R_{ab}R^2/3`$, is a special property of the geometry of the four-dimensional rotating Einstein universe<sup>1</sup><sup>1</sup>1Our argument has been given under consideration of the weak coupling calculation. This seems still available in the strong coupling, because the free energy density at weak coupling is only different from that for the strong coupling (and a high temperature limit) up to a constant factor for a leading term .. Therefore, it appears that there is not a precise relationship between addition of new counterterms and the choice of the undetermined coefficient of $`\mathrm{}R`$ in the field theory. Secondly, it has to be noted that the addition of new counterterms means that there may be ‘pulsative counterterms’ which could be turned on and off depending on given boundary geometries and/or topologies. However, considering the procedure of the derivation of counterterm action in , it must be available for all kinds of asymptotic AdS spaces with boundaries of arbitrary geometries and topologies as solutions to the Einstein’s equations (containing only the gravitational field without other bulk fields). Thus, it seems hard to put the pulsative counterterms into the counterterm action with a consistent description. It has to do that just by hand.
In this paper, we revisit this paradox. Our starting point is to elaborate on the on-shell action in the context of the ADM formulation. Taking into account this description, we show that in general the stress tensor may contain a new term of divergence of a spacelike unit normal acceleration and be a traceless one. Then we shall argue that this prescription interestingly may entail a local failure on the IR-UV connection; One loop effective action of the $`𝒩=4`$ SYM on four-dimensional rotating Einstein universe is UV finite, and correspondingly the effective action evaluated from bulk action is IR finite. However, the modified stress tensor derived from bulk theory may be traceless, while the SYM has non-vanishing trace of stress tensor.
$`(d+1)`$-dimensional gravitational action with cosmological constant $`\mathrm{\Lambda }=d(d1)/(2\mathrm{}^2)`$ is given by
$$S=\frac{1}{16\pi G}_Xd^{d+1}x\sqrt{g}\left(\widehat{R}+\frac{d(d1)}{\mathrm{}^2}\right)\frac{1}{8\pi G}_Xd^dx\sqrt{\gamma }\mathrm{\Theta },$$
(2)
where $`X`$ denotes $`d`$-dimensional boundary manifold with metric $`\gamma _{ab}`$ and $`\mathrm{\Theta }_{ab}`$ is the extrinsic curvature of the boundary defined by $`\mathrm{\Theta }_{ab}=\gamma _a^\mu _\mu n_b`$. $``$ denotes the covariant derivative on the $`(d+1)`$-dimensional manifold $`X`$ and $`n^\mu `$ is an outward unit normal to the boundary. $`\widehat{R}`$ is the bulk scalar curvature. The surface term in (2), so called Gibbons-Hawking term, is required for well defined variational principle. In this paper, we consider the bulk metric of the form
$$g_{\mu \nu }dx^\mu dx^\nu =N^2dr^2+\gamma _{ab}dx^adx^b,$$
(3)
where $`x^r`$ is the radial coordinate $`r`$, and $`N^2=N^2(r,x^a)`$, $`\gamma =\gamma (r,x^a)`$. In this coordinate system, the spacelike unit normal to the boundary is given by $`n_\mu =N\delta _\mu ^r`$.
According to the counterterm subtraction method, we introduce a counterterm action $`\stackrel{~}{S}`$ regularizing the action (2)
$`\stackrel{~}{S}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi G}}{\displaystyle _X}d^dx\sqrt{\gamma }\{{\displaystyle \frac{d1}{\mathrm{}}}+{\displaystyle \frac{\mathrm{}}{2(d2)}}R`$ (4)
$`+{\displaystyle \frac{\mathrm{}^3}{2(d2)^2(d4)}}(R_{ab}R^{ab}{\displaystyle \frac{d}{4(d1)}}R^2)+\mathrm{}\}.`$
Then, the regularized action $`S_p`$ is defined by $`S_pS+\stackrel{~}{S}`$. The line element of Kerr-AdS solutions ($`d3`$) interested in this paper is
$`ds^2`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_r}{\rho ^2}}\left(dt{\displaystyle \frac{a\mathrm{sin}^2\theta d\varphi }{\zeta }}\right)^2+{\displaystyle \frac{\mathrm{\Delta }_\theta \mathrm{sin}^2\theta }{\rho ^2}}\left(adt{\displaystyle \frac{(r^2+a^2)}{\zeta }}d\varphi \right)^2`$ (5)
$`+{\displaystyle \frac{\rho ^2}{\mathrm{\Delta }_r}}dr^2+{\displaystyle \frac{\rho ^2}{\mathrm{\Delta }_\theta }}d\theta ^2+r^2\mathrm{cos}^2\theta d\mathrm{\Omega }_{d3}^2,`$
where
$`\rho ^2`$ $`=`$ $`r^2+a^2\mathrm{cos}^2\theta ,\mathrm{\Delta }_r=(r^2+a^2)(1+r^2/\mathrm{}^2)2mGr^{4d},`$
$`\zeta `$ $`=`$ $`1a^2/\mathrm{}^2,\mathrm{\Delta }_\theta =1(a^2/\mathrm{}^2)\mathrm{cos}^2\theta ,`$ (6)
and $`m`$ and $`a`$ denote the black hole mass and angular momentum per unit mass, respectively. This is an AdS version of higher dimensional Kerr black holes .
The on-shell regularized action $`S_p^{cl}`$ of the five-dimensional Kerr-AdS in (5) does not contain a logarithmic divergence . Divergent part of the on-shell action is given by
$`S_{div}^{cl}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi G}}{\displaystyle _X}d^{d+1}x\sqrt{g}{\displaystyle \frac{d}{\mathrm{}^2}}{\displaystyle \frac{1}{8\pi G}}{\displaystyle _X}d^dx\sqrt{\gamma }\mathrm{\Theta },`$
$`=`$ $`{\displaystyle _X}d^dx{\displaystyle \frac{\sqrt{\mathrm{\Omega }_{d3}}}{8\pi G}}r^{d2}[{\displaystyle \frac{(d1)}{\mathrm{}^2}}r^2+(d1)(1+{\displaystyle \frac{a^2}{\mathrm{}^2}}(1{\displaystyle \frac{2\mathrm{cos}^2\theta }{d2}}))`$
$`{\displaystyle \frac{a^4}{r^4}}\mathrm{cos}^2\theta \mathrm{sin}^2\theta \mathrm{\Delta }_\theta +\mathrm{}+{\displaystyle \frac{a^2(a^2\mathrm{cos}^2\theta )^{(d6)/2}}{r^{d4}}}\mathrm{sin}^2\theta \mathrm{\Delta }_\theta ]{\displaystyle \frac{\mathrm{sin}\theta \mathrm{cos}^{d3}\theta }{\zeta }}.`$
In ( Alberta-Thy 07-00 May 2000 Holographic Stress Tensor for Kerr-AdS Black Holes and Local Failure on IR-UV Connection) and hereafter, we set $`m=0`$ in the metric (5). The terms including the mass is in fact finite on the asymptotic region and is irrelevant to aim of this paper. In addition, it is a necessary condition for counterterms that must be given in terms of only intrinsic boundary geometry.
On the other hand, the divergence structure of the on-shell action is tightly constrained by the Gauss-Codazzi equations in the sense that these are covariant expressions given in terms of the intrinsic and extrinsic boundary geometry . Thus, we expect that the ADM formulation gives us a hint for resolving the paradox above mentioned. In fact, using the ADM formulation, the on-shell action can be expressed by only the intrinsic boundary geometry up to redefinition of the radial coordinate . In this sense, we calculate the on-shell action again in the context of the ADM formulation.
The canonical form of the action (2) is
$$S=\frac{1}{16\pi G}_Xd^{d+1}xN\sqrt{\gamma }\left(\mathrm{\Theta }^2\mathrm{\Theta }^{ab}\mathrm{\Theta }_{ab}+R+\frac{d(d1)}{\mathrm{}^2}\right).$$
(8)
On the other hand, the Einstein equations contracted by the bulk metric can be written by
$$\mathrm{\Theta }^2\mathrm{\Theta }^{ab}\mathrm{\Theta }_{ab}R\frac{d(d1)}{\mathrm{}^2}=0,$$
(9)
and
$$\mathrm{\Theta }^{ab}\mathrm{\Theta }_{ab}n^\mu _\mu \mathrm{\Theta }_\mu b^\nu \frac{d}{\mathrm{}^2}=0,$$
(10)
where $`b^\mu n^\nu _\nu n^\mu `$ is an acceleration of the unit normal $`n^\mu `$. The first equation (9) is the normal-normal component of the Gauss-Codazzi equations and the second (10) can be identified with the tangential-tangential one (requiring the equation (9)). Substituting (9) into (8), the on-shell action is given by
$$S^{cl}=\frac{1}{8\pi G}_Xd^{d+1}xN\sqrt{\gamma }\left(R+\frac{d(d1)}{\mathrm{}^2}\right).$$
(11)
It must be noted that since we are concerned about the divergence structure of the on-shell action, the equation (10) is irrelevant in our calculation . However, the term of divergence of the acceleration in (10) is to play an important role in our prescription.
Now, we find that divergent part of the on-shell action (11) for $`d`$-dimensional Kerr-AdS solutions ($`d=4,6,\mathrm{}`$) contains a logarithmic term
$`S_{div}^{cl}`$ $`=`$ $`{\displaystyle _X}d^dx{\displaystyle \frac{\sqrt{\mathrm{\Omega }_{d3}}}{8\pi G}}r^{d2}[{\displaystyle \frac{(d1)}{\mathrm{}^2}}r^2+(d1)(1+{\displaystyle \frac{a^2}{\mathrm{}^2}}(1{\displaystyle \frac{2\mathrm{cos}^2\theta }{d2}}))`$ (12)
$`+{\displaystyle \frac{a^2}{r^2}}\left(d3+{\displaystyle \frac{2((d3)\mathrm{sin}^2\theta \mathrm{cos}^2\theta )}{d4}}{\displaystyle \frac{2a^2\mathrm{cos}^2\theta ((d2)\mathrm{sin}^2\theta \mathrm{cos}^2\theta )}{\mathrm{}^2(d4)}}\right)`$
$`{\displaystyle \frac{a^4}{r^4}}\left({\displaystyle \frac{2\mathrm{cos}^2\theta ((d3)\mathrm{sin}^2\theta \mathrm{cos}^2\theta )}{d6}}{\displaystyle \frac{2a^2\mathrm{cos}^4\theta ((d2)\mathrm{sin}^2\theta \mathrm{cos}^2\theta )}{\mathrm{}^2(d6)}}\right)`$
$`+\mathrm{}+2(\mathrm{cos}\theta )^{(d4)}\mathrm{ln}r\left({\displaystyle \frac{a}{r}}\right)^{(d4)}((d3)\mathrm{sin}^2\theta \mathrm{cos}^2\theta `$
$`{\displaystyle \frac{a^2\mathrm{cos}^2\theta }{\mathrm{}^2}}((d2)\mathrm{sin}^2\theta \mathrm{cos}^2\theta ))]{\displaystyle \frac{\mathrm{sin}\theta \mathrm{cos}^{d3}\theta }{\zeta }}.`$
The logarithmic divergence term in (12) apparently leads a conformal anomaly
$`𝒜`$ $`=`$ $`{\displaystyle \frac{a^2(a^2\mathrm{cos}^2\theta )^{(d4)/2}}{8\pi G}}\times `$ (13)
$`\left({\displaystyle \frac{(d3)\mathrm{sin}^2\theta \mathrm{cos}^2\theta a^2\mathrm{cos}^2\theta ((d2)\mathrm{sin}^2\theta \mathrm{cos}^2\theta )/\mathrm{}^2}{\rho r^{d3}\mathrm{\Delta }_r^{(m=0)}}}\right),`$
where we used a cutoff $`r^2/\mathrm{}^2`$ (c.f. ). As expected, for the case of $`d=4`$ the conformal anomaly in the leading contribution, $`𝒜_{d=4}`$, is equal to that evaluated in
$$𝒜_{d=4}=\frac{a^2\mathrm{}}{8\pi Gr^4}\left(\frac{a^2\mathrm{cos}^2\theta }{\mathrm{}^2}(3\mathrm{cos}^2\theta 2)\mathrm{cos}2\theta \right).$$
(14)
Finally, we are in position of describing the discrepancy of the on-shell actions in ( Alberta-Thy 07-00 May 2000 Holographic Stress Tensor for Kerr-AdS Black Holes and Local Failure on IR-UV Connection) and (12). Following the above observation, especially deriving the conformal anomaly (14), the discrepancy should be closely related to the paradox given in why the stress tensor is not traceless, while the on-shell action ( Alberta-Thy 07-00 May 2000 Holographic Stress Tensor for Kerr-AdS Black Holes and Local Failure on IR-UV Connection) does not contain a logarithmic divergence.
The origin of the discrepancy is found in the canonical form of the action (8). In fact, it contained two total derivative terms, one canceled the Gibbons-Hawking term, and the other was discarded by a simple algebraic relation given by
$`S_{bt}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi G}}{\displaystyle _X}d^{d+1}x\sqrt{g}_\mu (n^\nu _\nu n^\mu )`$ (15)
$`=`$ $`{\displaystyle \frac{1}{8\pi G}}{\displaystyle _X}d^dx\sqrt{\gamma }n_\mu (n^\nu _\nu n^\mu )=0.`$
However, it is easily shown that the above calculation is not correct. In the coordinate system of (3), the divergence of acceleration $`_\mu b^\mu `$ cannot be a surface term of the timelike boundary $`X`$, because it is not given by a total derivative term of the radial coordinate
$$\sqrt{g}_\mu b^\mu =_a(\sqrt{\gamma }\gamma ^{ab}_bN)=\sqrt{\gamma }D^aD_aN,$$
(16)
where $`D_a`$ is the covariant derivative defined on the timelike boundary. Thus, we have to keep this term in calculation of the on-shell action (11)
$$S^{cl}=\frac{1}{8\pi G}_Xd^{d+1}xN\sqrt{\gamma }\left(R+\frac{d(d1)}{\mathrm{}^2}+_\mu b^\mu \right).$$
(17)
The on-shell action (17) for Kerr-AdS solutions does not contain the logarithmic divergence term and recover that in ( Alberta-Thy 07-00 May 2000 Holographic Stress Tensor for Kerr-AdS Black Holes and Local Failure on IR-UV Connection). According to this observation, it appears that the stress tensor is modified by the term of divergence of unit normal acceleration and becomes a traceless one.
Definition of the stress tensor is
$$T^{ab}\frac{2}{\sqrt{\gamma }}\frac{\delta S^{cl}}{\delta \gamma _{ab}}=\frac{1}{8\pi G}\left(\mathrm{\Theta }^{ab}\gamma ^{ab}\mathrm{\Theta }\right).$$
(18)
Actually, in the second equality, it was assumed that $`\gamma ^{\mu \nu }\delta n_\mu =\delta \gamma ^{\mu \nu }n_\mu =0`$. This means that the boundary is fixed under the variations so that the variations of the normal dual-vector on the boundary are proportional to the normal dual-vector. (For an example, see .) However, this assumption is no longer proper for the type of metric (3), e.g., Kerr-AdS solutions, in which the radial lapse $`N`$ is a function of a boundary coordinate as well as the radial one, and this restriction has to be relaxed. Unfortunately, taking the relaxation of the assumption, direct calculation of the stress tensor seems not easy because of the particular algebraic form of the divergence of the acceleration. Taking into account the form of the on-shell action 17, we propose a form of new stress tensor as
$`T_{new}^{ab}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi G}}\left(\mathrm{\Theta }^{ab}\gamma ^{ab}\mathrm{\Theta }+{\displaystyle \frac{\alpha }{2}}\gamma ^{ab}{\displaystyle 𝑑rN_\mu b^\mu }\right)`$ (19)
$`=`$ $`{\displaystyle \frac{1}{8\pi G}}\left(\mathrm{\Theta }^{ab}\gamma ^{ab}\mathrm{\Theta }\alpha \gamma ^{ab}K\right),`$
where $`K\frac{1}{2}𝑑rD^cD_cN`$. (In the following, we are calling the terms that are just intuitively related to the $`K`$ term as ‘$`K`$ term’.) If the constant $`\alpha `$ in (19) is one, then the new regularized stress tensor $`T_p^{ab}T_{new}^{ab}+\stackrel{~}{T}^{ab}`$ is traceless.
The usual scheme, which an on-shell action does not contain a logarithmic divergence then the trace of stress tensor vanishes, make us lead to take $`\alpha =1`$. In fact, there is another reason why $`\alpha =1`$ is attractive. Kerr black hole solutions, which are asymptotically flat spacetime, can be obtained by taking the flat spacetime limit $`\mathrm{}\mathrm{}`$ in 5. In the case, one can see that the ‘old’ definition of the stress tensor (18) has a non-vanishing trace in leading contribution, which is the same with the anomaly in (13) taking the flat spacetime limit . (In the calculation, the counterterm action for $`d=4`$ is the form of
$$\stackrel{~}{S}=\frac{1}{8\pi G}_Xd^4x\sqrt{\gamma }\sqrt{\frac{3}{2}R},$$
(20)
(For the $`d=4`$ case, the counterterm action (20) is enough to eliminate the divergence appearing in the classical action<sup>2</sup><sup>2</sup>2For higher dimensional Kerr solutions, see .) and the counterterm stress tensor $`\stackrel{~}{T}^{ab}`$ is given by
$$\stackrel{~}{T}^{ab}=\frac{1}{8\pi G}\left(\mathrm{\Phi }(R^{ab}\gamma ^{ab}R)+D^aD^b\mathrm{\Phi }D^cD_c\mathrm{\Phi }\gamma ^{ab}\right),$$
(21)
where $`\mathrm{\Phi }\sqrt{3/(2R)}`$.) Taking into account the holographic principle and extrapolating duality in the AdS/CFT correspondence, we expact that there might be a quantum field theory living on the boundary that is dual to a (super)gravity on an asymptotically flat spacetime. However, it should not be a conformal field theory. Then, the stress tensor with non-vanishing trace would be problematic. In this sense, the choice of traceless stress tensor seems to be reasonable. The new stress tensor in (19) with $`\alpha =1`$ taking the asymptotically flat limit $`\mathrm{}\mathrm{}`$ gives a traceless regularized stress tensor.
Even though the new traceless stress tensor with $`\alpha =1`$ (19) is plausible with the fact that the on-shell action does not contain logarithmic divergence, it seems to be problematic on the AdS/CFT correspondence. As mentioned above, the one loop effective action of the $`𝒩=4`$ SYM on four-dimensional rotating Einstein universe is UV finite, so the equation (1) is still satisfied. However, the SYM has non-vanishing conformal anomaly. In this sense, it appears that the IR-UV connection is locally broken. In fact, this local failure occurs in the case of $`\alpha 0`$. In order words, the contribution of the $`K`$ term to the stress tensor gives rise to the local failure on the IR-UV connection.<sup>3</sup><sup>3</sup>3It must be noted that the $`K`$ term does not contribute to the total energy of the bulk theory in the leading contribution, and the Casimir energy derived from the bulk theory still matches to that of the boundary dual CFT. The $`K`$ term in the leading contribution is propertional to squared angular momentum, and apparently, the local failure is due to rotation of the bulk spacetime. In order to discuss this problem more, we consider some aspects of the $`K`$ term in canonical point of view.
In some sense, the $`K`$ term measures how much deviated the boundary geometry is from a round sphere. This reflects that the contribution appears in the tangential-tangential component of the Gauss-Codazzi equations (10). On the other hand, the canonical form of the action (8) including the $`K`$ term is written in terms of canonical variables
$$S=_Xd^{d+1}x\left(\pi ^{ab}\gamma _{ab}^{}N\frac{\sqrt{\gamma }}{8\pi G}D^aD_aN\right),$$
(22)
where $`\pi ^{ab}`$ is the conjugate momenta defined by $`\pi ^{ab}\delta S/\delta \gamma _{ab}^{}`$. The radial Hamiltonian density $``$ is given by
$$=\frac{16\pi G}{\sqrt{\gamma }}\left(\frac{\pi ^2}{d1}\pi _{ab}\pi ^{ab}\right)\frac{\sqrt{\gamma }}{16\pi G}\left(R+\frac{d(d1)}{\mathrm{}^2}\right).$$
(23)
Now, how can we understand the $`K`$ term in the canonical action (22)? First of all, the canonical action (22) can be rewritten by
$$S=𝑑r\left[_Xd^dx\left(\pi ^{ab}\gamma _{ab}^{}N\right)\frac{1}{8\pi G}_Xd^{d1}x\sqrt{|h|}u^aD_aN\right],$$
(24)
where $`h`$ is an induced metric of a $`(d1)`$-dimensional boundary $`X`$ and $`u^a`$ is a unit normal to the $`X`$. Thus, the $`K`$ term becomes a surface term of the Hamiltonian $`H`$
$$H=_Xd^dx+\frac{1}{8\pi G}_Xd^{d1}x\sqrt{|h|}u^aD_aN.$$
(25)
In usual, a surface term of a Hamiltonian plays an important role in physics, e.g., as the total energy of a system. In this paper, we have not find a physical description for this surface term. According to the AdS/CFT correspondence, the Hamiltonian constraint $`=0`$ (turning on bulk scalar fields) is equivalent to the renormalization group flow (RG-flow) equation of the boundary CFT . However, the surface term should not give any contribution on the CFT deformation, and moreover, we have been concentrated only on the asymptotic boundary in which bulk scalar fields vanish. Thus, even though it will be find that the surface term in (25) give a kind of (local) deformation of boundary dual CFT related to the local failure on the IR-UV connection, it should be different from the CFT deformations recently studied in the holographic RG-flow (For review, see and references therein). We leave the investigation of prescription for a possible description for the surface term in (25) and its relationship with the local failure on the IR-UV connection as a future work.
Acknowledgments: I thank A. Zelnikov, Y. Gusev, and S. Corley for helpful discussions. This work was supported by National Science and Engineering Research Council of Canada.
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# Generalized Field Theoretical Approach to General Relativity and Conserved Quantities in Anti-de Sitter Spacetimes
## I Introduction
In the last years, a lot of investigations have been done on asymptotically anti-de Sitter (AdS) spacetimes. This interest is widely connected with the development of string theory. For instance, the Maldacena conjecture relates conformal field theories in a $`d`$-dimensional space with supergravity or string theory on the product of $`(d+1)`$-dimensional AdS space with a compact manifold, and this result can be used to study the thermodynamics of AdS black holes . One important problem of AdS space concerns the definition of mass and angular momentum of asymptotically AdS spacetimes. Many papers have been written on this subject proposing different ways to calculate such quantities . Classically, they all agree in their final results, with the exception of the generalized Komar mass calculated in Ref. . In this paper we will propose a new method to calculate such quantities based on the field theoretical approach to General Relativity (GR) , which is applicable to either asymptotically flat or asymptotically AdS spacetimes for all $`d2`$, and discuss its advantages and drawbacks with respect to the other methods. Our development also clarifies some issues concerning the background metric defined in Refs. .
In the field theoretical approach to GR, the gravitational field is treated as a spin two field $`h^{\mu \nu }`$ propagating with self interaction in a background metric $`\gamma _{\mu \nu }`$, which is usually considered to be Ricci flat. The self interaction and interaction of the spin two field with other matter fields is constructed in such a way that the background metric is never perceived: it appears only through the combination $`\sqrt{g}g^{\mu \nu }=\sqrt{\gamma }(\gamma ^{\mu \nu }+h^{\mu \nu })`$, and this combination satisfies exactly the Einstein’s equations. In terms of $`h^{\mu \nu }`$ and $`\gamma _{\mu \nu }`$, the Einstein-Hilbert action and Einstein’s equation can be viewed as an action and equations of motion of a spin two field in a Ricci flat background metric. This allows one to define a true energy-momentum tensor of the gravitational field plus matter fields by varying the lagrangian with respect to the background metric. However, the theory in terms of $`h^{\mu \nu }`$ aquires a gauge freedom which is the manifestation of the invariance of GR under the action of the manifold mapping group (MMG), and the energy momentum tensor calculated in the way explained above is not a gauge invariant quantity. Nevertheless, quantities like the total energy and angular momentum do not suffer from this ambiguity once one specify asymptotic conditions. In fact, some calculations of total energies have been done for asymptotically flat spacetimes agreeing with well known results. Our task is to extend this procedure to asymptotically AdS spacetimes. This extension, however, is not straightforward. This is because of the presence of the cosmological constant in AdS spacetimes. One might think that this should not pose a problem because the formalism developped in Refs. are applicable to any matter field, and a cosmological constant can be viewed, for instance, as resulting from the energy density of some scalar field in its ground state. However, when we calculate the total energy of asymptotically AdS spacetimes considering that the background metric is still Ricci flat, we obtain preposterous results, even when we remove the energy of the pure AdS spacetime. The resulting energy does not yield the usual asymptotically flat results when we put the cosmological constant to zero. This result enforces us to take as the background metric some Einstein space satisfying $`\stackrel{}{R}_{\mu \nu }=2\gamma _{\mu \nu }\stackrel{}{\mathrm{\Lambda }}/(d1)`$.
In general, the total energy and angular momenta are calculated by the integration of a closed $`d`$-form $`J`$, $`dJ=0`$, on a d-dimensional spacelike volume. The main result of this paper is to exhibit the general expression of a $`(d1)`$-form $`\mathrm{\Omega }`$ which fulfills the condition $`J=d\mathrm{\Omega }`$ for the general case of background metrics satisfying $`\stackrel{}{R}_{\mu \nu }=2\gamma _{\mu \nu }\stackrel{}{\mathrm{\Lambda }}/(d1)`$, hence proving that $`J`$ is indeed globally exact in general. This enables one to write those conserved quantities as asymptotic $`(d1)`$-dimensional integrals, which rends their calculations feasible, and clarify the conditions for their gauge invariance. It turns out that the $`(d1)`$-form $`\mathrm{\Omega }`$ does not depend neither on $`\stackrel{}{\mathrm{\Lambda }}`$ nor on $`d`$. With the $`(d1)`$-form $`\mathrm{\Omega }`$ we have at our disposal an alternative and straightforward method to calculate any conserved quantity of any asymptotically flat or AdS spacetime in any number of dimensions greater than two. We calculated total energies and angular momenta for some known solutions in four and five dimensions using our expression for $`\mathrm{\Omega }`$, and they agree with standard results.
This paper is organized as follows: in the next section we generalize the field theoretical approach to the case of GR with a cosmological constant $`\mathrm{\Lambda }`$ and background metric satisfying $`\stackrel{}{R}_{\mu \nu }=2\gamma _{\mu \nu }\stackrel{}{\mathrm{\Lambda }}/(d1)`$ (note that, for the sake of generality, the two cosmological constants may be different). We obtain the total energy-momentum tensor. In section III we write down the conserved current $`J`$ associated with this total energy-momentum tensor, and we arrive at the $`(d1)`$-form $`\mathrm{\Omega }`$. We discuss the conditions for its gauge invariance, and we calculate some conserved quantities of well known examples using this method. We also discuss the necessity of using background metrics satisfying $`\stackrel{}{R}_{\mu \nu }=2\gamma _{\mu \nu }\stackrel{}{\mathrm{\Lambda }}/(d1)`$ when the effective geometry is asymptotically AdS with $`\stackrel{}{\mathrm{\Lambda }}`$ being the same as the effective cosmological constant in the effective geometry. We end up in section IV with conclusions and discussions.
## II The field theoretical approach for General Realtivity with a cosmological constant
Take the Einstein-Hilbert action depending on a metric $`g_{\mu \nu }`$ and a symmetric connection $`\mathrm{\Gamma }_{\mu \nu }^\alpha `$, together with a cosmological constant $`\mathrm{\Lambda }`$ term, and an action for matter fields without couplings with the connection. In this total action, make the substitutions
$$\sqrt{g}g^{\mu \nu }=\sqrt{\gamma }(\gamma ^{\mu \nu }+h^{\mu \nu }),$$
(1)
$$\mathrm{\Gamma }_{\mu \nu }^\alpha =C_{\mu \nu }^\alpha +K_{\mu \nu }^\alpha .$$
(2)
where $`\gamma _{\mu \nu }`$ and $`C_{\mu \nu }^\alpha `$ are the metric and Christoffel symbols, respectively, of a background space satisfying
$$\stackrel{}{R}_{\mu \nu }=\frac{2\stackrel{}{\mathrm{\Lambda }}}{d1}\gamma _{\mu \nu }$$
(3)
From now on indices will be raised and lowered by $`\gamma _{\mu \nu }`$. After discarding surface terms and terms independent on the fields $`h^{\mu \nu }`$ and $`K_{\mu \nu }^\alpha `$, we obtain:
$`S`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G}}{\displaystyle }[\stackrel{~}{h}^{\mu \nu }(K_{\mu \nu ;\alpha }^\alpha K_{\mu ;\nu })+(\stackrel{~}{\gamma }^{\mu \nu }+\stackrel{~}{h}^{\mu \nu })(KK)_{\mu \nu }`$ (5)
$`+{\displaystyle \frac{2}{d1}}\stackrel{}{\mathrm{\Lambda }}\stackrel{~}{h}^{\mu \nu }\gamma _{\mu \nu }2\mathrm{\Lambda }\sqrt{\overline{g}}]d^{d+1}x+{\displaystyle }_M(\mathrm{\Phi }_A,\overline{g}_{\alpha \beta })d^{d+1}x,`$
where
$$\stackrel{~}{\gamma }^{\mu \nu }\sqrt{\gamma }\gamma ^{\mu \nu },\stackrel{~}{h}^{\mu \nu }\sqrt{\gamma }h^{\mu \nu },$$
(6)
$$K_\mu K_{\alpha \mu }^\alpha ,$$
(7)
and
$$(KK)_{\mu \nu }K_{\mu \nu }^\alpha K_\alpha K_{\mu \beta }^\alpha K_{\nu \alpha }^\beta .$$
(8)
The semi-colon (;) represents the covariant derivative with respect to the background connections, and $`\mathrm{\Phi }_A`$ represents general matter fields. The terms $`\overline{g}_{\alpha \beta }`$, which is the inverse of $`g^{\mu \nu }`$, and $`\sqrt{\overline{g}}`$ appearing in the action must be understood as algebraic functions of finite degree polynomials on $`h^{\mu \nu }`$ and $`\gamma _{\mu \nu }`$ (see Ref. for their explicit form).
The dynamical equations for the gravitational field can be obtained by performing a Palatini-like variation of the action (5) with respect to the fields $`\stackrel{~}{h}^{\mu \nu }`$ and $`K_{\mu \nu }^\alpha `$ independently. We obtain the following equations:
$`K_{\mu \nu ;\alpha }^\alpha {\displaystyle \frac{1}{2}}(K_{\mu ;\nu }+K_{\nu ;\mu })+(KK)_{\mu \nu }=`$ (9)
$`{\displaystyle \frac{2}{d1}}(\mathrm{\Lambda }\overline{g}_{\mu \nu }\stackrel{}{\mathrm{\Lambda }}\gamma _{\mu \nu })+8\pi G\left(T_{\mu \nu }^M{\displaystyle \frac{1}{d1}}T^M\gamma _{\mu \nu }\right),`$ (10)
and
$$\stackrel{~}{h}_{;\alpha }^{\mu \nu }(\stackrel{~}{\gamma }^{\mu \nu }+\stackrel{~}{h}^{\mu \nu })K_\alpha +(\stackrel{~}{\gamma }^{\mu \rho }+\stackrel{~}{h}^{\mu \rho })K_{\rho \alpha }^\nu +(\stackrel{~}{\gamma }^{\nu \rho }+\stackrel{~}{h}^{\nu \rho })K_{\rho \alpha }^\mu =0,$$
(11)
where
$$T_{\mu \nu }^M=\frac{2}{\sqrt{\gamma }}\frac{\delta _M}{\delta \gamma ^{\mu \nu }},$$
(12)
and $`T^M`$ is its trace. Note that the matter energy-momentum tensor defined above is not the same as the usual matter energy-momentum tensor $`t_{\mu \nu }^M`$ defined using variations of the matter lagrangian with respect to the effective metric $`g^{\mu \nu }`$. Only the combinations
$$T_{\mu \nu }^M\frac{1}{d1}T_{\alpha \beta }^M\gamma ^{\alpha \beta }\gamma _{\mu \nu }=2\frac{\delta _M}{\delta \stackrel{~}{\gamma }^{\mu \nu }}=2\frac{\delta _M}{\delta \stackrel{~}{g}^{\mu \nu }}=t_{\mu \nu }^M\frac{1}{d1}t_{\alpha \beta }^Mg^{\alpha \beta }g_{\mu \nu }$$
(13)
are equal. Note also that the $`T_{\mu \nu }^M`$ defined above is not covariantly conserved with respect to the background connection.
Eq. (11) is equivalent to the determination of the connection $`\mathrm{\Gamma }_{\mu \nu }^\alpha `$ in terms of the Christoffel symbols of $`g_{\mu \nu }`$ while Eq. (9) is equivalent to the Einstein’s equations when we use Eqs. (1,2,3). After some rearrangements in Eqs. (9,11) (see Ref. for details), we obtain the following equation:
$`G_{\mu \nu }^L`$ $`=`$ $`(KK)_{\mu \nu }+{\displaystyle \frac{1}{2}}\gamma _{\mu \nu }(KK)_\alpha ^\alpha +Q_{\mu \nu ;\alpha }^\alpha `$ (15)
$`+{\displaystyle \frac{\mathrm{\Lambda }}{d1}}(2\overline{g}_{\mu \nu }\overline{g}_\alpha ^\alpha \gamma _{\mu \nu })+\stackrel{}{\mathrm{\Lambda }}\gamma _{\mu \nu }+8\pi GT_{\mu \nu }^M,`$
where
$$2G_{\mu \nu }^L=[\gamma _{\mu \nu }h^{\alpha \beta }+\gamma ^{\alpha \beta }h_{\mu \nu }\delta _\mu ^\alpha h_\nu ^\beta \delta _\nu ^\alpha h_\mu ^\beta ]_{;\alpha ;\beta },$$
(16)
and
$$\begin{array}{ccc}\hfill 2Q_{\mu \nu }^\alpha & =& \gamma _{\mu \nu }h^{\rho \sigma }K_{\rho \sigma }^\alpha +h_{\mu \nu }K^\alpha h_\mu ^\alpha K_\nu h_\nu ^\alpha K_\mu +\hfill \\ & +& h_\mu ^\rho (K_{\rho \nu }^\alpha K_{\rho \lambda }^\sigma \gamma ^{\alpha \lambda }\gamma _{\sigma \nu })+h_\nu ^\rho (K_{\rho \mu }^\alpha K_{\rho \lambda }^\sigma \gamma ^{\alpha \lambda }\gamma _{\sigma \mu })+\hfill \\ & +& h^{\alpha \rho }(K_{\rho \mu }^\sigma \gamma _{\sigma \nu }+K_{\rho \nu }^\sigma \gamma _{\sigma \mu }),\hfill \end{array}$$
(17)
The total energy-momentum tensor of the gravitational field plus matter described by the action (5) can be calculated as usual by varying the total lagrangian $`_T=_G+_M`$ with respect to the background metric, yielding
$`T_{\mu \nu }^T`$ $`=`$ $`{\displaystyle \frac{2}{\sqrt{\gamma }}}{\displaystyle \frac{\delta _T}{\delta \gamma ^{\mu \nu }}}={\displaystyle \frac{1}{8\pi G}}[(KK)_{\mu \nu }+{\displaystyle \frac{1}{2}}\gamma _{\mu \nu }(KK)^\alpha _\alpha +Q^\alpha _{\mu \nu ;\alpha }`$ (19)
$`+{\displaystyle \frac{\mathrm{\Lambda }}{d1}}(2\overline{g}_{\mu \nu }\overline{g}_\alpha ^\alpha \gamma _{\mu \nu })+{\displaystyle \frac{2\stackrel{}{\mathrm{\Lambda }}}{d1}}h_{\mu \nu }]+T_{\mu \nu }^M.`$
Equations (15) and (19) give rise to
$$G_{\mu \nu }^L+\frac{2\stackrel{}{\mathrm{\Lambda }}}{d1}h_{\mu \nu }=8\pi GT_{\mu \nu }^T+\stackrel{}{\mathrm{\Lambda }}\gamma _{\mu \nu }.$$
(20)
Due to the Riemannian nature of the background geometry, one can define a new conseved energy-momentum tensor as
$$T_{\mu \nu }=T_{\mu \nu }^T+\frac{\stackrel{}{\mathrm{\Lambda }}}{8\pi G}\gamma _{\mu \nu },$$
(21)
and write Eq. (20) as
$$G_{\mu \nu }^L+\frac{2\stackrel{}{\mathrm{\Lambda }}}{d1}h_{\mu \nu }=8\pi GT_{\mu \nu }.$$
(22)
There are some advantages in adopting $`T_{\mu \nu }`$ instead of $`T_{\mu \nu }^T`$. For instance, let us assume that the background metric is a solution of the Einstein’s equations when the matter fields are in the vacuum state, which is equivalent to say that $`h_{\mu \nu }=0`$ is a solution of Eqs. (9) and (11) in vacuum. This is a reasonable assumption, although not mandatory at this moment. In this case, $`T_{\mu \nu }(h_{\mu \nu }=0)`$ is identically zero while $`T_{\mu \nu }^T(h_{\mu \nu }=0)`$ has a residual background term, as can be seen imediately from Eqs. (20,22). As we will see in the next section, the new $`T_{\mu \nu }`$ is a renormalized energy-momentum tensor yielding authomatically finite results for conserved quantities.
Note that the covariant divergence of the left-hand-side of Eq. (22) is identically zero once the background metric satisfies Eq. (3). Hence, the conservation of $`T_{\mu \nu }`$ can also be obtained as a consequence of the field equations.
As the energy-momentum tensor $`T_{\mu \nu }`$ is a true tensor, it seems that it does not suffer from the ambiguities which are present in the usual definitions of pseudo-tensors in GR. However, the field theoretical approach to GR has an invariance under true gauge transformations that comes from the invariance of GR under the action of the manifold mapping group. The $`T_{\mu \nu }`$ above defined is not a gauge invariant quantity, as it will now be seen.
The coordinate transformation invariance of GR is translated to invariance under gauge transformations on $`h^{\mu \nu }`$ and $`K_{\mu \nu }^\alpha `$ in the field theoretical approach in the following way: consider the infinitesimal coordinate transformation $`x^\alpha =x^\alpha +\xi ^\alpha (x)`$, which changes the functional form of $`\stackrel{~}{g}^{\mu \nu }`$ as
$$\stackrel{~}{g}^{\mu \nu }(x)=\stackrel{~}{g}^{\mu \nu }(x)+\mathrm{\pounds }_\xi ^{(1)}\stackrel{~}{g}^{\mu \nu }(x),$$
(23)
where the Lie derivative $`\mathrm{\pounds }_\xi ^{(1)}\stackrel{~}{g}^{\mu \nu }`$ is given by
$$\mathrm{\pounds }_\xi ^{(1)}\stackrel{~}{g}^{\mu \nu }(x)=\stackrel{~}{g}_{,\lambda }^{\mu \nu }\xi ^\lambda +\stackrel{~}{g}^{\lambda \mu }\xi _{,\lambda }^\nu +\stackrel{~}{g}^{\lambda \nu }\xi _{,\lambda }^\mu \stackrel{~}{g}^{\mu \nu }\xi _{,\sigma }^\sigma .$$
(24)
In the case of a finite transformation, the change in $`\stackrel{~}{g}^{\mu \nu }`$ is given by
$$\stackrel{~}{g}^{\mu \nu }(x)=\stackrel{~}{g}^{\mu \nu }(x)+\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k!}\mathrm{\pounds }_\xi ^{(k)}\stackrel{~}{g}^{\mu \nu },$$
(25)
where $`\mathrm{\pounds }_\xi ^{(k)}`$ is the Lie derivative of order $`k`$ defined as
$$\mathrm{\pounds }_\xi ^{(k)}=\mathrm{\pounds }_\xi ^{(1)}[\mathrm{\pounds }_\xi ^{(k1)}].$$
(26)
Substituting in (25) the definition (1) we get
$$\stackrel{~}{g}^{\mu \nu }(x)=\stackrel{~}{\gamma }^{\mu \nu }(x)+\stackrel{~}{h}^{\mu \nu }(x)+\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k!}\mathrm{\pounds }_\xi ^{(k)}(\stackrel{~}{\gamma }^{\mu \nu }+\stackrel{~}{h}^{\mu \nu }).$$
(27)
The transformed metric density can be decomposed in two distinct ways:
$$\stackrel{~}{g}^{\mu \nu }(x)=\stackrel{~}{\gamma }^{\mu \nu }(x)+\stackrel{~}{h}^{\mu \nu }(x)$$
(28)
and
$$\stackrel{~}{g}^{\mu \nu }(x)=\stackrel{~}{\gamma }^{\mu \nu }(x)+\stackrel{~}{h}^{\mu \nu }(x),$$
(29)
where, by comparison with (27) one gets
$$\stackrel{~}{h}^{\mu \nu }(x)=\stackrel{~}{h}^{\mu \nu }(x)+\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k!}\mathrm{\pounds }_\xi ^{(k)}(\stackrel{~}{\gamma }^{\mu \nu }+\stackrel{~}{h}^{\mu \nu }),$$
(30)
$$\stackrel{~}{h}^{\mu \nu }(x)=\stackrel{~}{h}^{\mu \nu }(x)+\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k!}\mathrm{\pounds }_\xi ^{(k)}\stackrel{~}{h}^{\mu \nu },$$
(31)
and
$$\stackrel{~}{\gamma }^{\mu \nu }(x)=\stackrel{~}{\gamma }^{\mu \nu }(x)+\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k!}\mathrm{\pounds }_\xi ^{(k)}\stackrel{~}{\gamma }^{\mu \nu }.$$
(32)
Analogously, the gauge transformation corresponding to Eq. (28) for the field $`K_{\mu \nu }^\alpha (x)`$ is
$$K_{\mu \nu }^\alpha (x)=K_{\mu \nu }^\alpha (x)+\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k!}\mathrm{\pounds }_\xi ^{(k)}(C_{\mu \nu }^\alpha +K_{\mu \nu }^\alpha ),$$
(33)
while for the case corresponding to the transformation (29) reads
$$K_{\mu \nu }^\alpha (x)=K_{\mu \nu }^\alpha (x)+\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k!}\mathrm{\pounds }_\xi ^{(k)}K_{\mu \nu }^\alpha .$$
(34)
Eqs. (29), (31),(32) and (34) represent the usual transformations on tensorial fields resulting from a general mapping of the manifold on which they are defined. Hence, all tensors in the manifold, in particular the energy-momentum tensor defined above, will transform in the usual homogeneous way:
$$T^{\mu \nu }(x)=T^{\mu \nu }(x)+\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k!}\mathrm{\pounds }_\xi ^{(k)}T^{\mu \nu }(x)$$
(35)
The situation is completely different for the case of Eqs. (28), (30) and (33). Those transformations act only on the fields $`\stackrel{~}{h}^{\mu \nu }`$ and $`K_{\mu \nu }^\alpha `$, letting invariant the background metric $`\gamma ^{\mu \nu }(x)`$. In this sense, one can interpret them as true gauge transformations. It can be shown that the dynamical equations for the gravitational field (15) transform as a combination of themselves under the transformations (30) and (33), supposing that the background space satisfies $`\stackrel{}{R}_{\mu \nu }=2\stackrel{}{\mathrm{\Lambda }}\gamma _{\mu \nu }/(d1)`$. The new field $`\stackrel{~}{h}^{\mu \nu }(x)`$ is also a solution of the field equations, and it corresponds to the same physical field as $`\stackrel{~}{h}^{\mu \nu }(x)`$. In this case, the tensors do not transform in the usual way (35) but contains extra inhomogeneous terms which brings the possibility of annulling them. In this case, the energy-momentum tensor, for example, transforms according to:
$`T_{\mu \nu }(h^{},K^{})`$ $`=`$ $`T_{\mu \nu }(h,K)+{\displaystyle \frac{1}{16\pi G}}\{\widehat{G}_{\mu \nu }^L\left[{\displaystyle \frac{1}{\sqrt{\gamma }}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{k!}}\mathrm{\pounds }_\xi ^{(k)}(\stackrel{~}{\gamma }^{\alpha \beta }+\stackrel{~}{h}^{\alpha \beta })\right]`$ (37)
$`+\gamma _{\alpha \mu }\gamma _{\beta \nu }{\displaystyle \frac{4\stackrel{}{\mathrm{\Lambda }}}{(d1)\sqrt{\gamma }}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{k!}}\mathrm{\pounds }_\xi ^{(k)}(\stackrel{~}{\gamma }^{\alpha \beta }+\stackrel{~}{h}^{\alpha \beta })\},`$
where $`\widehat{G}_{\mu \nu }^L`$ is the operator which when acting on $`h^{\alpha \beta }`$ yields expression (16). Hence, it is always possible to find gauge transformations (30) and (33) which makes the energy-momentum tensor defined previously to be null. This is the analogue in the field theoretical approach to what happens with the pseudotensors in GR. Note that the new energy-momentum tensor in Eq. (37) is also covariantly conserved due to the properties of $`G_{\mu \nu }^L`$.
We have now completed our generalization of the field theoretical approach for the case when the background metric satisfies Eq. (3), which includes the Ricci flat case when we make $`\stackrel{}{\mathrm{\Lambda }}=0`$. Although the energy-momentum tensor defined above suffers from gauge ambiguities, the total energy or angular momentum of a gravitational field does not. As we will see in the next section, total conserved quantities are given by integration of a $`d`$-form on a $`d`$-dimensional hypersurface, which can be reduced to an integration of a $`(d1)`$-form on the $`(d1)`$-dimensional boundary of this hypersurface at the asymptotic limit. As one would like to preserve the asymptotic structure of the gravitational field, the gauge vectors $`\xi ^\alpha `$ must satisfy boundary conditions at the asymptotic $`(d1)`$-dimensional surfaces. They must tend to AdS Killing vectors at spatial infinity (see Ref. for details), and the true gauge transformations reduce to the identity at the boundaries. Hence, as the total conserved quantities can be calculated by integration on the asymptotic boundary, they are gauge independent.
## III The $`(d1)`$-form which gives total conserved quantities of asymptotically anti-de Sitter spacetimes
Consider the energy-momentum tensor $`T_{\mu \nu }`$ defined in Eqs. (19) and (21), and a Killing form $`\xi _\nu `$ of the background metric. Construct, as usual, the current
$$J^\mu =T^{\mu \nu }\xi _\nu ,$$
(38)
and define the $`d`$-form
$$J=\frac{1}{d!}J^\mu \eta _{\mu \alpha _1\mathrm{}\alpha _d}dx^{\alpha _1}\mathrm{}dx^{\alpha _d},$$
(39)
where $`\eta _{\mu \alpha _1\mathrm{}\alpha _d}=\sqrt{\gamma }ϵ_{\mu \alpha _1\mathrm{}\alpha _d}`$ and $`ϵ_{\mu \alpha _1\mathrm{}\alpha _d}`$ is the $`(d+1)`$-dimensional (metric independent) completely antisymmetric object.
Due to energy-momentum conservation together with the Killing equation, it follows that $`J_{;\mu }^\mu =0`$, which is equivalent to $`dJ=0`$, i.e., $`J`$ is a closed $`d`$-form. Hence,
$$_M𝑑J=_MJ=0,$$
(40)
where $`M`$ is a $`(d+1)`$-dimensional spacetime volume. If the boundary $`M`$ of $`M`$ is constituted of two $`d`$-dimensional spacelike hypersurfaces $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_2`$ labelled by the time parameters $`t_1`$ and $`t_2`$, respectively, and a $`d`$ dimensional timelike hypersurface $`B`$ at spatial infinity, and supposing that the fields and Killing vectors are such that $`J`$ is zero at $`B`$ (which is the case for the non-radiating and asymptotically AdS gravitational fields we will study in this paper), then Eq. (40) reduces to
$$_{\mathrm{\Sigma }_2}J_{\mathrm{\Sigma }_1}J=0,$$
(41)
and $`_\mathrm{\Sigma }J`$ is a conserved quantity. If the Killing vector field is timelike or associated with some rotational symmetry, then we will have a conserved total energy or a conserved total angular momentum, respectively. Looking at the equations of motion (22) and the definition of $`G_{\mu \nu }^L`$ in Eq. (16), we can write the conserved current $`J^\mu `$ on shell as:
$$J^\mu =\frac{1}{16\pi G}[(\gamma ^{\mu \nu }h^{\alpha \beta }+\gamma ^{\alpha \beta }h^{\mu \nu }\gamma ^{\alpha \mu }h^{\beta \nu }\gamma ^{\alpha \nu }h^{\beta \mu })_{;\alpha ;\beta }\xi _\nu +\frac{4\stackrel{}{\mathrm{\Lambda }}}{d1}h^{\mu \nu }\xi _\nu ].$$
(42)
Using the fact that the background metric satisfies $`\stackrel{}{R}_{\mu \nu }=2\gamma _{\mu \nu }\stackrel{}{\mathrm{\Lambda }}/(d1)`$, the last term in the right-hand-side of the above equation can be written in a suitable form:
$$\frac{4\stackrel{}{\mathrm{\Lambda }}}{d1}h^{\mu \nu }\xi _\nu =2(\xi _\nu h_\lambda ^\nu )^{[;\mu ;\lambda ]},$$
(43)
where $`A^{[;\mu ;\lambda ]}A^{;\mu ;\lambda }A^{;\lambda ;\mu }`$. In this form the current $`J^\mu `$ does not depend anymore on $`\stackrel{}{\mathrm{\Lambda }}`$ and $`d`$.
The following statement summarizes the mathematical achievement of this work: the skew-symmetric tensor
$$\mathrm{\Omega }^{\alpha \mu }=\frac{1}{64\pi G}(Q_{[;\beta }^{\alpha \mu \beta \nu }\xi _{\nu ]}Q^{\alpha \mu \beta \nu }\xi _{\nu _;\beta }),$$
(44)
where
$$Q^{\alpha \mu \beta \nu }\gamma ^{\mu \nu }h^{\alpha \beta }+\gamma ^{\alpha \beta }h^{\mu \nu }\gamma ^{\alpha \nu }h^{\beta \mu }\gamma ^{\beta \mu }h^{\alpha \nu },$$
(45)
satisfies
$$J^\mu =2\mathrm{\Omega }_{;\alpha }^{\mu \alpha }.$$
(46)
The tensor $`Q^{\alpha \mu \beta \nu }`$ has the same symmetries as the Riemman tensor. Defining the $`(d1)`$-form
$$\mathrm{\Omega }=\frac{1}{(d1)!}\mathrm{\Omega }^{\mu \nu }\eta _{\mu \nu \alpha _1\mathrm{}\alpha _{d1}}dx^{\alpha _1}\mathrm{}dx^{\alpha _{d1}},$$
(47)
then Eq. (46) is equivalent to
$$J=d\mathrm{\Omega },$$
(48)
which proves that $`J`$ is not only closed, but globally exact. Hence, the conserved quantities will be given by
$$Q=_\mathrm{\Sigma }\mathrm{\Omega }.$$
(49)
This total conserved quantity is gauge independent because it was reduced to an integral at the boundary at infinity. There, the true gauge tranformations (30) reduce to the identity because we are imposing the preservation of the asymptotic structure of the fields. Also, with the $`(d1)`$-form $`\mathrm{\Omega }`$ we do not need the knowledge of the gravitational field on the whole $`\mathrm{\Sigma }`$ but only its asymptotic behaviour on $`\mathrm{\Sigma }`$, which makes the calculation of the total conserved quantities much easier and general<sup>*</sup><sup>*</sup>*In Ref. another energy-momentum tensor is defined for flat background, and Eq. (22) is written in another form, with the presence of extra non linear terms in its left-hand-side. However, as total conserved quantities are calculated by means of asymptotic boundary integrals, only the linear terms given by $`\mathrm{\Omega }`$ in Eq. (49) are important..
Note that $`\mathrm{\Omega }`$ does not depend neither on $`d`$ nor on $`\mathrm{\Lambda }`$. All dependence on these parameters are contained in $`h^{\mu \nu }`$. However, we have to fix the background metric in some way because once a metric $`g_{\mu \nu }`$ is given, the determination of $`h^{\mu \nu }`$ will depend on the choice of $`\gamma _{\mu \nu }`$, as can be seen from Eq. (1). Note that even for asymptotically AdS spacetimes one could take Ricci flat background geometries just by making $`\stackrel{}{\mathrm{\Lambda }}=0`$ in Eq. (3): our results are independent on the choice of $`\stackrel{}{\mathrm{\Lambda }}`$. Let us then calculate the total energy of a simple asymptotically AdS spacetime, the Schwarzschild AdS solution in four dimensions, using a Ricci flat background, namely, a flat background. The Schwarzschild AdS solution is:
$$ds^2=\left[1+\left(\frac{r}{R}\right)^2\frac{2m}{r}\right]dt^2+\left[1+\left(\frac{r}{R}\right)^2\frac{2m}{r}\right]^1dr^2+r^2[d\theta ^2+\mathrm{sin}^2(\theta )d\varphi ^2]$$
(50)
where $`R`$ is the radius of curvature of such space, related to the effective cosmological constant $`\mathrm{\Lambda }_{eff}=\mathrm{\Lambda }+8\pi G\rho _V`$ by $`R=(3/\mathrm{\Lambda }_{eff})^{1/2}`$ ($`\rho _V`$ is the matter vacuum energy density). The flat background metric is taken in spherical coordinates. The non null $`h^{\mu \nu }`$ are
$$h^{tt}=\frac{r\rho ^22m}{r2m+r\rho ^2},$$
(51)
and
$$h^{rr}=\frac{2m}{r}+\rho ^2,$$
(52)
where $`\rho r/R`$. To calculate the total energy, we will take the timelike Killing vector field of flat spacetime, $`\xi ^\mu =\delta _t^\mu `$. When we calculate the total energy using Eqs. (44), (45), (47) and (49), we obtain (from now on we will make $`G=1`$):
$$E=\underset{r\mathrm{}}{lim}[\frac{r^2(m+r\rho ^2)}{2(r2m+r\rho ^2)^2}+\frac{m}{2}\frac{r^3}{R^2}]=\frac{m}{2}\underset{r\mathrm{}}{lim}\frac{r^3}{R^2}.$$
(53)
The last term is the energy of the anti-de Sitter solution, which is the Schwarzschild AdS solution with $`m=0`$. Subtracting it we find:
$$E_{ren}=\frac{m}{2}.$$
(54)
This result does not give the Schwarzschild mass, first calculated in the field theoretical approach in Ref. , when we put $`\mathrm{\Lambda }_{eff}`$ to zero after the calculation. For the usual asymptotically flat Schwarzschild solution, the energy in Eq. (53), setting $`\mathrm{\Lambda }_{eff}=\rho =0`$ from the beginning, reads $`E=m`$. Hence, we have an internal inconsistency: the total energy of the Schwarzschild AdS solution using a flat background does not yield the total energy of the asymptotically flat Schwarzschild solution (calculated within the same rules) when we put $`\mathrm{\Lambda }_{eff}`$ to zero. However, as we will see in the following, if we take an ADS background (with the same $`\mathrm{\Lambda }_{eff}`$) in the calculation of the energy of the Schwarzschild AdS solution, the energy in the limit $`\mathrm{\Lambda }_{eff}=0`$ yields the energy of the asymptotically flat Schwarzschild solution calculated in Ref. . This shows that the background metric is not arbitrary but is dictated by the asymptotic structure of the spacetime in question. If we take other simple examples, we can readily conclude that the background metric for asymptotic AdS spacetimes must not only satisfy $`\stackrel{}{R}_{\mu \nu }=2\stackrel{}{\mathrm{\Lambda }}\gamma _{\mu \nu }/(d1)`$, with $`\stackrel{}{\mathrm{\Lambda }}=\mathrm{\Lambda }_{eff}=\mathrm{\Lambda }+8\pi G\rho _V`$, but it must also be asymptotically AdS. A metric which is asymptotically AdS, satisfies Eq. (3), and is regular everywhere must be the AdS metric (an analogous reasoning can be used for asymptotically flat spacetimes). Hence, the asymptotic structure, together with regularity assumptions, fixes the background metric. Also, the background AdS metric must be in the same coordinate system as the AdS asymptotic geometry at infinity, but this coordinate system may be arbitrary because $`\mathrm{\Omega }^{\alpha \mu }`$ in Eq. (44) is a true tensor. This is equivalent to the imposition that only true gauge transformations which reduce to the identity at the spatial infinity are allowed. With these restrictions dictated by the asymptotic structure of the geometry under study, the conserved quantity (49) has no ambiguities: it is invariant under the allowed true gauge transformations, and the background metric is fixed by the asymptotic behaviour of the geometry.
Note that the condition $`\stackrel{}{\mathrm{\Lambda }}=\mathrm{\Lambda }_{eff}=\mathrm{\Lambda }+8\pi G\rho _V`$ is also the constraint one must impose on $`\stackrel{}{\mathrm{\Lambda }}`$ in order to have $`h_{\mu \nu }=0`$ as a solution of the vacuum field equations (9).
With these rules in mind, let us calculate some conserved quantities for asymptotically AdS solutions. In the examples below we will take for convenience, and without loss of generality, that $`\rho _V=0`$ and hence $`\stackrel{}{\mathrm{\Lambda }}=\mathrm{\Lambda }`$. Also, all the quantities calculated below are conserved because in all cases $`_BJ=0`$, where $`B`$ is the timelike portion of $`M`$ in Eq. (40), from where Eq. (41) follows.
1. The Kerr-anti-de Sitter spacetime.
The coordinates of the effective geometry $`g_{\mu \nu }`$ are such that it tends asymptotically to the anti-de Sitter metric in the form
$$d\stackrel{}{s}^2=[1+\left(\frac{r}{R}\right)^2]dt^2+[1+\left(\frac{r}{R}\right)^2]^1dr^2+r^2[d\theta ^2+\mathrm{sin}^2(\theta )d\varphi ^2]$$
(55)
where, as before, $`R=(3/\mathrm{\Lambda })^{1/2}`$. In this case, the leading order aymptotic components of the gravitational field $`h^{\mu \nu }`$ read
$$\{\begin{array}{ccc}\hfill h^{tt}& =& \frac{2mR^4}{r^5}[1\alpha ^2\mathrm{sin}^2\theta ]^{5/2}\hfill \\ \hfill h^{t\varphi }& =& \frac{2amR^2}{r^5}[1\alpha ^2\mathrm{sin}^2\theta ]^{5/2}\hfill \\ \hfill h^{\varphi \varphi }& =& \frac{2ma^2}{r^5}[1\alpha ^2\mathrm{sin}^2\theta ]^{5/2}\hfill \\ \hfill h^{rr}& =& \frac{2m}{r}[1\alpha ^2\mathrm{sin}^2\theta ]^{3/2}\hfill \\ \hfill h^{r\theta }& =& \frac{2ma^2}{r^4}[1\alpha ^2\mathrm{sin}^2\theta ]^{5/2}\mathrm{sin}\theta \mathrm{cos}\theta \hfill \\ \hfill h^{\theta \theta }& =& \frac{2ma^4}{r^7}[1\alpha ^2\mathrm{sin}^2\theta ]^{7/2}\mathrm{sin}^2\theta \mathrm{cos}^2\theta \hfill \end{array}$$
(56)
where $`\alpha =a/R`$, and $`a`$ is related to the angular momentum per unit mass of the Kerr-anti-de Sitter spacetime.
For the calculation of the gravitational energy, we take the timelike Killing vector field of the AdS spacetime $`\xi ^\mu =\delta _t^\mu `$, which is also a Killing vector field of the effective geometry, and insert it in Eq. (44). The only nonzero component of $`\mathrm{\Omega }^{\alpha \mu }`$ is $`\mathrm{\Omega }^{tr}`$, which yields for $`Q`$ in Eq. (49)
$$Q=E=\frac{m}{(1\alpha ^2)^2},$$
(57)
agreeing with the results using pseudotensors , the hamiltonian formalism , and the quasilocal stress tensor . Note that the flat space limit $`\mathrm{\Lambda }=0`$ in the above expression yields the energy of the asymptotically flat Kerr spacetime calculated within the same rules. There are no internal inconsistencies.
The Kerr-anti-de Sitter spacetime also has a conserved angular momentum. To calculate it, we now have to take the Killing vector field $`\xi ^\mu =\delta _\varphi ^\mu `$, which is again a Killing vector field of the effective geometry. We proceed in an analogous way obtaining
$$Q=L_\varphi =\frac{ma}{(1\alpha ^2)^2}$$
(58)
which coincide with calculations on Refs.
2. The Schwarzschild-anti-de Sitter spacetime in five dimensions.
The Schwarzschild-anti-de Sitter metric in five dimensions reads
$`ds^2`$ $`=`$ $`\left[1+\left({\displaystyle \frac{r}{R}}\right)^2\left({\displaystyle \frac{r_0}{r}}\right)^2\right]dt^2+\left[1+\left({\displaystyle \frac{r}{R}}\right)^2\left({\displaystyle \frac{r_0}{r}}\right)^2\right]^1dr^2`$ (60)
$`+r^2\{d\theta _1^2+\mathrm{sin}^2(\theta _1)[d\theta _2^2+\mathrm{sin}^2(\theta _2)d\varphi ^2]\},`$
where $`R=(6/\mathrm{\Lambda })^{1/2}`$, and the background metric is obtained from Eq. (60) by making $`r_0=0`$ in it. The non null $`h^{\mu \nu }`$ are
$$h^{tt}=\frac{l^2}{(1+\rho ^2)(1+\rho ^2l^2)},$$
(61)
and
$$h^{rr}=l^2,$$
(62)
where $`lr_0/r`$ and, as before, $`\rho r/R`$. Taking again the timelike Killing vector field of the $`AdS`$ spacetime $`\xi ^\mu =\delta _t^\mu `$, and inserting it in Eq. (44), we obtain for the total the gravitational energy of this spacetime the value
$$Q=E=\frac{3\pi r_0^2}{8}$$
(63)
which agree with the result of Ref. . This is the only non null conserved quantity of this geometry.
3. The near-horizon limit of the D3-brane
The effective five-dimensional metric now reads
$$ds^2=\left(\frac{r}{R}\right)^2\left\{\left[1\left(\frac{r_0}{r}\right)^4\right]dt^2+(dx^i)^2\right\}+\left[1\left(\frac{r_0}{r}\right)^4\right]^1\left(\frac{r}{R}\right)^2dr^2.$$
(64)
The background metric is obtained by making $`r_0=0`$ in Eq. (64). It is now written in a different coordinate system than in the precedent example but in accordance with the asymototic limit of (64). The non null $`h^{\mu \nu }`$ are
$$h^{tt}=\frac{l^4}{\rho ^2(1l^4)},$$
(65)
and
$$h^{rr}=l^4\rho ^2,$$
(66)
where $`l`$ and $`\rho `$ are defined as above. Taking again the timelike Killing vector field of the AdS spacetime, the gravitational energy now reads
$$Q=E=\frac{3r_0^4}{16\pi R^5}d^3x,$$
(67)
which agree with the result of Ref. . This is the only non null conserved quantity of this geometry.
## IV Conclusion
In this paper we have extended the field theoretical approach to GR to the case where the background metric satisfies $`\stackrel{}{R}_{\mu \nu }=2\gamma _{\mu \nu }\stackrel{}{\mathrm{\Lambda }}/(d1)`$. After that, we have obtained a $`(d1)`$-form $`\mathrm{\Omega }`$ which, when integrated on asymptotic $`(d1)`$-dimensional surfaces, yields the values of total energies and angular momenta of asymptotically $`(d+1)`$-dimensional AdS or flat spacetimes. Although the dynamics of the effective geometry does not depend on the background metric we choose, the values of those total conserved quantities are strongly affected by the choice we make. As we have shown in the text, if we do not choose judiciously the background metric we may obtain preposterous results for the gravitational energy. Hence, the ambiguity in the choice of the background metric may be eliminated only by going beyond the equations of motion and examining further concepts, like the consistency of calculations of total energy and angular momenta. These considerations indicate, together with regularity conditions, what is the background metric one should adopt.
The calculations of total conserved quantities using $`\mathrm{\Omega }`$ yield finite results, and are gauge independent once one does not violate asymptotic conditions. This is not true, however, for conserved quantities contained in finite regions of the background space. These calculations may suffer from gauge ambiguities because the effective geometry may be presented in many different coordinate systems with the same asymptotic limits, and hence, in finite regions, the true gauge transformations (30) are not trivial. It should be interesting to investigate under what subgroup of the true gauge transformations (30) is the $`(d1)`$-form $`\mathrm{\Omega }`$ given in Eqs. (44) and (47) invariant.
Let us now compare the $`(d1)`$-form $`\mathrm{\Omega }`$ given in Eqs. (44,47) and its integral with the quasilocal stress tensor of Ref. , and the surface integrals of Ref. . They have in common the presence of a background (reference) space which is fixed by the asymptotic behaviour of the effective geometry. However, for the $`(d1)`$-form $`\mathrm{\Omega }`$, the presence of the background space is much more important. Contrary to the other prescriptions, the surfaces where the integrals are performed are defined on the background space, and the Killing vector fields which are present in $`\mathrm{\Omega }`$ generates isometries of the background, not of the effective geometry. This last fact does not mean that we can have more conserved quantities than the number of Killing vectors of the effective geometry. If we use some Killing vector field of the background geometry which does not describe an isometry of the effective geometry than the integral $`_BJ`$, where $`B`$ is the timelike portion of $`M`$ in Eq. (40), is not zero and Eq. (41) does not follow. The quantity $`_\mathrm{\Sigma }J`$ is not conserved because there is a flux of $`J`$ through $`B`$.
The $`(d1)`$-form $`\mathrm{\Omega }`$ can be used to calculate conserved quantities or fluxes through the boundary $`B`$ for more involved effective geometries. The fact that $`\mathrm{\Omega }`$ is derived from a theory which describes the gravitational field as a spin-two field propagating on a fixed background may be useful to understand some aspects of the correspondence of conformal field theory in an AdS boundary and gravitational theory in AdS spaces.
ACKNOWLEDGEMENTS
We would like to thank the Cosmology Group of CBPF for useful discussions, and CNPq of Brazil for financial support.
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# Contents
## 1 Introduction
Exponentiated effective action (partition function, statistical sum) is defined as a functional of the coupling constants $`T`$ and background fields $`\phi `$ (the “vacuum configuration”), resulting from functional integration over quantum fields $`\varphi `$:
$`𝒵\{T|\phi \}={\displaystyle \mathrm{exp}\left(S_T(\phi +\varphi )\right)𝒟\varphi }`$ (1.1)
When all possible coupling constants $`T`$ are taken into account (i.e. the theory is maximally deformed), $`𝒵\{T\}`$ becomes a generating function of all the correlation functions in entire family of models. Such $`𝒵\{T\}`$ possesses a hidden group-theoretical structure and – as a manifestation of this – satisfies bilinear (Hirota-like) and differential (Laplace-like) equations, i.e. belongs to the class of generalized $`\tau `$-functions. One can consider $`𝒵\{T\}`$ as a function (section) on the moduli space $``$ of theories (parametrized by the coupling constants $`T`$). There are two important groups, acting transitively on $``$: the abelian group $`\mathrm{𝑆ℎ𝑖𝑓𝑡}`$ of shifts along $``$ and non-abelian group $`\mathrm{𝐷𝑖𝑓𝑓}`$ of diffeomorphisms of $``$. They act on partition functions in the same way:
$`𝒵\{T\}𝒵\{T+V(T)\},`$ (1.2)
but the composition rules are different:
$`𝒵\{T\}𝒵\{T+V_1(T)+V_2(T)\}`$ (1.3)
for $`\mathrm{𝑆ℎ𝑖𝑓𝑡}`$ and
$`𝒵\{T\}𝒵\{T+V_1(T)+V_2(T+V_1(T))\}`$ (1.4)
for $`\mathrm{𝐷𝑖𝑓𝑓}`$. In other words, the infinitesimal action of both groups is decribed by vector fields $`\widehat{V}\{T\}=V(T)/T`$, but the global action is by exponentiated vector field, $`\mathrm{exp}\widehat{V}`$ for $`\mathrm{𝐷𝑖𝑓𝑓}`$ and by the normal-ordered exponent $`:\mathrm{exp}\widehat{V}:`$ for abelian $`\mathrm{𝑆ℎ𝑖𝑓𝑡}`$. A map between $`\mathrm{𝐷𝑖𝑓𝑓}`$ and $`\mathrm{𝑆ℎ𝑖𝑓𝑡}`$ is provided by the relation
$`e^{\widehat{V}}=:e^{\widehat{\stackrel{~}{V}}}:`$ (1.5)
Looking from this perspective, one associates with every particular theory a group element $`g_T`$ with two basic properties
$`g_{T_1}=g_{T_{12}}g_{T_2},T_1=T_{12}T_2`$ (1.6)
and
$`\mathrm{\Delta }g_T=g_Tg_T,`$ (1.7)
and represents $`𝒵\{T\}`$ as a matrix element (generalized zonal-function or $`\tau `$-function): $`𝒵\{T\}=1|g_T|0`$ between a Gaussian theory, labeled by $`|0`$, and some other state $`1|`$, depending on particular realization of $`g_T`$. The transitive action of the group basically puts all the points in the moduli space on equal footing (in particular the Gaussian point is not distinguished among the others), and this can explain the surprising power of the free-field formalism in quantum field theory. The composition rule $`T_1=T_{12}T_2`$ depends on which of the two groups, $`\mathrm{𝑆ℎ𝑖𝑓𝑡}`$ or $`\mathrm{𝐷𝑖𝑓𝑓}`$, we want $`g_T`$ to belong to. In the case of abelian $`\mathrm{𝑆ℎ𝑖𝑓𝑡}`$ it is just an addition: $`T_1=T_{12}+T_2`$ (if the space of coupling constants is big enough and appropriate choice of coordinates in the moduli space is made). A more interesting non-abelian diffeomorphism group, and especially the stability subgroup $`\mathrm{𝐷𝑖𝑓𝑓}_{\mathrm{}}`$ of the Gaussian model in $``$, is relevant for description of one-parametric renormalization-group flows along $``$.
This old set of ideas is supported by evidence in matrix models , Seiberg-Witten theory and AdS/CFT correspondence <sup>1</sup><sup>1</sup>1 The AdS/CFT-correspondence claims that certain Yang-Mills partition functions are represented by the boundary dependence of the bulk actions of certain classical gravities. Among other things, this implies that they satisfy bilinear Hamilton-Jacobi equations, which should be nothing but an avatar of bilinear Hirota and Laplace-like equations for the effective actions. . The purpose of this paper is to claim that additional evidence is provided by the recent studies of A.Connes and D.Kreimer , who actually define the action of the operators $`g_T`$ in the space of graphs (Feynman diagrams). In what follows we basically repeat their reasoning, making use of convenient quantum models and separating the algebraic structures, relevant for Hirota-like equations and Bogolubov’s $`R`$-operation, from peculiarities of particular models and prejudices of conventional local field theory. Some graph-theory and combinatorial routine is omitted.
## 2 The Basics of Bilinear Relations
The basic model for generic studies of integrable structures in quantum field theory is the one of arbitrarily many scalars $`\varphi ^i`$ in $`0+0`$ dimensions:
$`𝒵\{T\}={\displaystyle e^{V_T(\varphi )}\underset{i}{}d\varphi ^i},`$ (2.1)
$`V_T(\varphi )={\displaystyle \underset{n}{}}{\displaystyle \frac{1}{n!}}\left({\displaystyle \underset{i_1,\mathrm{},i_n}{}}T_{i_1\mathrm{}i_n}^{(n)}\varphi ^{i_1}\mathrm{}\varphi ^{i_n}\right)`$ (2.2)
Partition function $`𝒵\{T\}`$ can be represented as a sum over all possible graphs without external legs (vacuum Feynman diagrams). One can exclude disconnected graphs by switching to $`\mathrm{log}𝒵\{T\}`$ – in exchange one gets additional $`1/n`$ coefficients. These can be eliminated by consideration of correlation functions: graphs with external legs. This is the usual routine of diagram technique .
At this point it can make sense to comment on the choice of the model (2.1). Among its particular reductions (truncations) are: the single-scalar model,
$`𝒵_{U(1)}\{T\}={\displaystyle \mathrm{exp}\left(\underset{n=0}{\overset{\mathrm{}}{}}\frac{T_n}{n!}\varphi ^n\right)𝑑\varphi };`$ (2.3)
the $`N\times N`$ matrix model with $`N^2`$ scalars, assembled into a matrix $`\varphi _{ab}`$,
$`𝒵_{U(N)}\{T\}={\displaystyle \mathrm{exp}\left(\underset{n=0}{\overset{\mathrm{}}{}}\frac{T_n}{n!}Tr\varphi ^n\right)\underset{a,b=1}{\overset{N}{}}d\varphi _{ab}};`$ (2.4)
the scalar field in $`d`$ space-time dimensions, where indices $`i`$ become continuous and $`T_{ij}^{(2)}`$ is taken to be Laplace operator; etc. (In the last example, if the space-time is non-compact, it is unavoidable to introduce the background fields $`\phi `$, like it is done in (1.1), to label the boundary conditions and/or the asymptotics at infinities.) The most essential difference of all these popular models from the universal one in (2.1) is that the vertices in (2.1) are of the most general form, e.g. in $`d`$-dimensional theory the $`\varphi ^n`$ coupling should allow any dependence of the coupling “constants” on all the $`n`$ $`d`$-momenta (while in local field models they are indeed constants or at best polynomia in momenta). In such a large moduli space one can distinguish between any two Feynman diagrams, looking at their $`T`$-dependencies (while for the model (2.3) all the diagrams with the same number of propagators, $`l`$, and vertices of valences $`k`$, $`v_k`$, give rise to the same expression, $`T_2^l_kT_k^{v_k}`$; and switching to the matrix model (2.4) introduces nothing more than extra factor $`N^\chi `$, depending on Euler characteristics $`\chi `$ of the corresponding fat graph, which is still not enough to distinguish between any two diagrams). Last, but not the least, such a large moduli space is preserved by renormalization group flow: effective actions at any stage of the flow remains in the class (2.1); also renormalizability is not a restriction on the form of the theory (once it is somehow regularized).
A further extension of (2.1), playing the same role of the universal model for fat graphs as the model (2.1) plays for the ordinary ones, is provided by the matrix model with 2-index fields $`\varphi ^{ij}`$ and the action
$`V_T(\varphi )={\displaystyle \underset{n}{}}{\displaystyle \frac{1}{n!}}\left({\displaystyle \underset{i_1,\mathrm{},i_n}{}}T_{i_1\mathrm{}i_n}^{(n)}\varphi ^{i_1i_2}\varphi ^{i_2i_3}\mathrm{}\varphi ^{i_ni_1}\right)`$ (2.5)
Its principal difference from (2.1) is that the couplings $`T_{i_1\mathrm{}i_n}^{(n)}`$ are no longer symmetric under permutations of indices $`i_1,\mathrm{},i_n`$, only cyclic symmetric. This model remains beyond our consideration in the present paper.
The model (2.1) can be also regarded in a different way. Given any particular quantum theory one can switch to its $`GL(N)`$ or $`GL(\mathrm{})`$ extension, just adding a vector index $`iI`$ to all the fields of the theory and ascribing the relevant tensor structure to all the coupling constants. For example, the $`d`$-dimensional $`\varphi ^3`$ theory can be substituted by
$`{\displaystyle \underset{i}{}D\varphi ^i(x)\mathrm{exp}_{d^dx}\left(T_{ij}^{(2)}\left((\varphi ^i)(\varphi ^j)m^2\varphi ^i\varphi ^j\right)g_3T_{ijk}^{(3)}\varphi ^i\varphi ^j\varphi ^k\right)}`$ (2.6)
without requiring that $`T`$-variables are $`x`$-dependent. Then the partition function is $`GL(\mathrm{})`$-invariant and can be expanded in a series over the basic $`GL(\mathrm{})`$-invariant functions, provided by (2.1), however the expansion coefficients are now sophisticated functions not only of graphs, but also of many other parameters, including external momenta and the spins of particles. Still, some basic properties can be seen at the level of graph theory alone – and this is the subject of our futher considerations.
We now return to the main line of discussion. The issue of our interest is the group (integrable) structure, hidden in partition functions $`𝒵\{T\}`$. This structure survives various reductions, e.g. the one to the matrix model (2.4), see , but many aspects are much more transparent in analysis of the universal model (2.1). Of course, the relation
$`V_{T+T^{}}=V_T+V_T^{}`$ (2.7)
for the potential under the integral does not imply that the average $`𝒵\{T\}`$ is a character of $`\mathrm{𝑆ℎ𝑖𝑓𝑡}`$ group, $`𝒵\{T+T^{}\}𝒵\{T\}𝒵\{T^{}\}`$, the true relation is between the group elements, $`g_{T+T^{}}=g_Tg_T^{}`$, and
$`𝒵\{T+T^{}\}=1|g_{T+T^{}}|0={\displaystyle \underset{states}{}}1|g_T|statesstates|g_T^{}|0`$ (2.8)
At the r.h.s. stands a non-trivial operator, acting on $`T`$ and $`T^{}`$ in the product $`𝒵\{T\}𝒵\{T^{}\}`$.
The question is, what are the relevant realizations of the Hilbert space of $`|states`$ and of the operators $`g_T`$ acting in it.
The simplest realization is implied by the functional integral and makes special use of the source-dependence of $`𝒵\{T\}`$. Namely, the $`T^{(1)}`$-terms in (2.2) can be identified with the sources for the fields $`\varphi ^i`$: $`_jT_j^{(1)}\varphi ^j=i_jJ_j\varphi ^j`$ and, as usual in the derivation of Hirota equations, one can use sources to construct a delta-function projector:
$`{\displaystyle e^{i(\varphi ^j\varphi ^j)J_j^{}}\underset{j}{}dJ_j^{}}{\displaystyle \underset{j}{}}\delta (\varphi ^j\varphi ^j),`$ (2.9)
so that
$`{\displaystyle 𝒵_{JJ^{}}\{T\}𝒵_J^{}\{T^{}\}𝑑J^{}}=`$
$`={\displaystyle 𝑑J^{}\left(𝑑\varphi e^{V_T(\varphi )}e^{i(JJ^{})\varphi }𝑑\varphi ^{}e^{V_T^{}(\varphi )}e^{i(JJ^{})\varphi ^{}}\right)}=𝒵_J\{T+T^{}\}`$ (2.10)
(we explicitly labeled the $`J`$\- ($`T^{(1)}`$-) dependence of the action, suppressed all the indices $`i`$ and made use of the addition formula (2.7)).
The simplest example arises if all $`T^{(n)}=0`$ for $`n>2`$. Then
$`𝒵_{Gauss}\{T\}{\displaystyle \frac{1}{\sqrt{detT^{(2)}}}}\mathrm{exp}\left({\displaystyle \frac{1}{4}}TrT^{(1)}{\displaystyle \frac{1}{T^{(2)}}}T^{(1)}\right),`$ (2.11)
and the bilinear relation
$`𝒵_{Gauss}\{T+T^{}\}={\displaystyle 𝑑T^{(1)}𝑑T^{(1)}\delta (T^{(1)}T^{(1)})𝒵_{Gauss}\{T\}𝒵_{Gauss}\{T^{}\}}`$ (2.12)
is just the completeness formula for Guassian propagators, which in the single-scalar case is widely-known in the form
$`{\displaystyle 𝑑x_2\frac{e^{x_{12}^2/4t_{12}}}{\sqrt{t_{12}}}\frac{e^{x_{23}^2/4t_{23}}}{\sqrt{t_{23}}}}{\displaystyle \frac{e^{x_{13}^2/4t_{13}}}{\sqrt{t_{13}}}}`$ (2.13)
(in our case $`x_{12}=x_1x_2=T^{(1)}`$, $`x_{23}=T^{(1)}`$, $`t_{12}=T^{(2)}`$, $`t_{23}=T^{(2)}`$).
In a similar way one can obtain bilinear integral “summation formulae” for the Eiry functions (if all $`T^{(n)}=0`$ for $`n>3`$) etc.
Instead of using the source-dependencies, one can exploit those on other coupling constants. All these dependencies are interrelated: in the universal model (2.1)
$`{\displaystyle \frac{𝒵}{T_{i_1\mathrm{}i_n}^{(n)}}}={\displaystyle \frac{^n𝒵}{T_{i_1}^{(1)}\mathrm{}T_{i_n}^{(1)}}},`$ (2.14)
in its various reductions one has more sophisticated Ward identities, like Virasoro and W-constraints in matrix models . Also, Legendre transform relates source- and background-field dependencies of generic $`𝒵\{T|\phi \}`$ in (1.1). For more discussion of interplay between the source-, coupling-constants and background-fields dependencies see , especially the example of generalized Kontsevich model . The Ward identities like (2.14) play important role in building explicit maps like (1.5) between the $`\mathrm{𝑆ℎ𝑖𝑓𝑡}`$ and $`\mathrm{𝐷𝑖𝑓𝑓}`$ groups.
Drawback of such functional-integral approaches to bilinear identities is that they do not immediately provide representations in terms of conventional (perturbative) correlation functions: at least some operator, like the source term $`_jJ_i\varphi ^i`$, should be exponentiated, i.e. one needs to consider a global deformation and the entire family of theories, not just infinitesimal vicinities of the given models $`g_T`$ and $`g_T^{}`$. Though there is nothing bad about this from the general perspective of string theory, such representations are not the best ones for the search of bilinear relations in conventional quantum field theories, where isolated points in the moduli spaces (i.e. isolated particular models) are usually analyzed. One possibility to obtain representations in terms of the ordinary Green functions is to use the Vermat-module realizations of $`|states`$ . Another – not unrelated – possibility is exploited by A.Connes and D.Kreimer (CK) in : it is to look at the contributions of particular graphs (Feynman diagrams).
## 3 Hilbert Space of Graphs and Operators $`g_T`$
Operators, acting in the Hilbert space of Feynman diagrams naturally appear if the Gaussian measure is explicitly extracted from $`\mathrm{exp}V_T`$ (as a starting point for perturbation expansion) and if the theory (2.1) is further “complexified”:
$`𝒵\{\stackrel{~}{T},T\}={\displaystyle \left\{\underset{i,\stackrel{~}{i}}{}d\varphi ^{i\stackrel{~}{i}}\mathrm{exp}\left(\frac{1}{2}\underset{i,j,\stackrel{~}{i},\stackrel{~}{j}}{}G_{ij}\stackrel{~}{G}_{\stackrel{~}{i}\stackrel{~}{j}}\varphi ^{i\stackrel{~}{i}}\varphi ^{j\stackrel{~}{j}}\right)\right\}e^{V_{\stackrel{~}{T},T}(\varphi )}},`$ (3.1)
$`V_{\stackrel{~}{T},T}(\varphi )={\displaystyle \underset{n}{}}{\displaystyle \frac{1}{n!}}\left({\displaystyle \underset{i_1,\mathrm{},i_n;\stackrel{~}{i}_1,\mathrm{},\stackrel{~}{i}_n}{}}\stackrel{~}{T}_{\stackrel{~}{i}_1,\mathrm{},\stackrel{~}{i}_n}^{(n)}T_{i_1\mathrm{}i_n}^{(n)}\varphi ^{\stackrel{~}{i}_1i_1}\mathrm{}\varphi ^{\stackrel{~}{i}_ni_n}\right)`$ (3.2)
The fields $`\varphi ^{i\stackrel{~}{i}}`$ are now labeled by a pair of indices, taking values in the sets $`I`$ and $`\stackrel{~}{I}`$, $`iI`$, $`\stackrel{~}{i}\stackrel{~}{I}`$, and coupling constants are factorized (assumed to be the “squared modules” of “holomorphic” $`T`$’s). In variance with ordinary complexification, we do not assume that the sets $`I`$ and $`\stackrel{~}{I}`$ are the same. In particular, we can return to the original model (2.1) by asking $`\stackrel{~}{I}`$ to consist of a single element and putting all $`\stackrel{~}{T}_{1\mathrm{}1}^{(n)}=1`$ and $`\stackrel{~}{G}_{11}=1`$ (below,when necessary, we just write $`\stackrel{~}{T}=1`$, implying that $`\stackrel{~}{I}=\{1\}`$, and $`𝒵\{T\}=𝒵\{\stackrel{~}{T}=1,T\}`$).
Expanding $`\mathrm{exp}V_{\stackrel{~}{T},T}(\varphi )`$ in (3.1) into formal series and applying the Wick theorem to Gaussian integrals, one obtains the expansion over vacuum Feynman diagrams $`\mathrm{\Gamma }^{(0)}`$, which has specific structure, called “holomorphic factorization”:
$`𝒵\{\stackrel{~}{T},T\}={\displaystyle \underset{\mathrm{\Gamma }^{(0)}}{}}{\displaystyle \frac{𝒵_\mathrm{\Gamma }\{\stackrel{~}{T}\}𝒵_\mathrm{\Gamma }\{T\}}{S_\mathrm{\Gamma }}}`$ (3.3)
Here $`S_\mathrm{\Gamma }`$ is the “symmetry factor” of the graph $`\mathrm{\Gamma }`$, for connected graph it is the order of the discrete group which permutes links, while keeping their ends fixed. For vacuum diagrams (graphs without external legs) $`S_\mathrm{\Gamma }`$ contains an additional factor $`\mathrm{𝑉𝑒𝑟𝑡}(\mathrm{\Gamma })`$ – the number of vertices in the graph . For disconnected graphs $`S(_i\mathrm{\Gamma }_i^{n_i})=_in_i!S_{\mathrm{\Gamma }_i}^{n_i}`$ Expression $`Z_\mathrm{\Gamma }\{T\}/S_\mathrm{\Gamma }`$ for the Feynman diagram $`\mathrm{\Gamma }`$ is a convolution of vertices and propagators, divided by $`S_\mathrm{\Gamma }`$. In particular, $`Z_\mathrm{\Gamma }\{T=1\}=1`$, and for the model (2.1) eq.(3.3) gives:
$`𝒵\{T\}=𝒵\{\stackrel{~}{T}=1,T\}={\displaystyle \underset{\mathrm{\Gamma }^{(0)}}{}}{\displaystyle \frac{𝒵_\mathrm{\Gamma }\{T\}}{S_\mathrm{\Gamma }}}`$ (3.4)
In Feynman diagrams $`G_{ij}`$ plays the role of inverse propagator. In what follows we lower and raise indices with the help of $`G_{ij}`$ and its inverse $`G^{ij}`$. In particular the switch from Green functions to the amputated correlators is nothing but lowering of indices on external legs. Coupling constants are defined to have lower indices.
Introduce now the Hilbert space $`^{(0)}`$ of all possible graphs (with vertices of any valence, connected or disconnected) with no external legs (vacuum Feynman diagrams), i.e. with every graph $`\mathrm{\Gamma }^{(0)}`$ we associate a state $`|\mathrm{\Gamma }^{(0)}`$. The scalar product is
$`\mathrm{\Gamma }^{(0)}|\mathrm{\Gamma }^{(0)}=S_\mathrm{\Gamma }\delta _{\mathrm{\Gamma },\mathrm{\Gamma }^{}}`$ (3.5)
Let us further define the coherent-like states
$`|T={\displaystyle \underset{\mathrm{\Gamma }^{(0)}}{}}{\displaystyle \frac{𝒵_\mathrm{\Gamma }\{T\}}{S_\mathrm{\Gamma }}}|\mathrm{\Gamma }`$ (3.6)
They are not orthonormal, instead (3.3) states that
$`𝒵\{\stackrel{~}{T},T\}=\stackrel{~}{T}|T`$ (3.7)
Thus with every particular model (a point in the moduli space $``$) we associate a point $`|T`$ in the Hilbert space of vacuum graphs, and partition function $`𝒵\{\stackrel{~}{T},T\}`$ (3.1) is just a scalar product of associated states. In particular, the holomorphic partition function $`𝒵\{T\}`$ in (2.1) is a scalar product of $`|T`$ with a special state $`\stackrel{~}{T}=1|`$, where $`\stackrel{~}{T}`$ has single-valued indices and all $`\stackrel{~}{T}=1`$.
Another special state is the Gaussian model, $`|Gauss=|0=|\mathrm{}`$, associated with all $`T=0`$ (with any set $`I`$ and non-vanishing metric $`G_{ij}`$). At Gaussian point the only contribution to (3.6) comes from the empty graph $`\mathrm{\Gamma }=\mathrm{}`$ (we assume the normalization $`𝒵\{T=0\}=1`$). There is an even more distinguished “trivial” model with $`𝒵\{T\}=\mathrm{exp}T^{(0)}`$, the transitive actions of $`\mathrm{𝑆ℎ𝑖𝑓𝑡}`$ and $`\mathrm{𝐷𝑖𝑓𝑓}`$ groups connect it to any other model, including the Gaussian one. However, the structures which are of interest for us are explicitly dependent on the metric $`G_{ij}`$, and the minimal model which takes this into account is the Gaussian one: arbitrary diffeomorphisms and shifts can be expanded over basic functions, provided by Gaussian model, but not by the “trivial” one.
Let us now introduce an operator $`g_T`$, acting in the Hilbert space of graphs, such that
$`|T=g_T|Gauss=g_T|\mathrm{}`$ (3.8)
Then
$`𝒵_\mathrm{\Gamma }\{T\}=\mathrm{\Gamma }|T=\mathrm{\Gamma }|g_T|\mathrm{}`$ (3.9)
In sec.8 below we use these matrix elements to convert functions of coupling constants $`T`$ into functions of graphs $`\mathrm{\Gamma }`$ and vice versa.
Eq.(3.9) does not fully specify the operator $`g_T`$. However, field theory implies a natural extension of (3.9) to all the matrix elements
$`\mathrm{\Gamma }_{i_1\mathrm{}i_n}^{(n)}|g_T|\gamma _{j_1\mathrm{}j_m}^{(m)}`$ (3.10)
in the enlarged Hilbert space $`=_n^{(n)}`$ of all the graphs $`\mathrm{\Gamma }^{(n)}`$ with any number $`n`$ of external legs with indices $`iI`$ ascribed to every leg.
The space $`^{(n)}`$ naturally appears if one considers the $`n`$-point correlation functions in the theory (3.1). Such correlator carries extra $`2n`$ indices and is decomposed into a sum over all the graphs with $`n`$ external legs:
$`𝒵\{\stackrel{~}{T},T\}^{\stackrel{~}{i}_1\mathrm{}\stackrel{~}{i}_n;i_1\mathrm{}i_n}={\displaystyle \underset{\mathrm{\Gamma }^{(n)}}{}}{\displaystyle \frac{𝒵_\mathrm{\Gamma }\{\stackrel{~}{T}\}^{\stackrel{~}{i}_1\mathrm{}\stackrel{~}{i}_n}𝒵_\mathrm{\Gamma }\{T\}^{i_1\mathrm{}i_n}}{S_\mathrm{\Gamma }}}`$ (3.11)
Similarly to the case of $`^{(0)}`$ one can now define a set of states $`|\mathrm{\Gamma }_{(n)}^{i_1\mathrm{}i_n}`$ with the scalar product
$`\mathrm{\Gamma }_{(n)}^{i_1\mathrm{}i_n}|\mathrm{\Gamma }_{(m)}^{}{}_{}{}^{j_1\mathrm{}j_m}=S_\mathrm{\Gamma }G^{i_1j_1}\mathrm{}G^{i_nj_n}\delta _{\mathrm{\Gamma },\mathrm{\Gamma }^{}}`$ (3.12)
(Since scalar product is non-vanishing only for coincident graphs, the number of external legs are also the same, and $`G^{i_kj_k}`$ couples the indices ascribed to the same $`k`$-th leg.) One can amputate external legs by lowering the indices with the help of the metric $`G_{ij}`$.<sup>2</sup><sup>2</sup>2 When one cuts a link in a graph, two new external legs are formed at the place of a single propagator $`G^1`$ and one glues them back with the help of the metric $`G`$, or, alternatively, amputate one leg in each pair. This can be done more symmetrically, if $`G=D^2`$: then one can associate with every external leg the matrix $`D^1`$, instead of the usual rule, ascribing the propagator $`G^1`$ to non-amputated leg and unity to the amputated one. In continuous case, when $`G`$ is Laplace operator, $`D`$ turns into a Dirac operator. Though the use of $`D`$ can make the bilinear relations below conceptually more symmetric, we ignore this possibility in the present text.
Now we introduce in $`^{(n)}`$ the state
$`|T^{(n)}={\displaystyle \underset{\mathrm{\Gamma }^{(n)};i_1,\mathrm{},i_n}{}}{\displaystyle \frac{𝒵_\mathrm{\Gamma }\{T\}^{i_1\mathrm{}i_n}}{S_\mathrm{\Gamma }}}|\mathrm{\Gamma }_{i_1\mathrm{}i_n}^{(n)}`$ (3.13)
(the indices are lowered with the help of the metric $`G`$) and finally the state
$`|T=_n|T^{(n)}`$ (3.14)
in entire $``$.
In order to define the matrix elements of $`g_T`$ between any two states in $``$ we need to introduce the notion of subgraph.
## 4 Subgraphs
There are two different notions of subgraph, relevant for our further discussion. Let $`\mathrm{\Gamma }^{(n)}`$ be a graph (connected or disconnected, possibly one-particle reducible) with $`n`$ external legs. It has vertices of any valence (including one and two).
1) Vertex-subgraps. Divide the set of vertices in two non-intersecting subsets and cut all the links, connecting vertices from different sets. If $`m`$ legs were cut, we decompose the original graph $`\mathrm{\Gamma }^{(n)}`$ into two disconnected graphs $`\gamma _1^{m+n_1}`$ and $`\gamma _2^{m+n_2}`$, such that $`n_1+n_2=n`$. We call them vertex-subgraphs of $`\mathrm{\Gamma }^n`$ and introduce a notation $`\gamma _2=\mathrm{\Gamma }/\gamma _1`$ (of course, also $`\gamma _1=\mathrm{\Gamma }/\gamma _2`$). The empty graph $`\gamma ^{(0)}=\mathrm{}`$ and $`\gamma ^{(n)}=\mathrm{\Gamma }^{(n)}`$ are vertex-subgraphs of $`\mathrm{\Gamma }^{(n)}`$. The number of vertices $`Vert(\gamma )+Vert(\mathrm{\Gamma }/\gamma )=Vert(\mathrm{\Gamma })`$.
2) Box-subgraphs. Pick a non-empty<sup>3</sup><sup>3</sup>3 The would-be empty box is not well defined. If there are no vertices inside the box, it still can be not empty: contain fragments of some links. To avoid such ambiguities we exclude empty graphs from the set of box-subgraphs of $`\mathrm{\Gamma }`$. When necessary, their contributions will be explicitly added to sums over the set $`\mathrm{\Gamma }`$ of box-subgraphs. subset of vertices and draw a box or a set of non-intersecting boxes around them. Boxes should not lie one inside another. Each box in the set should contain at least one vertex, and the subgraph inside the box should be connected. The sides of the box cut some links of original graph, in particular a link connecting two vertices from our subset can be cut (and these two vertices can belong to the same box or to two disconnected boxes). We call the subgraph $`\gamma `$ lying in this system of boxes a box-subgraph of $`\mathrm{\Gamma }^{(n)}`$. Its complement is no longer a box-subgraph: it can contain just a remnant of a double-cut link with no vertices. Instead of a complement, for a box-subgraph $`\gamma ^{(m)}`$ one can always define a contraction $`[\mathrm{\Gamma }^{(n)}/\gamma ^{(m)}]`$ obtained when each connected component of a box, which cuts links at $`k`$ places is substituted by a single valence-$`k`$ vertex. The resulting graph $`[\mathrm{\Gamma }^{(n)}/\gamma ^{(m)}]`$ has $`n`$ external legs, as the original $`\mathrm{\Gamma }^{(n)}`$ and the same number of connected components, $`\mathrm{𝐶𝑜𝑛}([\mathrm{\Gamma }/\gamma ])=\mathrm{𝐶𝑜𝑛}(\mathrm{\Gamma })`$. According to this definition, the empty graph $`\mathrm{}`$ is not a box-subgraph of $`\mathrm{\Gamma }`$, and there is no box-subgraph $`\gamma `$, such that $`[\mathrm{\Gamma }/\gamma ]=\mathrm{}`$. The number of vertices $`Vert(\gamma )+Vert([\mathrm{\Gamma }/\gamma ])=Vert(\mathrm{\Gamma })+\mathrm{𝐶𝑜𝑛}(\mathrm{\Gamma })`$.
The same graph $`\gamma `$ can happen to be a vertex-subgraph and a box-subgraph simultaneously, but the two sets $`𝒱\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }`$ (of vertex- and box-subgraphs respectively) are different. $`𝒱\mathrm{\Gamma }`$ is just a set-theory object: the set of all subsets of the set of vertices of $`\mathrm{\Gamma }`$, in particular there are always exactly $`2^{Vert(\mathrm{\Gamma })}`$ vertex-subgraphs. As to $`\mathrm{\Gamma }`$, this is a more sophisticated object, essentially depending on the graph structure (not just the set-theory one), in particular, the size of this set depends on the valences of vertices and on exact construction of the links.
The set $`𝒱\mathrm{\Gamma }`$ is related to the abelian group $`\mathrm{𝑆ℎ𝑖𝑓𝑡}`$ (and is relevant for description of bilinear identities), while $`\mathrm{\Gamma }`$ is related to the non-abelian group $`\mathrm{𝐷𝑖𝑓𝑓}_{\mathrm{}}`$, generated by vector fields on $``$ (and is relevant for description of Bogolubov’s recursion and renormalization flows).
### 4.1 Examples
It is now instructive to consider some examples. We mention three classes of simple graphs, useful for various illustrations.
1) Single-vertex graph $`\mathrm{\Gamma }_{p;n}^{(p2n)}`$ has one valence-$`p`$ vertex, $`n`$ propagators and $`p2n`$ external legs. $`\mathrm{\Gamma }_{p;0}^{(p)}`$ is the elementary (bare) vertex of valence $`p`$.
$`\mathrm{\Gamma }_{p;n}^{(p2n)}`$ has just two vertex-subgraphs: $`\gamma =\mathrm{}`$ and $`\gamma =\mathrm{\Gamma }_{p;n}^{(p2n)}`$ itself. The corresponding complements are $`\mathrm{\Gamma }_{p;n}^{(p2n)}/\mathrm{}=\mathrm{\Gamma }_{p;n}^{(p2n)}`$ and $`\mathrm{\Gamma }_{p;n}^{(p2n)}/\mathrm{\Gamma }_{p;n}^{(p2n)}=\mathrm{}`$.
At the same time, there are $`2^n`$ box-subgraphs (of which $`n+1`$ are topologically different): $`\gamma =\frac{n!}{k!(nk)!}\times \mathrm{\Gamma }_{p;k}^{(p2k)}`$, $`0kn`$ (binomial coefficient $`n!/k!(nk)!`$ denotes the multiplicity of the subgraph). The corresponding contractions $`\left[\mathrm{\Gamma }_{p;n}/\mathrm{\Gamma }_{p;k}\right]`$ are $`\frac{n!}{k!(nk)!}\times \mathrm{\Gamma }_{p2k;nk}^{(p2n)}`$.
2) Two-vertex graph $`\mathrm{\Gamma }_{p,q;n}^{(p+q2n)}`$ has one valence-$`p`$ and one valence-$`q`$ vertices, $`n`$ propagators between them (i.e. $`0np,q`$) and $`p+q2n`$ external legs.
It has $`4`$ vertex-subgraphs, $`\gamma =\mathrm{},\mathrm{\Gamma }_{p;0},\mathrm{\Gamma }_{q;0},\mathrm{\Gamma }_{p,q;n}`$, and $`2^n+2`$ box-subgraphs: $`\gamma =\mathrm{\Gamma }_{p;0},\mathrm{\Gamma }_{q,0},\frac{n!}{k!(nk)!}\mathrm{\Gamma }_{p,q;k},0kn`$. The corresponding $`[\mathrm{\Gamma }/\gamma ]=\mathrm{\Gamma }_{q;0},\mathrm{\Gamma }_{p;0},\frac{n!}{k!(nk)!}\times \mathrm{\Gamma }_{p+q2k;nk}`$.
3) Chain graph $`C_N`$ has $`N`$ valence-two vertices, connected chain-wise by $`N1`$ propagators. It has $`2`$ external legs.
$`C_N`$ has $`2^N`$ vertex-subgraphs and $`\beta _N`$ box-subgraphs.
$`N=1`$.
Vertex subgraphs:
$`\begin{array}{cccc}\gamma & =& \mathrm{},& C_1\\ C_1/\gamma & =& C_1,& \mathrm{}\end{array}`$ (4.3)
Box-subgraphs ($`\beta _1=1`$):
$`\begin{array}{ccc}\gamma & =& C_1\\ [C_1/\gamma ]& =& C_1\end{array}`$ (4.6)
$`N=2`$.
Vertex subgraphs:
$`\begin{array}{ccccc}\gamma & =& \mathrm{},& 2\times C_1,& C_2\\ C_2/\gamma & =& C_2,& 2\times C_1,& \mathrm{}\end{array}`$ (4.9)
Box-subgraphs ($`\beta _2=4`$):
$`\begin{array}{ccccc}\gamma & =& 2\times C_1,& C_2,& C_1C_1\\ [C_2/\gamma ]& =& 2\times C_2,& C_1,& C_2\end{array}`$ (4.12)
$`N=3`$.
Vertex subgraphs:
$`\begin{array}{cccccccc}\gamma & =& \mathrm{},& 2\times C_1,& C_1,& 2\times C_2,& C_1C_1,& C_3\\ C_3/\gamma & =& C_3,& 2\times C_2,& C_1C_1,& 2\times C_1,& C_1,& \mathrm{}\end{array}`$ (4.15)
Box-subgraphs ($`\beta _3=12`$):
$`\begin{array}{cccc}\gamma & =& 3\times C_1,& 2\times C_2,\\ [C_3/\gamma ]& =& 3\times C_3,& 2\times C_2,\end{array}`$ (4.18)
$`\begin{array}{ccccc}C_3,& 2\times (C_1C_1),& C_1C_1,& 2\times (C_1C_2),& C_1C_1C_1\\ C_1,& 2\times C_2,& C_3,& 2\times C_2,& C_3\end{array}`$ (4.21)
Every box-subgraph of $`C_N`$ is located in $`s`$ non-intersecting boxes, with $`k`$-th box beginning at link $`i_k`$ and ending at link $`j_k`$. The total number
$`\beta _N={\displaystyle \underset{s=1}{\overset{N}{}}}\beta (N;s)=`$
$`={\displaystyle \underset{s=1}{\overset{N}{}}}\left({\displaystyle \underset{0i_1<j_1i_2<j_2\mathrm{}i_s<j_sN}{}}1\right)=\mu _{N1}\mu _N`$ (4.22)
The number of chain graphs with exactly $`s`$ connected components is
$`\beta (N;s)={\displaystyle \frac{(N+s)!}{(2s)!(Ns)!}}`$ (4.23)
The $`\mu _N`$ are Fibonacchi-like numbers, satisfying recurrent relations:
$`\mu _{2k}=4\mu _{2k1}\mu _{2k3},\mu _{2k+1}=3\mu _{2k1}\mu _{2k3}`$ (4.24)
and initial conditions $`\mu _0=\mu _1=1`$. Consequently
$`(\mu _0,\mu _1,\mu _2,\mathrm{})=(1,1,4,3,11,8,29,21,76,55,\mathrm{})`$ (4.25)
and
$`(\beta _1,\beta _2,\mathrm{})=(1,4,12,33,88,232,609,1596,\mathrm{})`$ (4.26)
## 5 Vertex-subgraphs, action of $`g_T`$ in $``$ and bilinear relations
We are now ready to define the matrix elements of $`g_T`$. They are different from zero only for $`\gamma `$ which is a vertex-subgraph of $`\mathrm{\Gamma }`$ (consequently, $`g_T`$ is triangular and can not be Hermitean operator – this is natural, since it is an element of a group, not algebra,– moreover, triangularity implies that $`g_T^{}g_T`$). For simplicity we first assume that no external legs of $`\mathrm{\Gamma }^{(n)}`$ were cut to make the subgraph $`\gamma ^{(m)}`$. Then
$`\mathrm{\Gamma }_{(n)}^{i_1\mathrm{}i_n}|g_T|\gamma _{(m)}^{j_1\mathrm{}j_m}=Z_{\mathrm{\Gamma }/\gamma }^{i_1\mathrm{}i_nj_1\mathrm{}j_m}\{T\}`$ (5.1)
In other words, the matrix element is given by the expression for Feynman diagram $`\mathrm{\Gamma }/\gamma `$ in the theory $`g_T`$ without the usual symmetry factor $`1/S_{\mathrm{\Gamma }/\gamma }`$. This means that every link in $`\mathrm{\Gamma }/\gamma `$ carries a propagator $`G^{ij}`$ and every vertex of valence $`k`$ in $`\mathrm{\Gamma }/\gamma `$ contributes $`T_{i_1\mathrm{}i_k}^{(k)}`$. The indices are contracted and summed over. Since $`\mathrm{\Gamma }^{(n)}/\gamma ^{(m)}`$ has $`n`$ original and $`m`$ new-formed external legs, the whole matrix element has $`n+m`$ free indices. If some $`p`$ of external legs of $`\gamma `$ coincide with external legs of $`\mathrm{\Gamma }`$, the corresponding indices appear as $`\delta _j^i`$ (or $`G^{ij}`$) factors.
$`\mathrm{\Gamma }_{(n)}^{i_1\mathrm{}i_{np}l_1\mathrm{}l_p}|g_T|\gamma _{(m)}^{j_1\mathrm{}j_{mp}}{}_{k_1\mathrm{}k_p}{}^{}=Z_{\mathrm{\Gamma }/\gamma }^{i_1\mathrm{}i_{np}j_1\mathrm{}j_{mp}}\{T\}\delta _{k_1}^{l_1}\mathrm{}\delta _{k_p}^{l_p}`$ (5.2)
According to our definition, if $`\mathrm{\Gamma }^{(n)}`$ consists of two disconnected components $`\mathrm{\Gamma }^{(n_1)}`$ and $`\mathrm{\Gamma }^{(n_2)}`$, $`n=n_1+n_2`$, then the same is true about $`\gamma ^{(m)}`$, it also consists of disconnected $`\gamma ^{(m_1)}`$ and $`\gamma ^{(m_2)}`$ (both can still be disconnected), $`m=m_1+m_2`$, and
$`\mathrm{\Gamma }^{(n)}|g_T|\gamma ^{(m)}=\mathrm{\Gamma }^{(n_1)}|g_T|\gamma ^{(m_1)}\mathrm{\Gamma }^{(n_2)}|g_T|\gamma ^{(m_2)}`$ (5.3)
It is natural to introduce the product of disconnected graphs as their unification and then interpret (5.3) as the group-element property (1.7) of $`g_T`$.
From the definition it immediately follows that
$`{\displaystyle \underset{j}{}}\mathrm{\Gamma }_{i_1\mathrm{}i_n}^{(n)}|g_T|\gamma _{(m)}^{j_1\mathrm{}j_m}\gamma _{j_1\mathrm{}j_m}^{(m)}|g_T|\stackrel{~}{\gamma }_{(l)}^{k_1\mathrm{}k_l}=\mathrm{\Gamma }_{i_1\mathrm{}i_n}^{(n)}|g_T|\stackrel{~}{\gamma }_{(l)}^{k_1\mathrm{}k_l}`$ (5.4)
for any fixed triple of vertex-subgraphs $`\stackrel{~}{\gamma }\gamma \mathrm{\Gamma }`$ and given $`g_T`$.
The basic relation (1.6) now acquires the form:
$`{\displaystyle \underset{all\gamma 𝒱\mathrm{\Gamma }:\stackrel{~}{\gamma }\gamma \mathrm{\Gamma }}{}}\mathrm{\Gamma }|g_T|\gamma \gamma |g_T^{}|\stackrel{~}{\gamma }=\mathrm{\Gamma }|g_{T+T^{}}|\stackrel{~}{\gamma }`$ (5.5)
for any fixed $`\stackrel{~}{\gamma }𝒱\mathrm{\Gamma }`$ and any two $`g_T`$ and $`g_T^{}`$. In more detail, the multiplication relation states:
$`{\displaystyle \underset{m}{}}\left({\displaystyle \underset{all\gamma 𝒱\mathrm{\Gamma }:\stackrel{~}{\gamma }\gamma \mathrm{\Gamma }}{}}\left({\displaystyle \underset{j}{}}\mathrm{\Gamma }_{i_1\mathrm{}i_n}^{(n)}|g_T|\gamma _{(m)}^{j_1\mathrm{}j_m}\gamma _{j_1\mathrm{}j_m}^{(m)}|g_T^{}|\stackrel{~}{\gamma }_{(l)}^{k_1\mathrm{}k_l}\right)\right)=`$
$`=\mathrm{\Gamma }_{i_1\mathrm{}i_n}^{(n)}|g_{T+T^{}}|\stackrel{~}{\gamma }_{(l)}^{k_1\mathrm{}k_l}`$ (5.6)
Let us illustrate the relation (5.6) by a couple of examples.
### 5.1 Examples
1) Let $`\stackrel{~}{\gamma }=\mathrm{}`$ and take a double-vertex graph $`\mathrm{\Gamma }_{3,4,;2}^{(3)}`$ for $`\mathrm{\Gamma }`$. Then
$`\mathrm{\Gamma }_{i_0;i_1i_2}^{(3)}|g_T|\mathrm{}={\displaystyle \underset{m,n,\stackrel{~}{m},\stackrel{~}{n}}{}}T_{i_0mn}^{(3)}G^{m\stackrel{~}{m}}G^{n\stackrel{~}{n}}T_{i_1i_2\stackrel{~}{m}\stackrel{~}{n}}^{(4)}={\displaystyle \underset{mn}{}}T_{i_0}^{mn}T_{i_1i_2mn}`$ (5.7)
For the sake of brevity we omitted the labels $`(3)`$ and $`(4)`$ in coupling constants.
Four different vertex-subgraphs $`\gamma `$ contribute to the sum in (5.6):
$`\gamma ^{(0)}=\mathrm{};\gamma ^{(3)}=\mathrm{\Gamma }_{3;0}^{(3)};\gamma ^{(4)}=\mathrm{\Gamma }_{4;0}^{(4)}and\gamma ^{(3)}=\mathrm{\Gamma }_{3,4;2}^{(3)}`$ (5.8)
The corresponding
$`\mathrm{\Gamma }^{(3)}/\gamma ^{(0)}=\mathrm{\Gamma }^{(3)};\mathrm{\Gamma }^{(3)}/\gamma ^{(3)}=\mathrm{\Gamma }_{4;0}^{(4)};\mathrm{\Gamma }^{(3)}/\gamma ^{(4)}=\mathrm{\Gamma }_{3;0}^{(3)};\mathrm{\Gamma }^{(3)}/\mathrm{\Gamma }^{(3)}=\mathrm{}`$ (5.9)
Eq.(5.6) states that
$`{\displaystyle \underset{m,n}{}}(T_{i_0}^{mn}T_{i_1i_2mn}1+T_{i_0}^{mn}T_{}^{}{}_{i_1i_2mn}{}^{}+T_{}^{}{}_{i_0}{}^{mn}T_{i_1i_2mn}+`$
$`+1T_{}^{}{}_{i_0}{}^{mn}T_{i_1i_2mn}^{})={\displaystyle }_{m,n}(T+T^{})_{i_0}^{mn}(T+T^{})_{i_1i_2mn}`$ (5.10)
what is indeed true. Note that in this check it is important that the metric $`G_{ij}`$ is the same for all the three theories $`g_T`$, $`g_T^{}`$ and $`g_{T+T^{}}`$.
2) Let $`\mathrm{\Gamma }`$ be a chain graph
$`\mathrm{\Gamma }_N^{(2)}=C_N`$ (5.11)
with $`N`$ vertices of valence two, connected by $`N1`$ propagators. Then
$`\mathrm{\Gamma }_N^{(2)}|g_T|\mathrm{}_{ij}={\displaystyle \underset{\stackrel{i_1,\mathrm{}i_{N1}}{j_1,\mathrm{},j_{N1}}}{}}T_{ij_1}^{(2)}G^{j_1i_1}T_{i_1j_2}^{(2)}G^{j_2i_2}\mathrm{}G^{j_{N1}i_{N1}}T_{i_{N1}j}^{(2)}`$ (5.12)
In this case $`\gamma `$ and $`\stackrel{~}{\gamma }`$ in (5.6) can be any collections of disconnected chains of the same type with the total length of no more than $`N`$. If $`\stackrel{~}{\gamma }=\mathrm{}`$, there are as many as $`2^N`$ possible choices of vertex-subgraphs $`\gamma `$ in (5.6), specified by all possible subsets of $`N`$ crosses in (5.11). In particular, there are $`\frac{N!}{k!(Nk)!}`$ vertex-subgraphs with $`k`$ vertices (connected and disconnected), and in the single-scalar case the identity (5.6) is just the binomial formula:
$`{\displaystyle \frac{1}{G^{N1}}}{\displaystyle \underset{k=0}{\overset{N}{}}}{\displaystyle \frac{N!}{k!(Nk)!}}T_{(2)}^k(T_{(2)}^{})^{Nk}={\displaystyle \frac{(T_{(2)}+T_{(2)}^{})^N}{G^{N1}}}`$ (5.13)
One can easily restore the indices $`i`$ and also consider non-trivial subchains $`\stackrel{~}{\gamma }`$.
## 6 Two Hopf algebras of graphs
The universal set-theoretical Hopf algebra defines a product of two graphs $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ to be a disconnected graph with components $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$,
$`\mathrm{\Gamma }_1\mathrm{\Gamma }_2=\mathrm{\Gamma }_1\mathrm{\Gamma }_2\mathrm{𝑓𝑜𝑟}\mathrm{\Gamma }_1\mathrm{\Gamma }_2=\mathrm{}`$ (6.1)
(the role of unity is played by the empty graph $`\mathrm{}`$), and the coproduct
$`\mathrm{\Delta }_{ST}(\mathrm{\Gamma })={\displaystyle \underset{all\gamma 𝒱\mathrm{\Gamma }}{}}\gamma \mathrm{\Gamma }/\gamma `$ (6.2)
This Hopf algebra is both commutative and cocommutative, associative and coassociative. Because of its cocommutativity, it is not associated with any non-trivial Lie algebra (the dual algebra $`\mathrm{𝑆ℎ𝑖𝑓𝑡}`$, introduced in (5.6), is obviously commutative: $`g_Tg_T^{}=g_{T+T^{}}=g_T^{}g_T`$). One can define such a Hopf algebra not only on graphs, but on any set and its subsets and we call it the “set-theory” (ST) Hopf algebra.
For functions on graphs, taking values in some commutative associative ring $`𝒦`$, one can define ST multiplication:
$`(F_{ST}G)(\mathrm{\Gamma })=m((FG)(\mathrm{\Delta }_{ST}(\mathrm{\Gamma }))={\displaystyle \underset{\gamma 𝒱\mathrm{\Gamma }}{}}F(\gamma )G(\mathrm{\Gamma }/\gamma ),`$
(operation $`m`$ multiplies two components of the tensor product: $`m((F(\gamma _1)G(\gamma _2))=F(\gamma _1)G(\gamma _2)`$).
Using the specifics of graphs, one can substitute vertex-subgraphs in (6.2) by box-subgraphs and construct a non-cocommutative comultiplication . First of all, the matrix element of $`g_T`$ for contracted graph $`[\mathrm{\Gamma }/\gamma ]`$, obtained by contraction of a box-subgraph $`\gamma ^{(m)}`$ in $`\mathrm{\Gamma }^{(n)}`$, is given by
$`\left[\mathrm{\Gamma }/\gamma \right]_{i_1\mathrm{}i_n}|g_T|\mathrm{}={\displaystyle \underset{j}{}}\mathrm{\Gamma }_{i_1\mathrm{}i_n}^{(n)}|g_T|\gamma _{(m)}^{j_1\mathrm{}j_m}T_{j_1\mathrm{}j_m}^{(m)}`$ (6.4)
for connected $`\gamma ^{(m)}`$,
$`\left[\mathrm{\Gamma }/(\gamma _1\gamma _2)\right]_{i_1\mathrm{}i_n}|g_T|\mathrm{}=`$
$`={\displaystyle \underset{j,k}{}}\mathrm{\Gamma }_{i_1\mathrm{}i_n}^{(n)}|g_T|\gamma _{1(m_1)}^{j_1\mathrm{}j_{m_1}}\gamma _{2(m_2)}^{k_1\mathrm{}k_{m_2}}T_{j_1\mathrm{}j_{m_1}}^{(m_1)}T_{k_1\mathrm{}k_{m_2}}^{(m_2)}`$ (6.5)
for $`\gamma ^{(m)}`$ consisting of two connected parts, and so on.
The Connes-Kreimer (CK) comultiplication
$`\mathrm{\Delta }_{CK}\mathrm{\Gamma }=\mathrm{}\mathrm{\Gamma }+\mathrm{\Gamma }\mathrm{}+{\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}\gamma \left[\mathrm{\Gamma }/\gamma \right],\mathrm{𝐶𝑜𝑛}(\mathrm{\Gamma })=1,`$
$`\mathrm{\Delta }_{CK}(\mathrm{\Gamma }_1\mathrm{\Gamma }_2)=\mathrm{\Delta }_{CK}(\mathrm{\Gamma }_1)\mathrm{\Delta }_{CK}(\mathrm{\Gamma }_2)`$ (6.6)
and the CK product
$`(F_{CK}G)(\mathrm{\Gamma })=m((FG)(\mathrm{\Delta }_{CK}(\mathrm{\Gamma }))=`$
$`\stackrel{\mathrm{𝐶𝑜𝑛}(\mathrm{\Gamma })=1}{=}F(\mathrm{})G(\mathrm{\Gamma })+F(\mathrm{\Gamma })G(\mathrm{})+{\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}F(\gamma )G([\mathrm{\Gamma }/\gamma ])`$ (6.7)
are no longer cocommutative. Thus the dual algebra is the universal enveloping of non-trivial Lie algebra. This Lie algebra, $`\mathrm{𝑑𝑖𝑓𝑓}_{\mathrm{}}`$, is straightforwardly realized by vector fields on the moduli space $``$ of coupling constants. The universal model (2.1) provides a basis in $`\mathrm{𝑑𝑖𝑓𝑓}_{\mathrm{}}`$, labeled by connected graphs: for any connected $`\mathrm{\Gamma }^{(n)}`$ one explicitly defines $`\widehat{Z}_\mathrm{\Gamma }T`$ as
$`\widehat{Z}_\mathrm{\Gamma }={\displaystyle \underset{i}{}}\mathrm{\Gamma }_{i_1\mathrm{}i_n}^{(n)}|g_T|\mathrm{}{\displaystyle \frac{}{T_{i_1\mathrm{}i_n}^{(n)}}}={\displaystyle \underset{i}{}}Z_{i_1\mathrm{}i_n}^{(n)}\{T\}{\displaystyle \frac{}{T_{i_1\mathrm{}i_n}^{(n)}}}`$ (6.8)
In what follows we denote by hats the vector fields and other elements of the universal enveloping $`𝒰(\mathrm{𝑑𝑖𝑓𝑓}_{\mathrm{}})`$ to distinguish them from scalars and other elements of modules (representations) of $`\mathrm{𝑑𝑖𝑓𝑓}_{\mathrm{}}`$. The commutator
$`[\widehat{Z}_{\mathrm{\Gamma }_1},\widehat{Z}_{\mathrm{\Gamma }_2}]=\widehat{Z}_{[\mathrm{\Gamma }_1,\mathrm{\Gamma }_2]},`$ (6.9)
where commutator $`[\mathrm{\Gamma }_1,\mathrm{\Gamma }_2]`$ is a linear combination of all graphs $`\mathrm{\Gamma }`$, such that $`[\mathrm{\Gamma }/\mathrm{\Gamma }_1]=\mathrm{\Gamma }_2`$ – these enter with the coefficient $`+1`$,– or $`[\mathrm{\Gamma }/\mathrm{\Gamma }_2]=\mathrm{\Gamma }_1`$ – these enter with $`1`$. (i.e. one blows any valence-$`m`$ vertex in $`\mathrm{\Gamma }_2`$ by insertion of $`\mathrm{\Gamma }_1^{(m)}`$ and any valence-$`n`$ vertex in $`\mathrm{\Gamma }_1`$ by gluing in $`\mathrm{\Gamma }_2^{(n)}`$ and takes an algebraic sum over such graphs with insertions.)
Disconnected graphs are associated with higher-order differential operators, e.g.
$`\widehat{Z}_{\mathrm{\Gamma }^{(n_1)}\mathrm{\Gamma }^{(n_2)}}={\displaystyle \underset{i,j}{}}\mathrm{\Gamma }_{i_1\mathrm{}i_{n_1}}^{(n_1)}\mathrm{\Gamma }_{j_1\mathrm{}j_{n_2}}^{(n_2)}|g_T|\mathrm{}{\displaystyle \frac{^2}{T_{i_1\mathrm{}i_{n_1}}^{(n_1)}T_{j_1\mathrm{}j_{n_2}}^{(n_2)}}}=`$
$`=:\widehat{Z}_{\mathrm{\Gamma }^{(n_1)}}\widehat{Z}_{\mathrm{\Gamma }^{(n_2)}}:\widehat{Z}_{\mathrm{\Gamma }^{(n_1)}}\widehat{Z}_{\mathrm{\Gamma }^{(n_2)}}`$ (6.10)
In other words, we associate with disconnected graphs the normal ordered products of vector fields, corresponding to each connected component. This provides a differential operator of certain order, equal to the number of connected components.
The vacuum graphs with no external legs define vector fields in the $`T^{(0)}`$ direction:
$`\widehat{Z}_{\mathrm{\Gamma }^{(0)}}=Z_{\mathrm{\Gamma }^{(0)}}\{T\}{\displaystyle \frac{}{T^{(0)}}}`$ (6.11)
The matrix elements $`\mathrm{\Gamma }|g_T|\gamma `$ can be associated either with Beltrami differentials:
$`\mu _{\mathrm{\Gamma }/\gamma }={\displaystyle \underset{i,j}{}}\mathrm{\Gamma }_{i_1\mathrm{}i_n}|g_T|\gamma ^{j_1\mathrm{}j_m}dT_{j_1\mathrm{}j_m}^{(m)}{\displaystyle \frac{}{T_{i_1\mathrm{}i_n}^{(n)}}}`$ (6.12)
(for connected $`\mathrm{\Gamma }`$ and $`\gamma `$),
$`\mu _{\mathrm{\Gamma }/(\gamma _1\gamma _2)}={\displaystyle \underset{i,j,k}{}}\mathrm{\Gamma }_{i_1\mathrm{}i_n}|g_T|\gamma _1^{j_1\mathrm{}j_{m_1}}\gamma _2^{k_1\mathrm{}k_{m_2}}dT_{j_1\mathrm{}j_m}^{(m_1)}dT_{k_1\mathrm{}k_{m_2}}^{(m_2)}{\displaystyle \frac{}{T_{i_1\mathrm{}i_n}^{(n)}}}`$ (6.13)
(for connected $`\mathrm{\Gamma }`$, $`\gamma _1`$ and $`\gamma _2`$) etc; or with the vector fields
$`\widehat{Z}_{[\mathrm{\Gamma }/\gamma ]}={\displaystyle \underset{i,j}{}}\mathrm{\Gamma }_{i_1\mathrm{}i_n}|g_T|\gamma ^{j_1\mathrm{}j_m}T_{j_1\mathrm{}j_m}^{(m)}{\displaystyle \frac{}{T_{i_1\mathrm{}i_n}^{(n)}}}`$ (6.14)
(for connected $`\mathrm{\Gamma }`$ and $`\gamma `$). Note that the only difference between (6.12) and (6.14) is in the letter $`d`$ in front of $`T^{(m)}`$, but it makes a lot of difference.
Operators $`g_T`$ form a subgroup in abelian group $`\mathrm{𝑆ℎ𝑖𝑓𝑡}`$, which acts transitively on the moduli space $``$. They are complemented by the non-abelian subgroup $`\mathrm{𝐷𝑖𝑓𝑓}_{\mathrm{}}()`$ of $`\mathrm{𝐷𝑖𝑓𝑓}`$, which is generated by the vector fields $`\widehat{Z}_\mathrm{\Gamma }`$, defined in (6.8), and is the stability subgroup of the Gaussian point $`T^{(n)}=0`$ (since $`\mathrm{\Gamma }|g_{T=0}|\mathrm{}=\delta _{\mathrm{\Gamma },\mathrm{}}`$ and all $`Z_\mathrm{\Gamma }(T=0)=0`$). The moduli space itself can be represented as a homogeneous factor-space $`=\mathrm{𝐷𝑖𝑓𝑓}()/\mathrm{𝐷𝑖𝑓𝑓}_{\mathrm{}}()=\mathrm{𝑆ℎ𝑖𝑓𝑡}()/\mathrm{𝑆ℎ𝑖𝑓𝑡}_{\mathrm{}}()`$. The action of $`\mathrm{𝐷𝑖𝑓𝑓}_{\mathrm{}}()`$ on non-Gaussian models is relevant for description of renormalization group flows in $``$.
The Lie algebra $`\mathrm{𝑑𝑖𝑓𝑓}_{\mathrm{}}`$ of vector fields $`\widehat{Z}_\mathrm{\Gamma }`$ on entire $``$ has a variety of reductions to smaller Lie algebras on subspaces $`_{red}`$, i.e. there are Lie algebras associated with smaller families of models than the universal (2.1). For example, one can consider only interactions of a given valence, i.e. all $`T^{(n)}=0`$ for $`nk`$, then vector fields (6.8) associated with connected graphs with exactly $`k`$ external legs form a closed Lie subalgebra. Alternative reduction is to the tree graphs ($`Z_\mathrm{\Gamma }=0`$ for any $`\mathrm{\Gamma }`$ with loops).
One can also consider the finite sets of indices $`I=\{1,\mathrm{},N\}`$ and accordingly reduced moduli spaces $`^{(N)}`$, in this case
$`Vect(^{(N)})=Vect()|_{^{(N)}}`$ (6.15)
## 7 Inverse operators, projectors and R-Operation
Due to (5.3) partition functions are characters of the graph multiplication: $`\mathrm{\Gamma }_1\mathrm{\Gamma }_2|g_T|\mathrm{}=\mathrm{\Gamma }_1|g_T|\mathrm{}\mathrm{\Gamma }_2|g_T|\mathrm{}`$, $`\mathrm{}|g_T\mathrm{}=1`$.
Characters take values in some commutative associative ring $`𝒦`$ and satisfy:
$`F(\mathrm{})=1,`$
$`F(\mathrm{\Gamma }_1\mathrm{\Gamma }_2)=F(\mathrm{\Gamma }_1)F(\mathrm{\Gamma }_2),`$
$`(F_{ST}G)(\mathrm{\Gamma })=m((FG)(\mathrm{\Delta }_{ST}(\mathrm{\Gamma }))={\displaystyle \underset{\gamma 𝒱\mathrm{\Gamma }}{}}F(\gamma )G(\mathrm{\Gamma }/\gamma ),`$
$`(F_{CK}G)(\mathrm{\Gamma })=m((FG)(\mathrm{\Delta }_{CK}(\mathrm{\Gamma }))=`$
$`\stackrel{\mathrm{𝐶𝑜𝑛}(\mathrm{\Gamma })=1}{=}F(\mathrm{\Gamma })+G(\mathrm{\Gamma })+{\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}F(\gamma )G([\mathrm{\Gamma }/\gamma ])`$ (7.1)
(operation $`m`$ multiplies two components of the tensor product: $`m((F(\gamma _1)G(\gamma _2))=F(\gamma _1)G(\gamma _2)`$).
One can define the inverses (antipodes) of a character $`F`$, $`F_{ST}^1`$, $`F_{CK}^1`$, which satisfy $`F_{ST}^1_{ST}F(\mathrm{\Gamma })=\delta _{\mathrm{\Gamma },\mathrm{}}`$, $`F_{CK}^1_{CK}F(\mathrm{\Gamma })=\delta _{\mathrm{\Gamma },\mathrm{}}`$, by recursive formulas:
$`F_{ST}^1(\mathrm{})=F_{CK}^1(\mathrm{})=1,`$
$`F_{ST}^1(\mathrm{\Gamma })=F(\mathrm{\Gamma }){\displaystyle \underset{\gamma 𝒱\mathrm{\Gamma };\gamma \mathrm{},\mathrm{\Gamma }}{}}F_{ST}^1(\gamma )F(\mathrm{\Gamma }/\gamma ),`$
$`F_{CK}^1(\mathrm{\Gamma })\stackrel{\mathrm{𝐶𝑜𝑛}(\mathrm{\Gamma })=1}{=}F(\mathrm{\Gamma }){\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}F_{CK}^1(\gamma )F([\mathrm{\Gamma }/\gamma ])`$ (7.2)
According to (5.6), if $`F(\mathrm{\Gamma })=\mathrm{\Gamma }|g_T|\mathrm{}`$, then $`F_{ST}^1(\mathrm{\Gamma })=\mathrm{\Gamma }|g_T|\mathrm{}`$, but $`F_{CK}^1(\mathrm{\Gamma })`$ is given by a more sophisticated expression. In fact, eq.(7.2) for $`F_{CK}^1(\mathrm{\Gamma })`$ is closely associated with Bogolubov’s recursive formula, defining the $`R`$-operation .
Assume that the ring $`𝒦`$ as a linear space can be decomposed into two components, $`𝒦=𝒦_1𝒦_2`$, with the help of projectors $`𝒫_\pm `$, $`𝒫_\pm ^2=𝒫_\pm `$, $`𝒫_{}=I𝒫_+`$: $`𝒦_\pm =𝒫_\pm 𝒦`$. These projectors can be used to define the “$`𝒫`$-inverse” ($`𝒫`$-antipode) $`PF^1`$ of $`F(\mathrm{\Gamma })`$ :
$`𝒫_{}\left((PF_{}^1F)(\mathrm{\Gamma })\delta _{\mathrm{\Gamma },\mathrm{}}\right)=0,`$
$`𝒫_+\left(PF_{}^1(\mathrm{\Gamma })\right)=0`$ (7.3)
The second condition makes the definition of $`𝒫`$-antipode unambiguous. The ordinary inverse $`F_{}^1`$ is associated with the trivial projector $`𝒫_+=0`$. One can easily write down recursive formulae for the $`𝒫`$-antipodes for comultiplications $`_{ST}`$ and $`_{CK}`$ by applying $`𝒫_{}`$ to the r.h.s. of (7):
$`PF_{ST}^1(\mathrm{\Gamma })=𝒫_{}\left(F(\mathrm{\Gamma })+{\displaystyle \underset{\gamma 𝒱\mathrm{\Gamma };\gamma \mathrm{},\mathrm{\Gamma }}{}}PF_{ST}^1(\gamma )F(\mathrm{\Gamma }/\gamma )\right),`$
$`PF_{CK}^1(\mathrm{\Gamma })\stackrel{\mathrm{𝐶𝑜𝑛}(\mathrm{\Gamma })=1}{=}𝒫_{}\left(F(\mathrm{\Gamma })+{\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}PF_{CK}^1(\gamma )F([\mathrm{\Gamma }/\gamma ])\right)`$ (7.4)
For projector $`𝒫_+`$, possessing additional triangular property w.r.to multiplication in $`𝒦`$, namely
$`𝒦_+𝒦_+𝒦_+,𝒦_{}𝒦_{}𝒦_{}`$ (7.5)
(i.e. the product of any two elements from $`𝒦_+`$ lies again in $`𝒦_+`$ and similarly for $`𝒦_{}`$), the ST $`𝒫`$-inverse of $`F(\mathrm{\Gamma })`$ is a character whenever $`F(\mathrm{\Gamma })`$ is a character:
$`PF_{ST}^1(\mathrm{\Gamma }_1\mathrm{\Gamma }_2)=PF_{ST}^1(\mathrm{\Gamma }_1)PF_{ST}^1(\mathrm{\Gamma }_2)`$ (7.6)
if
$`F(\mathrm{\Gamma }_1\mathrm{\Gamma }_2)=F(\mathrm{\Gamma }_1)F(\mathrm{\Gamma }_2)\mathrm{\Gamma }_1,\mathrm{\Gamma }_2`$ (7.7)
Indeed, assume that this is true for all smaller vertex-subgraphs of $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$. Then
$`PF_{ST}^1(\mathrm{\Gamma }_1\mathrm{\Gamma }_2)=𝒫_{}(F(\mathrm{\Gamma }_1\mathrm{\Gamma }_2)+`$
$`+{\displaystyle \underset{\gamma _1\mathrm{\Gamma }_1}{}}{\displaystyle \underset{\gamma _2\mathrm{\Gamma }_2}{}}PF_{ST}^1(\gamma _1)PF_{ST}^1(\gamma _2)F(\mathrm{\Gamma }_1/\gamma _1)F(\mathrm{\Gamma }_2/\gamma _2)`$
$`F(\mathrm{\Gamma }_1)F(\mathrm{\Gamma }_2)PF_{ST}^1(\mathrm{\Gamma }_1)PF_{ST}^1(\mathrm{\Gamma }_2))`$ (7.8)
The last two items at the r.h.s. subtract the contributions from $`\gamma _1\gamma _2=\mathrm{}`$ and $`\gamma _1\gamma _2=\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$. The double sum in (7.8) is equal to the product of two sums, defining the ST $`𝒫`$-inverses of $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$, which (the sums) are both $`𝒫`$-positive. Due to triangularity the product is also $`𝒫`$-positive and is eliminated by $`𝒫_{}`$. Therefore (7.8) states that
$`PF_{ST}^1(\mathrm{\Gamma }_1\mathrm{\Gamma }_2)=𝒫_{}\left(PF_{ST}^1(\mathrm{\Gamma }_1)PF_{ST}^1(\mathrm{\Gamma }_2)\right)=PF_{ST}^1(\mathrm{\Gamma }_1)PF_{ST}^1(\mathrm{\Gamma }_2)`$ (7.9)
The last equality is again implied by triangularity, since both $`𝒫`$-inverses are $`𝒫`$-negative.
Not every projector is triangular, for example projection on positive numbers in the ring of reals is not triangular: the product of two negatives is no longer negative. A natural triangular projector exists in a ring of Laurent series $`\left\{A=_{k=N}^{\mathrm{}}a_kz^k\right\}`$: $`𝒫_+A=_{k=0}^{\mathrm{}}a_kz^k`$. The difference $`=𝒫_+𝒫_{}`$ is the $`r`$-matrix, widely used in the theory of integrable systems and its applications (see, for example, and also ). To get a field-theory model with such $`𝒦`$ one can, for example, consider the $`z`$-dependent couplings $`T^{(n)}=_{k=N}^{\mathrm{}}T_k^{(n)}z^k`$ in (2.1). In the study of continuous field theory $`z`$ rather enters through regularization of infinite sums (integrals) over indices $`i`$ in (2.1): it can be identified with $`dd_{crit}`$ for dimensional regularization or with $`1/M`$ for Pauli-Villars regularization etc.
According to (7.3), the $`R`$-operation
$`F(\mathrm{\Gamma })RF_{}(\mathrm{\Gamma })=(PF_{}^1F)(\mathrm{\Gamma }),`$
$`𝒫_{}\left(RF_{}(\mathrm{\Gamma })\right)=0,`$ (7.10)
acting on the space of functions of graphs, converts any function into a $`𝒫`$-positive (“finite”) one. Moreover, since $``$-product of characters is again a character, it converts characters into characters. The main claim of is that eq.(7.10) for $`_{CK}`$ can be considered as group-theory interpretation of Bogolubov’s recursion formula . In sec.8 we shall see that more relevant in generic case is the corepresentation $`\widehat{}_{CK}`$.
Of course, from algebraic perspective there is nothing special about continuous theory, divergencies and dimensional regularization: the only things that matter are algebraic structures and triangular projectors.
## 8 Representations of $`\mathrm{𝑑𝑖𝑓𝑓}_{\mathrm{}}`$ and $`𝒰(\mathrm{𝑑𝑖𝑓𝑓}_{\mathrm{}})`$ in differential operators on $``$
Returning to the beginning of sec.7, model (2.1) provides $`Z_\mathrm{\Gamma }\{T\}`$ for the role of characters, if these quantities are considered as functions of $`\mathrm{\Gamma }`$, and $`T`$-dependence is not taken into account. Similar treatment can be given to vacuum diagrams in generic (2.1). However, it is more adequate to treat $`Z_\mathrm{\Gamma }\{T\}=\mathrm{\Gamma }|T`$ as describing a transformation between the functions of $`T`$ and functions of $`\mathrm{\Gamma }`$ (not obligatory characters), as suggested in sec.3 above.
Consider the action of vector fields on linear modules over $``$. Namely, a vector field
$`\widehat{V}={\displaystyle \underset{n}{}}\left({\displaystyle \underset{i_1,\mathrm{},i_n}{}}V_{i_1\mathrm{}i_n}^{(n)}(T){\displaystyle \frac{}{T_{i_1\mathrm{}i_n}}}\right)`$ (8.1)
can act on a function of $`T`$-variables $`F\{T\}`$ with or without free indices:
$`F\{T\}\widehat{V}\{T\}F\{T\}`$ (8.2)
Now we can exploit the power given by the use of the universal model (2.1). It provides a large enough set of functions on $``$ to establish the one-to-one correspondence between linear combinations of graphs and invariant functions of coupling constants (while for smaller models the set of such functions is much smaller: graphs label different types of contracting indices, and there should be many enough indices to distinguish between different contractions). Because of this, every invariant (i.e. with all indices $`i`$ contracted with the help of the metric $`G_{ij}`$) function on $``$ can be uniquely decomposed into a sum over graphs of the basic functions $`Z_\mathrm{\Gamma }\{T\}`$, introduced in sec.3 (of course, such expansions survive certain reductions of $``$, but this is a separate story). Actually, functions are decomposed into sums over vacuum graphs $`\mathrm{\Gamma }^{(0)}`$ without external legs,
$`F\{T\}={\displaystyle \underset{\mathrm{\Gamma }^{(0)}}{}}F(\mathrm{\Gamma })Z_\mathrm{\Gamma }\{T\}`$ (8.3)
with $`T`$-independent coefficients $`F_\mathrm{\Gamma }`$; vector fields – over connected graphs with any number of external legs,
$`\widehat{V}\{T\}={\displaystyle \underset{n1}{}}\left({\displaystyle \underset{connected\mathrm{\Gamma }^{(n)}}{}}V(\mathrm{\Gamma })\widehat{Z}_\mathrm{\Gamma }\{T\}\right);`$ (8.4)
the $`k`$-differentials – over graphs with $`k`$ connected components and non-vanishing number of external legs in each component,
$`:\widehat{W}_k\{T\}:=`$
$`={\displaystyle \underset{n_1,\mathrm{}n_k1}{}}({\displaystyle \underset{\stackrel{connected}{\mathrm{\Gamma }_1^{(n_1)},\mathrm{},\mathrm{\Gamma }_k^{(n_k)}}}{}}W(\mathrm{\Gamma }_1\mathrm{}\mathrm{\Gamma }_k):\widehat{Z}_{\mathrm{\Gamma }_1}\{T\}\mathrm{}\widehat{Z}_{\mathrm{\Gamma }_k}\{T\}:);`$ (8.5)
generic elements of the universal module (generic differential operators on $``$) – over all possible graphs (with any number of connected components and external legs). In what follows $`F\{T\}`$ can be arbitrary element of the universal module. Also, we assume that for a vector field the coefficients $`V(\mathrm{\Gamma })`$ are defined for all graphs $`\mathrm{\Gamma }`$, just $`V(\mathrm{\Gamma })=0`$ if $`\mathrm{\Gamma }`$ is not connected (of course, $`V(\mathrm{\Gamma })`$ is not a character, characters are associated with group elements $`G_{\widehat{V}}=e^{\widehat{V}}`$, not vector fields themselves).
The result of the action of $`\widehat{V}`$ on $`F`$ can also be decomposed in the basis $`Z_\mathrm{\Gamma }\{T\}`$,
$`\widehat{V}\{T\}F\{T\}={\displaystyle \underset{\mathrm{\Gamma }}{}}(\widehat{V}F)(\mathrm{\Gamma })Z_\mathrm{\Gamma }\{T\}`$ (8.6)
and one obtains a relation between the coefficients $`(\widehat{V}F)(\mathrm{\Gamma })`$, $`F(\mathrm{\Gamma })`$ and $`V(\mathrm{\Gamma })`$: since
$`{\displaystyle \underset{\mathrm{\Gamma }}{}}(\widehat{V}F)(\mathrm{\Gamma })Z_\mathrm{\Gamma }\{T\}={\displaystyle \underset{\mathrm{\Gamma }^{}}{}}{\displaystyle \underset{\gamma }{}}V(\gamma )F(\mathrm{\Gamma }^{})\left(\widehat{Z}_\gamma (T)Z_\mathrm{\Gamma }^{}\{T\}\right),`$ (8.7)
and
$`\widehat{Z}_\gamma \{T\}Z_\mathrm{\Gamma }^{}\{T\}={\displaystyle \underset{\mathrm{\Gamma }:\mathrm{\Gamma }^{}=[\mathrm{\Gamma }/\gamma ]}{}}Z_\mathrm{\Gamma }\{T\}`$ (8.8)
we get a convolution formula
$`(\widehat{V}F)(\mathrm{\Gamma })={\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}V(\gamma )F([\mathrm{\Gamma }/\gamma ])=(V\widehat{}_{CK}F)(\mathrm{\Gamma })`$ (8.9)
Operation $`\widehat{}_{CK}`$,
$`(W\widehat{}_{CK}F)(\mathrm{\Gamma })=m((WF)(\widehat{\mathrm{\Delta }}_{CK}\mathrm{\Gamma }))=`$
$`\stackrel{\mathrm{𝐶𝑜𝑛}(\mathrm{\Gamma })=1}{=}W(\mathrm{})F(\mathrm{\Gamma })+{\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}W(\gamma )F([\mathrm{\Gamma }/\gamma ]),`$ (8.10)
is expressed in terms of the corepresentation of the CK Hopf algebra of graphs,
$`\widehat{\mathrm{\Delta }}_{CK}\mathrm{\Gamma }=\mathrm{}\mathrm{\Gamma }+{\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}\gamma [\mathrm{\Gamma }/\gamma ],\mathrm{𝐶𝑜𝑛}(\mathrm{\Gamma })=1,`$ (8.11)
the same way as $`_{CK}`$, eq.(6.7), is expressed through the comultiplication
$`\mathrm{\Delta }_{CK}\mathrm{\Gamma }=\mathrm{}\mathrm{\Gamma }+\mathrm{\Gamma }\mathrm{}+{\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}\gamma [\mathrm{\Gamma }/\gamma ],\mathrm{𝐶𝑜𝑛}(\mathrm{\Gamma })=1`$ (8.12)
In (8.9) $`\widehat{V}`$ is a vector field, therefore $`V(\mathrm{})=0`$. The difference between comultiplication $`\mathrm{\Delta }`$ and corepresentation $`\widehat{\mathrm{\Delta }}`$ is in associativity conditions:
$`(\mathrm{\Delta }id)\mathrm{\Delta }=(id\mathrm{\Delta })\mathrm{\Delta }`$ (8.13)
for $`\mathrm{\Delta }`$ and
$`(\mathrm{\Delta }id)\widehat{\mathrm{\Delta }}=(id\widehat{\mathrm{\Delta }})\widehat{\mathrm{\Delta }}`$ (8.14)
for $`\widehat{\mathrm{\Delta }}`$.
Repeated application of formula (8.9) defines the action of products and normal ordered products of vector fields on $`F`$. For two vectors, since
$`\widehat{Z}_{\gamma _1}\left(\widehat{Z}_{\gamma _2}Z_{\mathrm{\Gamma }^{\prime \prime }}\right)={\displaystyle \underset{\mathrm{\Gamma }:[\mathrm{\Gamma }/\gamma _1]=\mathrm{\Gamma }^{}}{}}\left({\displaystyle \underset{\mathrm{\Gamma }^{}:[\mathrm{\Gamma }^{}/\gamma _2]=\mathrm{\Gamma }^{\prime \prime }}{}}Z_\mathrm{\Gamma }\right)`$ (8.15)
and
$`:\widehat{Z}_{\gamma _1}\widehat{Z}_{\gamma _2}:Z_\mathrm{\Gamma }^{}={\displaystyle \underset{\mathrm{\Gamma }:[\mathrm{\Gamma }/(\gamma _1\gamma _2)]=\mathrm{\Gamma }^{}}{}}Z_\mathrm{\Gamma },`$ (8.16)
we have:
$`(\widehat{V}_1\widehat{V}_2F)(\mathrm{\Gamma })={\displaystyle \underset{\gamma _2\mathrm{\Gamma }}{}}V_2(\gamma _2)\left({\displaystyle \underset{\gamma _1[\mathrm{\Gamma }/\gamma _2]}{}}V_1(\gamma _1)F(\left[[\mathrm{\Gamma }/\gamma _2]/\gamma _1\right])\right)`$ (8.17)
and
$`(:\widehat{V}_1\widehat{V}_2:F)(\mathrm{\Gamma })={\displaystyle \underset{\gamma _1\gamma _2\mathrm{\Gamma }}{}}V_1(\gamma _1)V_2(\gamma _2)F([\mathrm{\Gamma }/(\gamma _1\gamma _2)])=`$
$`=(:V_1V_2:\widehat{}_{CK}F)(\mathrm{\Gamma })`$ (8.18)
Note that in these formulas $`\left[[\mathrm{\Gamma }/\gamma _2]/\gamma _1\right][\mathrm{\Gamma }/(\gamma _1\gamma _2)]`$: a graph $`\gamma _1`$ can appear after contraction $`[\mathrm{\Gamma }/\gamma _2]`$ is made. Relatively simple formulae in terms of the corepresentation $`\widehat{}_{CK}`$ exist only for the normal ordered elements $`:\widehat{W}:𝒰(\mathrm{𝑑𝑖𝑓𝑓}_{\mathrm{}})`$:<sup>4</sup><sup>4</sup>4 Note that the only component of $`𝒰(\mathrm{𝑑𝑖𝑓𝑓}_{\mathrm{}})`$ which has $`W(\mathrm{})0`$ is a counity, i.e. $`W\{T\}=const`$. Non-trivial functions $`\widehat{f}\{T\}𝒰(\mathrm{𝑑𝑖𝑓𝑓}_{\mathrm{}})`$, and ordinary product of functions is not expressible in terms of coproduct $`_{CK}`$. Instead the scalar $`\widehat{f}\{T\}=_{\mathrm{\Gamma }^{(0)}}f(\mathrm{\Gamma })Z_\mathrm{\Gamma }\{T\}`$ acts on $`F\{T\}`$ as
$$(\widehat{f}F)(\mathrm{\Gamma })=\underset{\mathrm{\Gamma }_1,\mathrm{\Gamma }_2:\mathrm{\Gamma }=\mathrm{\Gamma }_1\mathrm{\Gamma }_2}{}f(\mathrm{\Gamma }_1)F(\mathrm{\Gamma }_2)$$
In particular, for connected $`\mathrm{\Gamma }`$, $`(\widehat{f}F)(\mathrm{\Gamma })=f(\mathrm{\Gamma })+F(\mathrm{\Gamma })`$.
$`(:\widehat{W}:F)(\mathrm{\Gamma })\stackrel{\mathrm{𝐶𝑜𝑛}(\mathrm{\Gamma })=1}{=}:W:(\mathrm{})F(\mathrm{\Gamma })+{\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}:W:(\gamma )F([\mathrm{\Gamma }/\gamma ])=`$
$`=(:W:\widehat{}_{CK}F)(\mathrm{\Gamma })`$ (8.19)
Of special interest for us are specific elements of $`𝒰(\mathrm{𝑑𝑖𝑓𝑓}_{\mathrm{}})`$, which are the group elements, and form the diffeomorphism group $`\mathrm{𝐷𝑖𝑓𝑓}_{\mathrm{}}`$.
Given a vector field $`\widehat{V}=V^\alpha _\alpha `$ ($`\alpha `$ is a multiindex, labeling connected graph with indices or any linear combinations of such graphs), one can make an element of $`\mathrm{𝐷𝑖𝑓𝑓}_{\mathrm{}}()`$ by exponentiation:
$`G_{\widehat{V}}=e^{\widehat{V}}={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\widehat{V}^n}{n!}}`$ (8.20)
However, it is not normal ordered, and the action of $`G_{\widehat{V}}`$ on $`F(\mathrm{\Gamma })`$ is described by sophisticated expression:
$`(G_{\widehat{V}}F)(\mathrm{\Gamma })={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}\left({\displaystyle \underset{\{\gamma \}_n}{}}V(\gamma _1)\mathrm{}V(\gamma _n)F([\mathrm{\Gamma }/\{\gamma \}_n])\right)`$ (8.21)
where $`\{\gamma \}_n`$ denotes a hierarchy of subgraphs $`\gamma _n\mathrm{\Gamma }`$, $`\gamma _{n1}[\mathrm{\Gamma }/\gamma _n]`$, $`\mathrm{}`$, $`\gamma _1\left[\left[\mathrm{}\left[[\mathrm{\Gamma }/\gamma _n]/\gamma _{n1}\right]/\mathrm{}\right]/\gamma _2\right]`$.
One can instead expand $`G_{\widehat{V}}`$ in normal order constituents with the help of a forest formula:
$`G_{\widehat{V}}=e^{\widehat{V}}={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\widehat{V}^n}{n!}}=`$
$`=1+V^\alpha _\alpha +{\displaystyle \frac{1}{2}}V^\gamma _\gamma V^\alpha _\alpha +{\displaystyle \frac{1}{6}}V^\gamma _\gamma V^\beta _\beta V^\alpha _\alpha +\mathrm{}=`$
$`=1+(V^\alpha +{\displaystyle \frac{1}{2}}V^\gamma (_\gamma V^\alpha )+{\displaystyle \frac{1}{6}}V^\gamma (_\gamma V^\beta )(_\beta V^\alpha )+`$
$`+{\displaystyle \frac{1}{6}}V^\beta V^\gamma (_\beta _\gamma V^\alpha )+\mathrm{})_\alpha +`$
$`+{\displaystyle \frac{1}{2}}\left(V^\alpha +{\displaystyle \frac{1}{2}}V^\gamma (_\gamma V^\alpha )+\mathrm{}\right)\left(V^\beta +{\displaystyle \frac{1}{2}}V^\gamma (_\gamma V^\beta )+\mathrm{}\right)_\alpha _\beta +`$
$`+{\displaystyle \frac{1}{6}}(V^\alpha +\mathrm{})(V^\beta +\mathrm{})(V^\gamma +\mathrm{})_\alpha _\beta _\gamma +\mathrm{}=`$
$`=1+{\displaystyle \underset{}{}}{\displaystyle \frac{1}{\mathrm{𝑇𝑟𝑒𝑒}()!}}:{\displaystyle \underset{𝒯}{}}{\displaystyle \frac{\widehat{V}_𝒯}{\sigma _𝒯𝒯!}}:`$ (8.22)
The forest $``$ is an ordered set of rooted trees. Rooted tree has a single external leg (root), all other external legs end at the valence-one vertices. $`\mathrm{𝑇𝑟𝑒𝑒}()`$ is the number of trees in the forest, and $`\mathrm{𝑉𝑒𝑟𝑡}(𝒯)`$ is the number of vertices in the tree $`𝒯`$. For every rooted tree $`\sigma _𝒯`$ is the symmetry factor (the order of the discrete group which interchanges subtrees, leaving the tree intact), while the tree-factorial is defined iteratively: $`𝒯!=\mathrm{𝑉𝑒𝑟𝑡}(𝒯)_a𝒯_a!`$, where $`𝒯_a`$ are the root subtrees formed after the root is cut away. In every vertex of a tree stands the vector field $`\widehat{V}`$, acting on the neighbor vertex downwards (in the direction to the root), and not further. Then with every tree we associate a vector field $`\widehat{V}_𝒯`$, which contains the $`\mathrm{𝑉𝑒𝑟𝑡}(𝒯)`$’s power of $`\widehat{V}`$ and $`\mathrm{𝐿𝑖𝑛𝑘}(𝒯)`$ derivatives ($`\mathrm{𝐿𝑖𝑛𝑘}(𝒯)`$ is the number of links in the tree). For example, for the 1-vertex ($`𝒯_1`$), 2-vertex ($`𝒯_2`$) and 3-vertex/2-branch ($`𝒯_Y`$) trees:
$`V(𝒯_1)=V^\alpha _\alpha ,`$
$`V(𝒯_2)=V^\alpha (_\alpha V^\beta )_\beta ,`$
$`V(𝒯_Y)=V^\alpha V^\beta (_\alpha _\beta V^\gamma )_\gamma ,`$ (8.23)
etc. With a forest we associate a differential operator, which is a normal-ordered product of vector fields $`\widehat{V}_𝒯`$ over the trees (as usual, normal ordering means, that all derivatives are written to the right of $`V^\alpha `$’s, this is a coordinate-dependent operation, e.g. $`:\widehat{V}^3:=V^\alpha V^\beta V^\gamma _\alpha _\beta _\gamma `$).
One can apply (8.19) to obtain an alternative expression to (8.21) in terms of the corepresentation $`\widehat{}_{CK}`$. Complexity of the formula is now encoded in the sum over forests. One can efficiently handle this complexity by the following trick. Since $`G_{\widehat{V}}=e^{\widehat{V}}`$ is a diffeomorphism of $``$, for any $`F\{T\}`$ we have:
$`(G_{\widehat{V}}F)\{T\}=F\{T+\stackrel{~}{V}(T)\}`$ (8.24)
with
$`\stackrel{~}{V}_{i_1\mathrm{}i_n}^{(n)}=\left(e^{\widehat{V}}1\right)T_{i_1\mathrm{}i_n}^{(n)}=\left({\displaystyle \underset{𝒯}{}}{\displaystyle \frac{\widehat{V}_𝒯}{\sigma _𝒯𝒯!}}\right)T_{i_1\mathrm{}i_n}^{(n)}`$ (8.25)
The first equality in (8.25) is obtained by substitution of $`T_{i_1\mathrm{}i_n}^{(n)}`$ instead of $`F\{T\}`$ in (8.24), the second equality is implied by the forest formula (8.22), because a normal product of two or more vector fields annihilates $`T_{i_1\mathrm{}i_n}^{(n)}`$. Now introduce a new vector field
$`\widehat{\stackrel{~}{V}}\{T\}={\displaystyle \underset{n}{}}\left({\displaystyle \underset{i_1,\mathrm{},i_n}{}}\stackrel{~}{V}_{i_1\mathrm{}i_n}^{(n)}{\displaystyle \frac{}{T_{i_1\mathrm{}i_n}^{(n)}}}\right)={\displaystyle \underset{connected\mathrm{\Gamma }}{}}\stackrel{~}{V}(\mathrm{\Gamma })\widehat{Z}_\mathrm{\Gamma }\{T\},`$ (8.26)
such that the shift operator
$`G_{\widehat{V}}=e^{\widehat{V}}=:e^{\widehat{\stackrel{~}{V}}}:`$ (8.27)
Equality (8.27) is implied by (8.24) and by Taylor expansion
$`F\{T+\stackrel{~}{V}(T)\}=:e^{\widehat{\stackrel{~}{V}}}:F\{T\}`$ (8.28)
Now we can make use of (8.19) to obtain a simple substitute for (8.21):
$`(G_{\widehat{V}}F)(\mathrm{\Gamma })\stackrel{\mathrm{𝐶𝑜𝑛}(\mathrm{\Gamma })=1}{=}F(\mathrm{\Gamma })+{\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}G_{\widehat{V}}(\gamma )F([\mathrm{\Gamma }/\gamma ])=`$
$`={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}({\displaystyle \underset{\stackrel{nonintersecting}{\gamma _1,\mathrm{},\gamma _n\mathrm{\Gamma }}}{}}\stackrel{~}{V}(\gamma _1)\mathrm{}\stackrel{~}{V}(\gamma _n)F([\mathrm{\Gamma }/(\gamma _1\mathrm{}\gamma _n)])`$ (8.29)
## 9 Bogolubov’s recursion and renormalized Lagrangian
One can apply diffeomorphisms in moduli space to “improve” partition functions. This is important if one wants to eliminate undesired dependence on one or another parameter of the theory, like ultraviolet cut-off in continuous local field models. Basically, one needs to project the entire moduli space $``$ onto certain subspace $`_{ren}`$ of “renormalized models”. The problem is that parameter-dependence arises in partition functions, and arbitrary elimination of unwanted parameters from particular correlators can break the relation to Lagrangian formalism and moduli space. It is exactly the problem, which is resolved by Bogolubov’s $`R`$-operation , and which can be most straightforwardly described in terms of diffeomorphisms of $``$.
The $`R`$-operation can be formulated as follows:
Given (triangular) projectors $`𝒫_\pm `$ in the ring $`𝒦`$ (where matrix elements and partition functions are taking values) and a function $`F\{T\}`$ (with or without free indices, i.e. any element of the universal module over $``$), one finds a specific diffeomorphism $`G_{\widehat{P}}\mathrm{𝐷𝑖𝑓𝑓}_{\mathrm{}}`$ which makes $`F\{T\}`$ $`𝒫`$-positive:
$`𝒫_{}\left(F\{T+\stackrel{~}{P}(T)\}\right)=0,`$ (9.1)
$`F\{T+\stackrel{~}{P}(T)\}=e^{\widehat{P}}F\{T\}=(G_{\widehat{P}}F)\{T\}`$ (9.2)
To define such diffeomorphism unambiguously, one imposes additional constraint on $`\widehat{P}\{T\}`$, for example,
$`𝒫_+\left(\stackrel{~}{P}_{i_1\mathrm{}i_n}^{(n)}\right)=0n;i_1,\mathrm{},i_n`$ (9.3)
(see eq.(9.6) below for a more adequate constraint). Some constraint of this type is needed to distinguish between “renormalizations”, needed to eliminate $`𝒫`$-negative contributions to the correlation functions from arbitrary diffeomorphisms of $``$, which can map $`𝒫`$-positive models into other $`𝒫`$-positive ones.
Eqs.(9.1-9.3) define the Bogolubov’s $`R`$-operation for any projector $`𝒫_+`$. One can apply the machinery of the previous sections to rewrite (9.2) either in terms of Gauss-Birkhoff decomposition of the shift operator,
$`g_{T+\stackrel{~}{P}(T)}=g_Tg_{\stackrel{~}{P}(T)}`$ (9.4)
(decomposing $`𝒫`$-positive renormalized model into the bare one and $`𝒫`$-negative counterterm model), or in terms of CK algebra of functions on graphs. In the last case one can use any of the three representations (8.21), (8.22) or (8.29). The most convenient is the third choice, and it is exactly the one providing the Bogolubov’s recursion formula. Eq.(8.29) can indeed be rewritten in the form of a recurrent relation for $`\stackrel{~}{P}(\mathrm{\Gamma })`$, expressing it through $`\stackrel{~}{P}(\gamma )`$ for smaller box-subgraphs $`\gamma `$ (with less vertices), provided $`F(\mathrm{\Gamma })`$ does not vanish on elementary vertices $`[\mathrm{\Gamma }/\mathrm{\Gamma }]`$. Indeed, one can extract from the r.h.s. of (8.29) two items: one with $`n=0`$ and another with $`n=1`$ and $`\gamma =\mathrm{\Gamma }`$. Then we obtain:<sup>5</sup><sup>5</sup>5 In notation of $`G_{\widehat{P}}(\mathrm{\Gamma })=C(\mathrm{\Gamma })`$.
$`\stackrel{~}{P}(\mathrm{\Gamma })F([\mathrm{\Gamma }/\mathrm{\Gamma }])\stackrel{\mathrm{𝐶𝑜𝑛}(\mathrm{\Gamma })=1}{=}𝒫_{}\left(F(\mathrm{\Gamma })+{\displaystyle \underset{\gamma \mathrm{\Gamma };\gamma \mathrm{\Gamma }}{}}G_{\widehat{P}}(\gamma )F([\mathrm{\Gamma }/\gamma ])\right)=`$
$`=𝒫_{}\left(F(\mathrm{\Gamma })+{\displaystyle \underset{n=1}{}}\left({\displaystyle \underset{\gamma _1,\mathrm{},\gamma _n^{}\mathrm{\Gamma }}{}}\stackrel{~}{P}(\gamma _1)\mathrm{}\stackrel{~}{P}(\gamma _n)F\left([\mathrm{\Gamma }/(\gamma _1\mathrm{}\gamma _n)]\right)\right)\right)`$ (9.5)
Here $`\gamma ^{}\mathrm{\Gamma }`$ means that the sum goes over non-intersecting box-subgraphs $`\gamma _1,\mathrm{},\gamma _n\mathrm{\Gamma }`$. We also assumed that (9.3) is in fact substituted by a more sophisticated constraint
$`𝒫_+\left(\stackrel{~}{P}(\mathrm{\Gamma })F([\mathrm{\Gamma }/\mathrm{\Gamma }])\right)=0`$ (9.6)
Then we can omit $`𝒫_{}`$ acting on the l.h.s. of (9.5). Eq.(9.5) provides a recursion formula for $`\stackrel{~}{P}(\mathrm{\Gamma })`$ if $`F([\mathrm{\Gamma }/\mathrm{\Gamma }])0`$ whenever the r.h.s. of (9.5) is non-vanishing. In fact, this is a necessary requirement for renormalizability of the theory (of a particular reduction of the universal model (2.1)): all the elementary vertices $`[\mathrm{\Gamma }/\mathrm{\Gamma }]`$ should be included into the bare Lagrangian, if they have the structure which can be generated in perturbation theory with $`𝒫`$-negative coefficients (in traditional language of quantum field theory: if there are divergent diagrams with a given number of external legs and external-momenta dependence, an elementary vertex with such valence and momentum dependence should be included into the bare Lagrangian).
Recursive formula (9.5) has a formal solution in terms of its own forest formula, involving decorated rooted trees. For connected graph $`\mathrm{\Gamma }`$ consider a sequence of embedded box-subgraphs, complementary to $`\{\gamma \}_n`$ in (8.21), $`\{\{\gamma \}\}_n:\gamma _0=\mathrm{\Gamma },\gamma _1\mathrm{\Gamma },\gamma _2\gamma _1\mathrm{\Gamma },\mathrm{},\gamma _n\gamma _{n1}\mathrm{}\gamma _1\mathrm{\Gamma }`$. It corresponds to a collection of non-intersecting boxes, which can now (in variance with the set, used in the definition of particular box-subgraph) lie one inside another. Such collection allows one to build a decorated rooted tree $`𝒯`$ . If $`\mathrm{\Gamma }`$ is disconnected, there will be trees, associated with every connected component. With lower site of each box one associates a vertex of the tree, two vertices are connected by a link if one of the corresponding boxes lies immediately inside another (i.e. there are no boxes in between the two). The root link ends at a vertex, associated with $`\gamma _0=\mathrm{\Gamma }`$. According to this construction, every vertex of the tree is associated with connected box-subgraph $`\widehat{\gamma }_k\gamma _k`$ ($`\gamma _k`$ need not be connected), and there is exactly one link, going downwards (towards the root, i.e. associated with the neighbor bigger box) and connecting $`\widehat{\gamma }_k`$ to some $`\widehat{\gamma }_{k1}`$, and unrestricted number of links, going upwards and connecting $`\widehat{\gamma }_k`$ to some collection $`\widehat{\gamma }_{k+1}^1,\mathrm{},\widehat{\gamma }_{k+1}^{s(k)}\gamma _{k+1}`$. The solution to (9.5) associates with every vertex $`\widehat{\gamma }_k`$ an operator $`\left(F([\widehat{\gamma }_k/\widehat{\gamma }_k])\right)^1(𝒫_{})F([\widehat{\gamma }_k/(\widehat{\gamma }_{k+1}^1\mathrm{}\widehat{\gamma }_{k+1}^{s(k)})])`$, where projector $`𝒫_{}`$ acts upwards along the branches of the tree. The root vertex $`\widehat{\gamma }_0`$ (i.e. a connected component of $`\mathrm{\Gamma }`$) contributes just $`F([\widehat{\gamma }_0/(\widehat{\gamma }_1^1\mathrm{}\widehat{\gamma }_1^{s(0)})])`$. In these terms the result of $`R`$-operation can be written as follows :
$`(G_{\widehat{P}}F)(\mathrm{\Gamma })=`$
$`={\displaystyle \underset{_\mathrm{\Gamma }}{}}{\displaystyle \underset{𝒯_\mathrm{\Gamma }}{}}({\displaystyle \underset{\stackrel{vertices}{of𝒯}}{\overset{}{}}}{\displaystyle \frac{1}{F([\widehat{\gamma }_k/\widehat{\gamma }_k])}}(𝒫_{})F([\widehat{\gamma }_k/(\widehat{\gamma }_{k+1}^1\mathrm{}\widehat{\gamma }_{k+1}^{s(k)})])`$ (9.7)
Arrow over the product sign means that the product is ordered along the branches.
Importance of Bogolubov’s recursion in the space of function $`F\{T\}`$ is that it converts partition functions ($`\tau `$-functions) into partition functions, while arbitrary subtraction procedure, like the naive $`F_\mathrm{\Gamma }\{T\}𝒫_{}(F_\mathrm{\Gamma }\{T\})`$, does not have this property: it may not be represented as an action of $`\mathrm{𝐷𝑖𝑓𝑓}`$ and no operator $`g_{T+\stackrel{~}{P}(T)}`$ results from such a subtraction.
It deserves noting that $`\left(F\{T+\stackrel{~}{P}(T)\}\right)_\mathrm{\Gamma }F_\mathrm{\Gamma }\{T+\stackrel{~}{P}(T)\}`$. For example, for the simplest chain graph $`C_1`$ (one valence-two vertex)
$`Z_{C_1}^{ij}\{T+\stackrel{~}{P}(T)\}=T_{(2)}^{ij}+\stackrel{~}{P}_{(2)}^{ij}\{T\}=`$
$`=Z_{C_1}^{ij}\{T\}+\stackrel{~}{P}_{C_1}^{ij}\{T\}+\stackrel{~}{P}_{C_2}^{ij}\{T\}+\mathrm{},`$ (9.8)
while
$`\left(Z^{ij}\{T+\stackrel{~}{P}(T)\}\right)_{C_1}=Z_{C_1}^{ij}\{T\}+\stackrel{~}{P}_{C_1}^{ij}\{T\}`$ (9.9)
Because of this difference one sometime says that renormalization of Lagrangian does not make contribution of each individual graph $`𝒫`$-positive (in the sense that sometime $`𝒫_{}\left(Z_\mathrm{\Gamma }\{T+\stackrel{~}{P}(T)\}\right)0`$), while $`R`$-operation does (in the sense that always $`𝒫_{}\left(Z\{T+\stackrel{~}{P}(T)\}\right)_\mathrm{\Gamma }=0`$). However, as we just explained, if interpreted properly, renormalization of Lagrangian and $`R`$-operation are just the same.
From here on – if one wants to continue – one needs to split the universal model (2.1) into smaller universality classes, which differ by the choice and properties of the sets $`I`$ (where indices $`i`$ in (2.1) take values), especially by the ways the possibly-divergent sums over indices (e.g. integrals over momenta) are regularized (it still makes sense to keep the full set of coupling constants $`T_{i_1\mathrm{}i_n}^{(n)}`$). The most interesting projectors exploit particular properties of particular $`I`$’s. They can act non-trivially on the basic functions $`Z_\mathrm{\Gamma }\{T\}`$, not only on the coefficient functions $`F(\mathrm{\Gamma })`$ (this actually happens in the case of regularized continuous field models, at least in the naive approach). For particular projectors the counter-terms $`\stackrel{~}{P}(\mathrm{\Gamma })`$ can vanish for certain classes of graphs (for divergency-eliminating projectors in renormalizable field models contributing are only graphs with loops and restricted number of external legs). Given $`I`$ and $`𝒫_\pm `$, one can say that the $`R`$-operation (9.2) provides a full set of $`𝒫`$-positive functions on $`(I)`$: a linear basis is provided by the set of $`Z_\mathrm{\Gamma }\{T+\stackrel{~}{P}(T)\}`$ (generically this space is smaller than the one with the basis $`Z_\mathrm{\Gamma }\{T\}`$).
## 10 Conclusion
We described the relation between the algebraic structures, introduced by A.Connes and D.Kreimer, and the generic bilinear relations (Hirota equations) for effective actions in quantum field theory. We discussed two groups acting on the moduli space $``$ of theories: one, essentially commutative $`\mathrm{𝑆ℎ𝑖𝑓𝑡}`$, acting transitively on $``$ and responsible for bilinear relations; another, the non-commutative stability subgroup of the Gaussian point $`\mathrm{𝐷𝑖𝑓𝑓}_{\mathrm{}}`$ in the diffeomorphism group $`\mathrm{𝐷𝑖𝑓𝑓}`$, is associated with Lie algebra of vector fields on $``$, it is related to the CK Hopf algebra of graphs, to Bogolubov’s $`R`$-operation and to renormalization group flows. Bogolubov’s $`R`$-operation is defined in terms of projector operators and can be expressed as renormalization of the action ($`T`$-dependent shift of the coupling constants $`T`$). This study provides a long awaited support to the idea of hidden integrability of non-perturbative quantum phenomena from the field of conventional field theory (Feynman diagram technique).
It also opens a way for the study of analogous phenomena in perturbative string theory, where graphs are substituted by open Riemann surfaces and CK Hopf algebra has interesting generalizations (an infinitesimal deformation in that direction is to the Hopf algebra of fat graphs, associated with the universal matrix model (2.5).
The old belief that the moduli space $``$ of theories and diffeomorphism group $`\mathrm{𝐷𝑖𝑓𝑓}`$ are indeed very similar to conventional simple moduli spaces, studied in mathematics and elementary string theory, gains a new support from the observations in earlier papers of D.Kreimer . However, this subject is beyond the scope of the present paper.
## 11 Acknowledgments
We are indebted for discussions to A.Gorsky, A.Levin, A.Losev and V.Novikov.
Our work is partly supported by RFBR grants 98-01-00328 (A.G.), 98-02-16575 (A.M.), by INTAS grant 97-0103 (K.S.), by Scientific Schools Grant 00-15-96557 (A.G.) and by the Russian President’s Grant 00-15-99296 (A.M.).
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# Quasi-one-dimensional spin chains in CuSiO3: an EPR study
## Abstract
Temperature dependent EPR studies were performed on a single crystal of CuSiO<sub>3</sub>. This recently discovered compound is isostructural with the spin-Peierls compound CuGeO<sub>3</sub>. The EPR signals show different characteristics than those of CuGeO<sub>3</sub> and are due to Cu<sup>2+</sup> spins located along quasi one-dimensional chains. For $`T>8.2`$ K the spin susceptibility closely follows the predictions of a $`S=1/2`$ one-dimensional Heisenberg antiferromagnet with $`J/k_B=21`$ K. Below $`T=8.2`$ K the spin susceptibility immediately drops to zero indicating long range magnetic order.
The linear spin-chain system CuGeO<sub>3</sub> is the first and up to now the only inorganic compound that exhibits a spin-Peierls transition. Regarding the magnetic properties the partial substitution of Ge by Si was an important subject in terms of studying frustration effects and the coexistence of the spin-Peierls state with long-range antiferromagnetic order. To characterize the nature of antiferromagnetic interactions in Si doped and pure CuGeO<sub>3</sub> electron paramagnetic resonance (EPR) of Cu<sup>2+</sup> spins provided important results. In pure CuGeO<sub>3</sub> the EPR parameters differ from those of conventional one-dimensional Heisenberg antiferromagnets. The antisymmetric Dzyaloshinsky-Moriya (DM) exchange interaction was claimed to explain this difference. In Si doped CuGeO<sub>3</sub> coexistence of spin-Peierls and antiferromagnetic order is reported for Si concentrations below $`1\%`$. For higher Si concentrations (up to 50%) a long range antiferromagnetic ground state is observed. However, for $`T>15`$ K the temperature dependence of the EPR parameters does not change significantly for Si-doping concentrations up to 7%. This paper reports the first EPR results on pure CuSiO<sub>3</sub> which are very different from pure and slightly Si-doped CuGeO<sub>3</sub>.
The EPR measurements were performed at X-band frequency with a Bruker ELEXSYS spectrometer. For cooling a continuous-flow helium cryostat (Oxford) was used. The spectra were taken from a non-oriented single crystal of CuSiO<sub>3</sub> at temperatures between 4 K and 300 K. DC magnetization measurements at low fields $`H10`$ kOe were carried out on a commercial SQUID magnetometer. The single crystal of CuSiO<sub>3</sub> was synthesized by dehydration of the mineral dioptase. Reported EPR spectra of CuO did not show up in our EPR spectra which indicated the high quality of our crystal. Figure 1 shows the temperature dependence of the EPR linewidth and EPR $`g`$ factor (determined from the EPR resonance field) of the investigated CuSiO<sub>3</sub> crystal. The inset of Fig. 1 shows a representative EPR spectrum of CuSiO<sub>3</sub> at $`T=40`$ K (solid line, derivative of the EPR absorbed power). At all temperatures the spectra could be nicely fitted with a Lorentzian derivative (dashed line in the inset of Fig. 1).
The EPR linewidth linearly decreases down to $`T100`$ K at a rate 0.5 Oe/K. This contrasts to a much steeper and non-linear decrease of the linewidth in CuGeO<sub>3</sub> which the antisymmetric DM interaction was used for to explain the linewidth. Therefore in CuSiO<sub>3</sub> a DM interaction seems to be less important. Anisotropic exchange interactions also contribute to the EPR linewidth. As the deviation of the O(2)-Cu-O(2) angle from $`90^{}`$ is smaller in CuSiO<sub>3</sub> than in CuGeO<sub>3</sub> the anisotropic exchange interactions in both compounds should be different. For temperatures above $`T12`$ K the EPR $`g`$ factor has a nearly temperature independent value of $`g=2.156\pm 0.001`$ which is commonly observed for Cu<sup>2+</sup>-ions in an octahedral environment. This is consistent with the crystal structure of CuSiO<sub>3</sub> which is reported to be the same as in CuGeO<sub>3</sub> (orthorhombic, Pbmm) and where the Cu<sup>2+</sup>-ions are located within strongly elongated oxygen octahedra.
Figure 2 shows the temperature dependence of the EPR intensity $`I_{EPR}(T)`$ which is determined by integration of the spectra. $`I_{EPR}(T)`$ is proportional to the spin susceptibility of Cu<sup>2+</sup> and can be well compared with the magnetic susceptibility $`\chi (T)`$ as shown in the inset of Fig. 2. However, below $`T=8.2`$4 K the EPR intensity reduces rapidly to zero, indicating an ordering phenomenon rather than a spin-Peierls state which produces an exponential decrease of the intensity. This is also evidenced by the none vanishing $`\chi (T0)`$ which is usually due to an anisotropic antiferromagnetic state. Above $`T=8.2`$ K $`\chi (T)`$ and $`I_{EPR}(T)`$ are very well described by theoretical calculations for an $`S=1/2`$ one-dimensional Heisenberg antiferromagnet (1D HAF) without frustration effects. This leads to an Cu-O(2)-Cu exchange of $`J/k_B=21`$ K which is much smaller than in CuGeO<sub>3</sub> ($`J/k_B160`$ K) as can be expected from the smaller Cu-Cu distances and the smaller O(2)-Cu-O(2) in CuSiO<sub>3</sub> . For high temperatures $`\chi (T)`$ and $`I_{EPR}(T)`$ nicely follow a Curie-Weiss law with a Weiss-temperature of $`\mathrm{\Theta }=7.2`$ K, indicating weak antiferromagnetic coupling.
Figure 3 shows a characterization of the temperature dependence of the EPR linewidth. The high temperature part is estimated with a linear function $`\mathrm{\Delta }H_{lin}(T)=0.5T\mathrm{Oe}/\mathrm{K}+300`$ Oe. This linear part was subtracted from $`\mathrm{\Delta }H`$ in order to obtain the broadening $`\mathrm{\Delta }H_{crit}`$ when the temperature is lowered towards a critical temperature $`T_{crit}=7.5`$ K. A power law $`\mathrm{\Delta }H_{crit}(TT_{crit})^\alpha `$ approximately describes the linewidth with $`\alpha =0.25`$ at low temperatures and above $`T=8.2`$ K. However, at $`T8.2`$ K the type of broadening obviously changes as a noticeable deviation from a power law occurs. This is indicated by the short dashed line in Fig. 3. The linewidth strongly increases nearly below the same temperature ($`T8.2`$ K) where a strong increase of the $`g`$ factor is observed as well (see Fig. 1). Hence the change of line broadening is indicative for the onset of magnetic ordering which yields strong internal fields and therefore inhomogeneous line broadening effects. Measurements of the specific heat give strong evidence for long range antiferromagnetic order. From the critical behavior of the linewidth it is not possible to compare CuSiO<sub>3</sub> unambiguously with typical antiferromagnets neither for the one-dimensional case like CuCl$`{}_{2}{}^{}`$2NC<sub>5</sub>H<sub>5</sub> ($`\alpha =0.5`$) nor for the three-dimensional case like GdB<sub>6</sub> ($`\alpha =1.5`$).
In summary our EPR results on CuSiO<sub>3</sub> do not show evidences for a spin-Peierls state below $`T=8.2`$ K. For low temperatures the EPR intensity and EPR linewidth are rather explained by long-range magnetic ordering phenomena. For temperatures above $`T=8.2`$ K the EPR intensity is proportional to the magnetic susceptibility and can be reproduced well with a behavior of an one-dimensional antiferromagnet. Antisymmetric and anisotropic exchange interactions contribute differently to the EPR parameters in CuSiO<sub>3</sub> and CuGeO<sub>3</sub>. Further clarification should be provided by measurements at higher temperatures and at defined crystal orientations which presently are in progress.
We acknowledge fruitful discussions with H.-A. Krug von Nidda and support by SFB 484.
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# On billiard weak solutions of nonlinear PDE’s and Toda flows 11footnote 1PACS numbers 03.40.Gc, 11.10.Ef, 68.10.-m, AMS Subject Classification 58F07, 70H99, 76B15
## 1 Introduction
Camassa and Holm introduced and studied classes of soliton-type weak solutions, for an integrable nonlinear equation derived in the context of a shallow water model. In particular, they described the soliton dynamics in terms of a system of Hamiltonian equations for the locations of the “peaks” of the solution, the points at which its spatial derivative changes sign. In other words, each peaked solution, or peakon, can be associated with a mechanical system of moving particles. New systems of this type were obtained in Calogero and Calogero and Francoise . The $`r`$-matrix approach was applied to the Lax pair formulation of an $`n`$-peakon system by Ragnisco and Bruschi , who also pointed out the connection of this system with the classical Toda lattice. A discrete version of the Adler-Kostant-Symes factorization method was used by Suris to study a discretization of the peakons lattice, realized as a discrete integrable system on a certain Poisson submanifold of $`\mathrm{gl}(n)`$ equipped with $`r`$-matrix Poisson bracket.
In Alber et al. the existence of peakons was linked to the presence of poles in the energy dependent Schrődinger operators associated with integrable evolution equations. Namely, it was shown that the presence of a pole in the potential is essential for a special limiting procedure which allows for the formation of “billiard’ weak solutions. By using the algebraic-geometric method, these billiard solutions are related to finite dimensional integrable dynamical systems with reflections. In this way, profiles of billiard weak solutions are associated with billiard motion inside quadrics with and without the presence of a Hooke’s potential. The points of impact on the quadrics correspond to the peaks of the profiles of weak solutions of the nonlinear PDE. Billiard solutions include new quasi-periodic and soliton-like solutions, as well as peaked solitons with compact support. This method can be used for a number of equations including the shallow water equation, the Dym type equation, as well as $`N`$-component systems with poles, together with all the equations in their respective hierarchies.
In Section 2 of this paper we derive equations for the motion of the peaks of billiard solutions in the context of the algebraic-geometric approach. We construct solutions of these equations by connecting them to a flow on an appropriate Riemann surface which leads to a classical Jacobi inversion problem.
In Section 3 we show that orbits of the Toda flow in $`sl(n)`$ into which the peakons lattice is mapped by means of the Lax matrix introduced in Ragnisco and Bruschi admits a natural generalization to the case of an arbitrary Lie algebra $`𝔤`$. Restrictions of the Toda flow in $`𝔤`$ to these orbits can be viewed as generalized peakons lattices. Using the methods developed in Gekhtman and Shapiro we give an intrinsic description of these orbits and prove complete integrability of the corresponding flows. Unlike in the $`sl(n)`$ case, Chevalley invariants of $`𝔤`$ are not sufficient to ensure integrability in the general situation, and so we produce the requisite number of additional first integrals. Lastly, we explain how to construct Darboux coordinates for generalized peakon orbits and consider $`sl(n)`$ and $`G_2`$ as examples.
## 2 The Dynamics of the Peak Points
As shown in Alber et al. , quasi-periodic solutions of the shallow water equation
$$U_t+3UU_x=U_{xxt}+2U_xU_{xx}+UU_{xxx},$$
(2.1)
on the infinite line can be represented in the form
$$U(x,t)=\mu _1+\mathrm{}+\mu _n𝔪$$
(2.2)
where $`\mu `$’s are solutions of the following systems of equations in $`x`$ and $`t`$,
$$\mu _i^{}=\frac{\mu _i}{x}=\mathrm{Sign}(\mu _i,m_{2i1},m_{2i},s_i)\frac{\sqrt{R(\mu _i)}}{\mu _i_{ji}^n(\mu _i\mu _j)}$$
(2.3)
and
$$\dot{\mu }_i=\frac{\mu _i}{t}=\mathrm{Sign}(\mu _i,m_{2i1},m_{2i},d_i)\frac{(\mu _i\mathrm{\Sigma })\sqrt{R(\mu _i)}}{\mu _i_{ji}^n(\mu _i\mu _j)},i=1,\mathrm{},n,\mathrm{\Sigma }=\mu _1+\mathrm{}+\mu _n,$$
(2.4)
with
$$R(\mu )=\mu \underset{r=1}{\overset{2n+1}{}}(\mu m_r).$$
(2.5)
Here the constants $`m_r`$ depend on the initial condition for the equation (2.1), $`𝔪=_r^{2n+1}m_r`$, and $`\mathrm{Sign}(y,l,r,s)`$ is a function on the Riemann surface, $`\mathrm{\Gamma }=\{w^2=R(\mu )\},`$ which switches from $`1`$ to $`1`$ and back each time $`y`$ reaches $`l`$ or $`r`$, endpoints of a cut on the Riemann surface. These ordinary differential equations (ODE’s) for the $`\mu `$-variables provide half of the equations of a finite dimensional Hamiltonian system. Notice that periodic solutions form a subclass of the quasi-periodic solutions. Also, soliton solutions can be obtained from quasi-periodic solutions by shrinking pair wize the roots $`m_r`$ of the spectral polynomial $`R(E)`$.
Peaked quasi-periodic solutions of equation (2.1) can be constructed from solutions of (2.3) and (2.4) by using the limiting process $`m_10`$ (see Alber et al. ) and the trace formula (2.2). In what follows ODE’s governing the time evolution of the peak locations are obtained and a connection with the canonical variables of Camassa and Holm is described.
The $`x`$ and $`t`$ evolution of the $`\mu `$ variables can also be connected to a constrained motion of a particle on an $`n`$-dimensional hypersurface imbedded in $`^{n+1}`$ which is parameterized by the constants $`m_r`$, under the action of a harmonic force field (see Alber et al. ). The limit $`m_10`$ corresponds to “flattening” this surface along one direction which results in harmonically forced motion in a region of $`^n`$ confined by an $`(n1)`$-dimensional boundary. The particle collides with the boundary and moves along segments of an $`n`$-dimensional billiard. The reflection on the boundary causes, after using the trace formula (2.2), the appearance of peaks in the corresponding PDE solutions of equation (2.1).
On the level of (2.3) and (2.4), the boundary of the billiard can be shown to be described by the following condition: $`\mu _1(x,t)=0`$. Reflection on the boundary is described by the sign switch from 1 to -1 (or from -1 to 1) of first derivative of $`\mu _1`$ in (2.3) and (2.4).
### The Periodic Case.
The main steps in deriving equations for the peak locations can be demonstrated by considering the following basic example. Let us assume $`n=2`$, and $`m_1=0`$. Equation $`\mu _1(x,t)=0`$ defines a function of time $`x=q(t)`$. An equation for the time dependence of $`q(t)`$ can be found by differentiating $`(\mu _1(q(t),t)=0)`$ with respect to time. We have
$$\frac{d}{dt}\left(\mu _1(q(t),t)\right)=0=\frac{\mu _1}{x}(q(t),t)\dot{q}+\frac{\mu _1}{t}(q(t),t).$$
(2.6)
The commuting $`x`$\- and $`t`$-flows of $`\mu _1`$ are described by
$`{\displaystyle \frac{\mu _1}{x}}=\mathrm{Sign}(\mu _1){\displaystyle \frac{\sqrt{C_4(\mu _1)}}{\mu _1\mu _2}},`$ (2.7)
$`{\displaystyle \frac{\mu _1}{t}}=\mathrm{Sign}(\mu _1){\displaystyle \frac{(\mu _2)\sqrt{C_4(\mu _1)}}{\mu _1\mu _2}},`$
$$C_4(\mu )=(\mu m_2)(\mu m_3)(\mu m_4)(\mu m_5).$$
Notice that the spatial derivative of $`\mu _1`$ is not defined at $`x=q(t)`$, so that in the above formula $`\mu _1/x`$ evaluated at $`(q(t),t)`$ needs to be found in the limit as $`xq`$ either from the left, or from the right. The sign uncertainty connected with this choice factors out from the final formula for $`\dot{q}`$. This combined with (2.6) and the billiard condition at the boundary $`\mu _1(q(t),t)=0`$ gives
$$\dot{q}=\mu _2(q(t),t).$$
(2.8)
We also need an evolution equation for the quantity $`\mu _2(q(t),t)`$. Notice that from the trace formula (2.2) and the expression for the boundary $`\mu _1(q(t),t)=0`$, this is the quantity that determines the amplitude of the peak $`U(q(t),t)`$. The time derivative of $`\mu _2(q(t),t)`$ is
$$\frac{d}{dt}\left(\mu _2(q(t),t)\right)=\frac{\mu _2}{x}(q(t),t)\dot{q}+\frac{\mu _2}{t}(q(t),t),$$
(2.9)
which used with the commuting $`x`$\- and $`t`$-flows equations for $`\mu _2`$ yields, in analogy with system (2.8),
$`{\displaystyle \frac{dy}{dt}}(t)`$ $`=\mathrm{Sign}(y){\displaystyle \frac{\sqrt{C_4(y)}}{y}}\left(\dot{q}+B_1(y)\right)`$
$`=\mathrm{Sign}(y){\displaystyle \frac{\sqrt{C_4(y)}}{y}}\dot{q}.`$ (2.10)
Here we have introduced the notation
$$y(t)=\mu _2(q(t),t)$$
and used again the definition of $`q(t)`$ so that
$$B_1(\mu _2(q(t))B_1(y)=\mu _1(q(t),t)=0.$$
Combining this equation for $`\dot{y}`$ with the evolution equation (2.8) for $`q`$ gives the following system for the variables $`y`$ and $`q`$,
$`\dot{y}=\mathrm{Sign}(y)\sqrt{(ym_2)(ym_3)(ym_4)(ym_5)}`$
$`\dot{q}=y`$ (2.11)
Notice that the evolution for $`y`$ is decoupled from that of $`q`$. The inversion problem corresponding to the equation for $`y`$ determines a periodic function of $`t`$ with nonzero average. The evolution of $`q`$ is then given by the combination of a linear growth in time with slope given by the average of $`y`$, plus periodic oscillations. We also remark that for weak solutions of equation (2.1) having isolated discontinuities in the first derivative, the jump condition yields
$$\frac{dq}{dt}=U(q(t),t),$$
and using the trace formula (2.2) this can be seen to coincide with the second equation in system (2).
### The Genus 2 Quasi-periodic Case.
We describe a billiard genus 2 quasi-periodic solution. The spectral polynomial in this case is of the $`5`$th order and will be denoted as $`C_5(\mu )`$. The boundary is again described by $`\mu _1(q(t),t)=0`$. We shall write
$$y_1=\mu _2(q(t),t)\text{and}y_2=\mu _3(q(t),t).$$
From the weak solution approach it follows that $`\dot{q}=U(q_1)=y_1+y_2`$. We differentiate $`y_1`$ and $`y_2`$ to obtain
$$\begin{array}{c}\dot{y}_1=\frac{d}{dt}\left[\mu _2(q,t)\right]=\frac{\mu _2}{x}\dot{q}+\frac{\mu _2}{t}=\frac{\mathrm{Sign}(\mu _2)\sqrt{C_6(\mu _2)}}{\mu _2(\mu _2\mu _3)}(\dot{q}\mu _3)=\frac{\mathrm{Sign}(y_1)\sqrt{C_6(y_1)}}{y_1y_2}\hfill \\ \\ \dot{y}_2=\frac{d}{dt}\left[\mu _3(q,t)\right]=\frac{\mu _3}{x}\dot{q}+\frac{\mu _3}{t}=\frac{\mathrm{Sign}(\mu _3)\sqrt{C_6(\mu _3)}}{\mu _3(\mu _3\mu _2)}(\dot{q}\mu _2)=\frac{\mathrm{Sign}(y_2)\sqrt{C_6(y_2)}}{y_2y_1}.\hfill \end{array}\}$$
(2.12)
This is a well-defined system associated with a genus $`2`$ Riemann surface. The corresponding problem of inversion involves only holomorphic differentials and therefore the quantities $`y_1`$,$`y_2`$ and $`q(t)`$ can be expressed in terms of standard $`\theta `$-functions on the Riemann surface.
### Peakon solutions.
We now specialize the above formalism to the limiting case of soliton solutions of equation (2.1) on the real line, in which each pair $`m_{2i},m_{2i+1}`$ is taken to limit to the $`a_i`$, $`i=1,\mathrm{},n`$. For just one $`\mu `$-variable, using the trace formula $`U=\mu a`$ at the billiard boundary results in $`U(q(t),t)=a`$. The single peakon solution determined by this procedure is therefore
$$U(x,t)=ae^{|x+at|},$$
(2.13)
which is a traveling wave soliton-type solution. Camassa and Holm found
$$U(x,t)=p(t)e^{|xq(t)|},$$
(2.14)
with an additional link between $`p(t)`$ and $`q(t)`$ (Hamiltonian structure) which in the $`1`$-peakon case results in $`p(t)`$ being a constant and $`q(t)`$ being a linear function.
### 2-peakon solution.
In this case, one obtains the following system that describes $`2`$-peakon profiles,
$$\begin{array}{c}\frac{\mu _1}{x}=\mathrm{Sign}(\mu _1)\frac{(\mu _1a_1)(\mu _1a_2)}{(\mu _1\mu _2)}\hfill \\ \\ \frac{\mu _2}{x}=\mathrm{Sign}(\mu _2)\frac{(\mu _2a_1)(\mu _2a_2)}{(\mu _2\mu _1)}\hfill \end{array}\}$$
(2.15)
where $`\mu _1`$ and $`\mu _2`$ are evaluated between $`a_1`$ and $`0`$ and $`a_2`$ and $`0`$, respectively. The corresponding time-flow is
$$\begin{array}{c}\frac{\mu _1}{t}=\mathrm{Sign}(\mu _1)B_1(\mu _1)\frac{(\mu _1a_1)(\mu _1a_2)}{(\mu _1\mu _2)}\hfill \\ \\ \frac{\mu _2}{t}=\mathrm{Sign}(\mu _2)B_1(\mu _2)\frac{(\mu _2a_1)(\mu _2a_2)}{(\mu _2\mu _1)}\hfill \end{array}\}$$
(2.16)
and the first order polynomial $`B_1`$ for equation (2.1) gives
$$B_1(\mu _1)=\mu _2+a_1+a_2,B_1(\mu _2)=\mu _1+a_1+a_2.$$
Define the boundary by introducing functions $`q_1(t)`$ and $`q_2(t)`$ such that
$$\mu _1(q_1(t),t)=0,\mu _2(q_2(t),t)=0.$$
(2.17)
Define functions
$$y_1=\mu _2(q_1(t),t),y_2=\mu _1(q_2(t),t).$$
(2.18)
Differentiating $`\mu _1(q_1(t),t)`$ and $`\mu _2(q_2(t),t)`$ in (2.17) results in
$$\begin{array}{c}\frac{d\mu _1}{dt}=0=\frac{\mu _1}{x}\left(\frac{dq_1}{dt}+B_1(\mu _1)\right)=\frac{\mu _1}{x}\left(\frac{dq_1}{dt}y_1+a_1+a_2\right)\hfill \\ \\ \frac{d\mu _2}{dt}=0=\frac{\mu _2}{x}\left(\frac{dq_2}{dt}+B_1(\mu _2)\right)=\frac{\mu _2}{x}\left(\frac{dq_2}{dt}y_2+a_1+a_2\right)\hfill \end{array}\}$$
(2.19)
which coincides with the jump conditions for weak solutions
$$\begin{array}{c}\frac{dq_1}{dt}=U(q_1)=\left(\mu _1+\mu _2a_1a_2\right)|_{q_1}=\mu _2(q_1)a_1a_2=y_1a_1a_2\hfill \\ \\ \frac{dq_2}{dt}=U(q_2)=\left(\mu _1+\mu _2a_1a_2\right)|_{q_2}=\mu _1(q_2)a_1a_2=y_2a_1a_2.\hfill \end{array}\}$$
(2.20)
Finally, differentiate $`y_1`$ and $`y_2`$ to find
$$\begin{array}{c}\frac{dy_1}{dt}=\frac{\mu _2}{x}\frac{dq_1}{dt}+\frac{\mu _2}{t}=\frac{\mu _2}{x}\left(\frac{dq_1}{dt}+B_1(\mu _2)\right)\hfill \\ \\ \frac{dy_2}{dt}=\frac{\mu _1}{x}\frac{dq_2}{dt}+\frac{\mu _1}{t}=\frac{\mu _1}{x}\left(\frac{dq_2}{dt}+B_1(\mu _1)\right)\hfill \end{array}\}$$
and so
$$\begin{array}{c}\frac{dy_1}{dt}=\mathrm{Sign}(y_1)(y_1a_1)(y_1a_2)\hfill \\ \\ \frac{dy_2}{dt}=\mathrm{Sign}(y_2)(y_2a_1)(y_2a_2).\hfill \end{array}\}$$
(2.21)
Thus, the equations of evolution for $`y_1`$ and $`y_2`$ decouple from those of $`q_1`$ and $`q_2`$. The decoupled equations (2.21) can be solved first and $`q_1`$, $`q_2`$ subsequently determined from (2.20) by quadratures.
It is interesting to examine the connection of the set of variables $`q_i,y_i`$, $`i=1,2`$, with the $`q_i,p_i`$, $`i=1,2`$, introduced by Camassa and Holm . By definition, the $`q`$’s are the same in both sets (the $`(q,y)`$ system does have a canonical Hamiltonian form). As to the $`y`$’s, notice that
$$\begin{array}{c}U(q_1)=y_1a_1a_2=p_1+p_2e^{|q_1q_2|}\hfill \\ \\ U(q_2)=y_2a_1a_2=p_2+p_1e^{|q_1q_2|}\hfill \end{array}\}$$
(2.22)
which provides a definition of the transformation of the $`y`$ variables in terms of $`p`$’s and $`q`$’s. This together with (2.20) yields the first set of equations for the Hamiltonian system derived by Camassa and Holm ,
$$\begin{array}{c}\frac{dq_1}{dt}=p_1+p_2e^{|q_1q_2|}\hfill \\ \\ \frac{dq_2}{dt}=p_2+p_1e^{|q_1q_2|}.\hfill \end{array}\}$$
(2.23)
The constants of motion $`a_1`$ and $`a_2`$ can be expressed in terms of $`p`$’s and $`q`$’s via the first integrals
$$P_{12}=p_1+p_2=(a_1+a_2),$$
and
$$H_{12}=\frac{1}{2}(p_1^2+p_2^2)+p_1p_2e^{|q_1q_2|}=\frac{1}{2}(a_1^2+a_2^2),$$
which are the total momentum and Hamiltonian for the $`(q,p)`$-flow, respectively. Using these expressions, the constants of motion $`a_1`$ and $`a_2`$ can be eliminated from (2.22) and the explicit variable transformation of variables for $`y_1`$, $`y_2`$ in terms of the $`(q,p)`$ system is
$$\begin{array}{c}y_1=p_2(e^{|q_1q_2|}1)\hfill \\ y_2=p_1(e^{|q_1q_2|}1).\hfill \end{array}\}$$
(2.24)
Differentiating these equations with respect to time, using (2.21), (2.20), and the transformation (2.24) itself, results in a system of equations for $`\dot{p}_1`$ and $`\dot{p}_2`$, which can be solved to yield
$$\begin{array}{c}\frac{dp_1}{dt}=\mathrm{sgn}(q_2q_1)p_1p_2e^{|q_1q_2|}\hfill \\ \\ \frac{dp_2}{dt}=\mathrm{sgn}(q_1q_2)p_2p_1e^{|q_1q_2|}.\hfill \end{array}\}$$
Equations (2.23) and (2.24) are precisely those that follow from $`H_{12}`$ with a canonical Hamiltonian structure in the $`(p,q)`$ variables.
The case of three or more derivative-shock singularities $`x_i,y_i`$, $`i=1,2,\mathrm{},N3`$ proceeds in complete analogy with the case $`N=2`$ above. Once again, the $`q`$-flows decouple from those of the $`y`$’s, however the equations in the system governing the $`y`$-flow are now coupled and it is not immediately obvious that this system is integrable. A closer inspection however reveals that the $`y`$-flow shares the same structure as that of the $`\mu `$-variables flow and is therefore integrable by a similar argument.
## 3 Generalized Peakons Lattices
In this section we discuss the connection between the peakons lattice and the Toda flows. The fact that the peakons lattice can be realized as a special case of the Toda flows in $`gl(n)`$ was discovered in Ragnisco and Bruschi and later used by Suris to introduce a discrete time peakons lattice. Let us review these results, concentrating, for simplicity, on a particular case of the peakons lattice with a Hamiltonian
$$H(p,q)=\frac{1}{2}\underset{i=1}{\overset{n}{}}p_i^2+\underset{1i<jn}{}p_ip_j\mathrm{exp}(q_jq_i),$$
(3.1)
that corresponds to the case when particles $`q_i`$ are ordered as follows: $`q_1>q_2>\mathrm{}>q_n`$. It turns out that (3.1) possesses a Lax representation with the Lax matrix
$$L=L(p,q)=\underset{i=1}{\overset{n}{}}p_iE_{ii}+\underset{1i<jn}{}\sqrt{p_ip_j}\mathrm{exp}\left(\frac{1}{2}(q_jq_i)\right)(E_{ij}+E_{ji}),$$
(3.2)
and an auxiliary matrix equal to the half of the skew-symmetric part of $`L`$, the peakons lattice thus being a restriction of the Toda flow to the set of matrices of the form (3.2). It was shown that this set forms a Poisson submanifold w.r.t. the associated $`r`$-matrix bracket on $`gl(n)`$. Complete integrability of (3.1) is provided by the Chevalley invariants of $`gl(n)`$, i.e. spectral invariants of $`L`$.
In what follows we give a uniform Lie-algebraic procedure which allows us to construct peakon-type orbits in any simple Lie algebra. On each of these orbits the Toda flow is shown to be a completely integrable system. Notice that for proving this we provide additional first integrals which supplement the Chevalley invariants of the algebra. We also indicate a convenient way of constructing Darboux coordinates for orbits under consideration.
Let us first recall the Hamiltonian formalism for the generalized (symmetric) Toda flows (cf. Kostant , Goodman and Wallach , Reyman and Semenov-Tian-Shansky ). Let $`𝔤`$ be the normal real form of a simple Lie algebra of rank $`r`$, $`𝐆`$ a corresponding Lie group and $`𝔥`$ a Cartan subalgebra of $`𝔤`$. We denote by $`\mathrm{\Phi }`$ the root system of $`𝔤`$, and by $`\mathrm{\Phi }^+`$ (resp. $`\mathrm{\Phi }^{}`$) the set of all positive (resp. negative) roots. We also fix a Chevalley basis $`\{e_\alpha ,\alpha \mathrm{\Phi };h_i,i=1,\mathrm{},r\}`$ in $`𝔤`$ . All properties of the root systems and Chevalley bases that we need can be found in Humphreys and Onishchik and Vinberg .
Consider a direct sum decomposition $`𝔤=𝔨+𝔟_+`$, where $`𝔟_+`$ is the upper Borel subalgebra and $`𝔨`$ is the maximal compact subalgebra of $`𝔤`$. The dual space $`𝔟_+^{}`$ of $`𝔟_+`$ can be identified with the space $`S`$ of “symmetric” elements of $`𝔤`$ that have a form
$$L=\underset{j=1}{\overset{r}{}}b_jh_j+\underset{\alpha \mathrm{\Phi }}{}a_\alpha (e_\alpha +e_\alpha ).$$
(3.3)
$`S`$ is an annihilator of $`𝔨`$ w.r.t the Killing form $`,`$.
The pull-back of the Lie-Poisson bracket on $`𝔟_+^{}`$ equips $`S`$ with the Poisson bracket
$$\{f_1,f_2\}_S(L)=\frac{1}{2}L,[\pi _+f_1(L),\pi _+f_2(L)],$$
(3.4)
where gradients are defined w.r.t. the Killing form and $`\pi _+`$ is a projection on $`𝔟_+`$ parallel to $`𝔨`$. Equations of motion of the generalized Toda flow on $`S`$ are generated by a Hamiltonian $`H(L)=\frac{1}{2}L,L`$ and have the Lax form
$$\dot{L}=[L,\frac{1}{2}\pi _+(L)]=[\frac{1}{2}\pi _{}(L),L]$$
(3.5)
where $`\pi _{}=Id\pi _+`$ . The bracket (3.4) is a restriction to $`S`$ of the so called $`r`$-matrix bracket on $`𝔤`$:
$$\{f_1,f_2\}_r(X)=\frac{1}{2}X,[\pi _+f_1(X),\pi _+f_2(X)][\pi _{}f_1(X),\pi _{}f_2(X)].$$
(3.6)
Symplectic leaves of the Poisson manifold $`(S,\{,\}_S)`$ coincide with the orbits of the coadjoint action of the upper Borel subgroup $`𝐁_+`$ of $`𝐆`$:
$$\text{Ad}_b^{}(L)=\pi _S\text{Ad}_{b^1}(L),$$
(3.7)
where $`\pi _S`$ is a projection on $`S`$ parallel to $`𝔟_+`$. On every indecomposable symplectic leaf, restrictions of the Chevalley invariants $`I_1=H,I_2,\mathrm{},I_r`$ of $`𝔤`$ form a family of independent Poisson commuting integrals for the Toda flow. This family is maximal, however, only for a few distinguished orbits, that were classified in Goodman and Wallach and Perelomov and Kamalin
To describe a particular type of coadjoint orbits, which in the $`sl(n)`$ case will be shown to be associated with the peakons lattice, we need the following definitions introduced in Gekhtman and Shapiro . Let $`m`$ be the maximal positive root and
$$h_m=[e_m,e_m].$$
(3.8)
Define
$`𝔤^{}=\text{Ker}_{\text{ad}_{e_m}}\text{Ker}_{\text{ad}_{e_m}}`$
$`F=\text{Span}\{e_\alpha :\alpha \mathrm{\Phi }^{}\text{and}(m,\alpha )0\}`$ (3.9)
$`\stackrel{~}{F}=\text{Span}\{F,h_m\}`$
$`V=Fe_m.`$
Then $`𝔤^{}`$ is a semisimple subalgebra of $`𝔤`$ and $`F`$ is a Heisenberg subalgebra of $`𝔤`$, i. e. $`V`$ is spanned by root vectors $`e_{\alpha _i},e_{m\alpha _i},i=1,\mathrm{}N`$, such that
$`[e_{\alpha _i},e_{\alpha _j}]=[e_{m\alpha _i},e_{m\alpha _j}]=[e_m,e_{\alpha _i}]=[e_m,e_{m\alpha _i}]=0,`$
$`[e_{\alpha _i},e_{m\alpha _j}]=c_i\delta _i^je_m,,`$ (3.10)
where $`c_i`$ are some positive constants. Note also that
$$[h_m,e_{\alpha _i}]=e_{\alpha _i},[h_m,e_{m\alpha _i}]=e_{m\alpha _i},[h_m,e_{\pm m}]=\pm 2e_{\pm m},$$
(3.11)
Therefore (3.4) and (3.11) yield the following Poisson brackets
$$\{y_0,x_0\}_S=x_0,\{y_0,x_i\}_S=\frac{1}{2}x_i,\{y_0,y_i\}_S=\frac{1}{2}y_i,\{y_i,x_i\}_S=\frac{1}{2}x_0,i=1,\mathrm{}N,$$
(3.12)
for the linear functions
$$x_i=L,e_{\alpha _i},y_i=\frac{1}{c_i}L,e_{m+\alpha _i},x_0=L,e_m,y_0=L,h_m.$$
(3.13)
All other brackets between functions (3.13) are zero.
Let $`h_0𝔥`$ be orthogonal to $`h_m`$ w.r.t the Killing form and let $`𝔒_{m,h_0}`$ be the coadjoint orbit of $`𝐁_+`$ through $`(e_m+h_0+e_m)`$.
###### Theorem 3.1
$`𝔒_{m,h_0}`$ can be parameterized by elements of $`\stackrel{~}{F}`$ :
$$𝔒_{m,h_0}=\{L=\pi _S\left(\zeta \frac{1}{2x(\zeta )}[\text{ad}_{e_m}v(\zeta ),v(\zeta )]\right):\zeta \stackrel{~}{F}_{},x(\zeta )>0\},$$
(3.14)
where element $`\zeta \stackrel{~}{F}_{}`$ is decomposed as
$$\zeta =x(\zeta )e_m+y(\zeta )h_m+v(\zeta ),$$
and $`v(\zeta )V.`$
The Toda flow (3.5) is completely integrable on $`𝔒_{m,h_0}`$ with a maximal family of functions in involution given by the coefficients of polynomials
$$I_j\left(L+\frac{\lambda }{x(\zeta )}e_m\right)=\underset{k}{}I_{jk}(L)\lambda ^k,j=1,\mathrm{}r,$$
(3.15)
where $`I_j`$ are the Chevalley invariants of $`𝔤`$.
The proof makes use of the so-called $`1`$-chop map, a Poisson map from $`𝔤`$ to $`𝔤^{}`$, which was introduced for $`sl(n)`$ in Singer and generalized for an arbitrary Lie algebra in Gekhtman and Shapiro . This map and the functions (3.15) play a key role in the proof of the complete integrability of the Toda flows on generic orbits (see Deift et al. and Ercolani et al. for the $`sl(n)`$ case and Gekhtman and Shapiro for a general case).
It also follows from Theorem 2.1 that for constructing Darboux coordinates on $`𝔒_{m,h_0}`$, it is sufficient to construct Darboux coordinates for a rather simple Poisson algebra generated by (3.12).
The following example demonstrates the connection with the peakons lattice.
Example 2.2. Let $`𝔤=sl(n)`$. Then $`e_m=E_{1n},e_m=E_{n1},h_m=E_{11}E_{nn}`$,
$$\stackrel{~}{F}=\{\zeta =\left[\begin{array}{ccc}y& 0& 0\\ v_1& 0& 0\\ x& v_2^T& y\end{array}\right]:x,y,v_1,v_2^{n2}\}.$$
Since Tr$`(L)`$ is invariant under the action (3.7), we can define $`S`$ to be the space of all symmetric matrices, not necessarily with zero trace. One can choose $`h_0`$ to be any diagonal matrix of the form $`h_0=`$diag$`(\kappa ,d,\kappa )`$, where $`d`$ is a $`(n2)\times (n2)`$ diagonal matrix. It follows from (3.14) that
$$𝔒_{m,h_0}=\{L=\left[\begin{array}{ccc}y+\kappa & v_1^T& x\\ v_1& d+\pi _S(x^1v_1v_2^T)& v_2\\ x& v_2^T& \kappa y\end{array}\right]:x>0,y,v_1,v_2^{n2}\}.$$
(3.16)
If $`d=0`$, an open set $`\{\kappa <y<\kappa ;v_{1i}>0,v_{2i}>0,i=1,\mathrm{},n2\}𝔒_{m,h_0}`$ admits Darboux coordinates $`p_i`$ and $`q_i`$ defined as follows:
$$p_1=y+\kappa ,p_n=\kappa y,p_{i+1}=\frac{v_{1i}v_{2i}}{x},i=1,\mathrm{},n2,$$
and
$$\mathrm{exp}(q_nq_1)=\frac{x^2}{\kappa ^2y^2},\mathrm{exp}(q_iq_1)=\frac{xv_{1i}}{(y+\kappa )v_{2i}},i=2,\mathrm{},n1.$$
With this parameterization matrix $`L`$ in (3.16) coincides with the Lax matrix (3.2). Note that $`p_1+\mathrm{}+p_n=const`$ on $`𝔒_{m,h_0}`$ and that coordinates $`q_i`$ are defined only up to a translation. Note that a different choice of Darboux coordinates for $`𝔒_{m,0}`$ was proposed in Kamalin and Perelomov as an example of the general construction of canonical coordinates on coadjoint orbits.
Example 2.3. Let $`𝔤`$ be an exceptional algebra of type $`G_2`$. Denote the short and long simple roots of $`𝔤`$ by $`\nu _1,\nu _2`$ resp., and root vector corresponding to a negative root $`\alpha =(i\nu _1+j\nu _2)`$ by $`e_{ij}`$. Then $`F`$ is spanned by vectors $`e_m=e_{32}`$ and $`e_{i1},i=1,\mathrm{}4`$. Thus, for any $`h_0`$ the orbit $`𝔒_{m,h_0}`$ is $`6`$-dimensional, while there are only two independent Chevalley invariants. This illustrates the necessity of using functions (3.15) to establish complete integrability.
Put $`h_0=0`$ and choose Darboux coordinates $`q_i,p_i`$ for Poisson algebra (3.12) with $`N=3`$ in the form:
$$q_1=x_0,p_1=\frac{y_0}{x_0},q_2=\frac{x_1}{3\sqrt{x_0}},p_2=\frac{6y_1}{\sqrt{x_0}},q_3=\frac{y_2}{\sqrt{x_0}},p_3=\frac{2x_2}{\sqrt{x_0}}.$$
Then Theorem 2.1 gives the following integrable polynomial Hamiltonian quadratic in momenta:
$$H(p,q)=\frac{1}{12}(p_3q_2+6q_3^2)^2+\frac{1}{8}(p_3q_3+p_2q_2)^2+\frac{q_1}{12}(p_3^2+12p_2^2+3p_1^2)+q_1^2+3q_3^2q_1+3q_3^4$$
(3.17)
A detailed analysis of the system generated by (3.17), as well as integrable systems associated with peakon-type orbits in other simple Lie algebras will be given in a forthcoming publication, Alber and Gekhtman .
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# 1 Introduction
## 1 Introduction
The production cross sections for all the processes at hadron-collider experiments are controlled by strong interaction physics and, hence, by its underlying field theory, QCD (see recent overviews in Refs. ). Studies of QCD at the Tevatron and the LHC have two main purposes . First, they are important to test the predictions of QCD, to measure its fundamental parameters (e.g. the strong coupling $`\alpha _\mathrm{S}`$) and to extract quantitative information on its non-perturbative dynamics (e.g. the distribution of partons in the proton). Second, they are relevant to a precise estimate of the background to other Standard Model processes and to signals of new physics.
This contribution is not a comprehensive review of QCD at high-energy hadron colliders. It is based on a selection of the topics presented in my introductory lectures at this Workshop. The selection highlights the QCD subjects that were most discussed during the Workshop and includes a pedagogical overview of some of the corresponding theoretical tools.
After the introduction of the general theoretical framework, I summarize in Sect. 2 the present knowledge on the parton densities and its impact on QCD predictions for hard-scattering processes at the Tevatron and the LHC. In Sect. 3, I then discuss some issues related to processes that are sensitive to the gluon density and, hence, to its determination. Section 4 presents a dictionary of different approaches (fixed-order expansions, resummed calculations, parton showers) to perturbative QCD calculations. The dictionary continues in Sect. 5, where I review soft-gluon resummation and discuss some recent phenomenological applications of threshold resummation to hadron collisions.
The QCD framework to describe any inclusive hard-scattering process,
$$h_1(p_1)+h_2(p_2)H(Q,\{\mathrm{}\})+X,$$
(1)
in hadron–hadron collisions is based on perturbation theory and on the factorization theorem of mass singularities. The corresponding cross section is computed by using the factorization formula
$`\sigma (p_1,p_2;Q,\{\mathrm{}\})`$ $`=`$ $`{\displaystyle \underset{a,b}{}}{\displaystyle _{x_{\mathrm{min}}}^1}𝑑x_1𝑑x_2f_{a/h_1}(x_1,\mu _F^2)f_{b/h_2}(x_2,\mu _F^2)\widehat{\sigma }_{ab}(x_1p_1,x_2p_2;Q,\{\mathrm{}\};\mu _F^2)`$ (2)
$`+`$ $`𝒪\left((\mathrm{\Lambda }_{QCD}/Q)^p\right).`$
The colliding hadrons $`h_1`$ and $`h_2`$ have momenta $`p_1`$ and $`p_2`$, $`H`$ denotes the triggered hard probe (vector bosons, jets, heavy quarks, Higgs bosons, SUSY particles and so on) and $`X`$ stands for any unobserved particle produced by the collision. The typical scale $`Q`$ of the scattering process is set by the invariant mass or the transverse momentum of the hard probe, and the notation $`\{\mathrm{}\}`$ stands for any other relevant scale and kinematic variable of the process. For instance, in the case of $`W`$ production we have $`Q=M_W`$ and $`\{\mathrm{}\}=\{Q_{},y,\mathrm{}\}`$, where $`M_W,Q_{}`$ and $`y`$ are the mass of the vector boson, its transverse momentum and its rapidity, respectively.
The factorization formula (2) involves the convolution of the partonic cross sections $`\widehat{\sigma }_{ab}`$ (where $`a,b=q,\overline{q},g)`$ and the parton distributions $`f_{a/h}(x,\mu _F^2)`$ of the colliding hadrons. If the hard probe $`H`$ is a hadron or a photon, the factorization formula has to include an additional convolution with the corresponding parton fragmentation function $`d_{a/H}(z,\mu _F^2)`$.
The term $`𝒪\left((\mathrm{\Lambda }_{QCD}/Q)^p\right)`$ on the right-hand side of Eq. (2) generically denotes non-perturbative contributions (hadronization effects, multiparton interactions, contributions of the soft underlying event, and so on). Provided the hard-scattering process (1) is sufficiently inclusive<sup>§</sup><sup>§</sup>§ More precisely, it has to be defined in an infrared- and collinear-safe manner., $`\widehat{\sigma }_{ab}`$ is computable as a power series expansion in $`\alpha _\mathrm{S}(Q^2)`$ and the non-perturbative contributions are (small) power-suppressed corrections (i.e. the power $`p`$ is positive) as long as the hard-scattering scale $`Q`$ is larger than few hundred MeV, the typical size of the QCD scale $`\mathrm{\Lambda }_{QCD}`$.
The parton densities $`f_{a/h}(x,\mu _F^2)`$ are phenomenological distributions that describe how partons are bounded in the colliding hadrons. Although they are not calculable in QCD perturbation theory, the parton densities are universal (process-independent) quantities. The scale $`\mu _F`$ is a factorization scale introduced in Eq. (2) to separate the bound-state effects from the perturbative interactions of the partons. The physical cross section $`\sigma (p_1,p_2;Q,\{\mathrm{}\})`$ does not depend on this arbitrary scale, but parton densities and partonic cross sections separately depend on $`\mu _F`$. In particular, higher-order contributions to $`\widehat{\sigma }_{ab}(x_1p_1,x_2p_2;Q,\{\mathrm{}\};\mu _F^2)`$ contain corrections of relative order $`(\alpha _\mathrm{S}(Q^2)\mathrm{ln}Q^2/\mu _F^2)^n`$. If $`\mu _F`$ is very different from $`Q`$, these corrections become large and spoil the reliability of the perturbative expansion. Thus, in practical applications of the factorization formula (2), the scale $`\mu _F`$ is set approximately equal to the hard scale $`Q`$ and variations of $`\mu _F`$ around this central value are used to estimate the uncertainty of the perturbative expansion.
The lower limit $`x_{\mathrm{min}}`$ of the integrations over the parton momentum fractions $`x_1`$ and $`x_2`$, as well as the values of $`x_1`$ and $`x_2`$ that dominate the convolution integral in Eq. (2), are controlled by the kinematics of the hard-scattering process. Typically we have $`x_{\mathrm{min}}>Q^2/S`$, where $`S=(p_1+p_2)^2`$ is the square of the centre-of-mass energy of the collision. If the hard probe is a state of invariant mass $`M`$ and rapidity $`y`$, the dominant values of the momentum fractions are $`x_{1,2}(Me^{\pm y})/\sqrt{S}`$ (see Fig. 1). Thus varying $`M`$ and $`y`$ at fixed $`\sqrt{S}`$, we are sensitive to partons with different momentum fractions. Increasing $`\sqrt{S}`$ the parton densities are probed in a kinematic range that extends towards larger values of $`Q`$ and smaller values of $`x_{1,2}`$.
## 2 Parton densities
The parton densities are an essential ingredient to study hard-scattering collisions. Once the partonic cross sections have been perturbatively computed, cross section measurements can be used to determine the parton densities. Then, they can in turn be used to predict cross sections for other hard-scattering processes.
The dependence of the parton densitiesIn the following the parton densities of the proton $`f_{a/p}`$ are simply denoted by $`f_a`$ and those of the antiproton are obtained by using charge-conjugation invariance, i.e. $`f_{a/\overline{p}}=f_{\overline{a}/p}=f_{\overline{a}}`$. $`f_a(x,\mu ^2)`$ on the momentum fraction $`x`$ and their absolute value at any fixed scale $`\mu `$ are not computable in perturbation theory. However, the scale dependence is perturbatively controlled by the DGLAP evolution equation
$$\frac{df_a(x,\mu ^2)}{d\mathrm{ln}\mu ^2}=\underset{b}{}_x^1\frac{dz}{z}P_{ab}(\alpha _\mathrm{S}(\mu ^2),z)f_a(x/z,\mu ^2).$$
(3)
The kernels $`P_{ab}(\alpha _\mathrm{S},z)`$ are the Altarelli–Parisi (AP) splitting functions. As the partonic cross sections in Eq. (2), the AP splitting functions can be computed as a power series expansion in $`\alpha _\mathrm{S}`$:
$$P_{ab}(\alpha _\mathrm{S},z)=\alpha _\mathrm{S}P_{ab}^{(LO)}(z)+\alpha _\mathrm{S}^2P_{ab}^{(NLO)}(z)+\alpha _\mathrm{S}^3P_{ab}^{(NNLO)}(z)+𝒪(\alpha _\mathrm{S}^4).$$
(4)
The leading order (LO) and next-to-leading order (NLO) terms $`P_{ab}^{(LO)}(z)`$ and $`P_{ab}^{(NLO)}(z)`$ in the expansion are known . These first two terms are used in most of the QCD studies. Having determined $`f_a(x,Q_0^2)`$ at a given input scale $`\mu =Q_0`$, the evolution equation (3) can be used to compute the parton densities at different perturbative scales $`\mu `$ and larger values of $`x`$.
The parton densities are determined by performing global fits to data from deep-inelastic scattering (DIS), Drell–Yan (DY), prompt-photon and jet production. The method consists in parametrizing the parton densities at some input scale $`Q_0`$ and then adjusting the parameters to fit the data. The parameters are usually constrained by imposing the positivity of the parton densities $`(f_a(x,\mu ^2)0)`$ and the momentum sum rule $`(_a_0^1𝑑xxf_a(x,\mu ^2)=1)`$.
The present knowledge on the parton densities of the proton is reviewed in Refs. . Their typical behaviour is shown in Fig. 2. All densities decrease at large $`x`$. At small $`x`$ the valence quark densities vanish and the gluon density dominates. The sea-quark densities also increase at small $`x`$ because they are driven by the strong rise of the gluon density and the splitting of gluons in $`q\overline{q}`$ pairs. Note that the quark densities are not flavour-symmetric either in the valence sector $`(u_vd_v)`$ or in the sea sector $`(\overline{u}\overline{d})`$.
In addition to having the best estimate of the parton densities, it is important to quantify the corresponding uncertainty. This is a difficult issue. The uncertainty depends on the kinematic range in $`x`$ and $`Q^2`$. Moreover, it cannot be reliably estimated by simply comparing the parton densities obtained by different global fits. In fact, a lot of common systematic and common assumptions affect the global-fit procedures. Recent attempts to obtain parton densities with error bands that take into account correlations between experimental errors are described in Refs. . Some important theoretical uncertainties that are still to be understood are also discussed in Ref. .
The overall conclusion is that the quark densitiesUncertainties on the determination of the quark densities at very high $`x`$ are discussed in Refs. . are reasonably well constrained and determined by DIS and DY processes, while the gluon density is certainly more uncertain . At small $`x`$ $`(x<10^3)`$, the gluon density $`f_g`$ is at present constrained by a single process, namely DIS at HERA. Thus, large higher-order corrections of the type $`(\alpha _\mathrm{S}\mathrm{ln}1/x)^n`$ could possibly affect the extraction of $`f_g`$. Assuming that $`f_g`$ is well determined at small $`x`$, the momentum sum rule reasonably constrains $`f_g`$ at intermediate values of $`x`$ $`(x10^2)`$. Jet production at the Tevatron at low to moderate values of the jet transverse energy $`E_T`$ can also be useful in constraining the gluon distribution in the range $`0.05<x<0.2`$. At large $`x`$ $`(x10^1)`$, the most sensitive process to $`f_g`$ is prompt-photon production. Since, at present, prompt-photon data are not well described/predicted by perturbative QCD calculations, they cannot be used for a precise determination of $`f_g`$. Further discussion on these points is given in Sect. 3.
The conclusion that the gluon density is not well known can also be drawn by inspection (see Fig. 3) of the differences between the most updated analyses performed by the CTEQ Collaboration and the MRST group.
The differences between the MRST gluons and the CTEQ ones are due to the fact that the two groups used different data sets. The various gluon densities are very similar at small $`x`$, because in this region both groups used the HERA data. The MRST group includes prompt-photon data in the global fit: these data constrain the gluon directly at $`x>10^1`$ and indirectly (by the momentum sum rule) at $`x10^2`$. The CTEQ group does not use prompt-photon data, but it includes Tevatron data on the one-jet inclusive cross section. These data give a good constraint on $`f_g`$ in the region $`10^2<x<10^1`$.
There are also differences within the MRST and CTEQ sets. The various gluon densities of the MRST set correspond to different values of the non-perturbative transverse-momentum smearing that can be introduced to describe the differences among the prompt-photon data that are available at several centre-of-mass energies. The CTEQ5M and CTEQ5HJ gluons correspond to different assumptions on the parametrization of the functional form of $`f_g(x,Q_0^2)`$ at large $`x`$; the CTEQ5M set corresponds to the minimum-$`\chi ^2`$ solution of the fit while the CTEQ5HJ set (with a slightly higher $`\chi ^2`$) provides the best fit to the high-$`E_T`$ tail of the CDF and D0 jet cross sections.
This brief illustration shows that the differences in the most recent parton densities are mainly due to either inconsistencies between data sets and/or poor theoretical understanding of them. A more quantitative picture of the dependence on $`x`$ and $`Q^2`$ of the gluon density uncertainty is presented in Fig. 4.
We can see that the DIS and DY data sets weakly constrain $`f_g`$ for $`x>10^1`$. Since the AP splitting functions lead to negative scaling violation at large $`x`$, when $`f_g(x,Q^2)`$ is evolved at larger scales $`Q`$ according to Eq. (3) the gluon uncertainty is diluted: it propagates at smaller values of $`x`$ and its size is reduced at fixed $`x`$.
Figure 5 shows the typical predictions for hard-scattering cross sections at the Tevatron and the LHC, as obtained by using the parton densities of the MRST set. These predictions have to be supplemented with the corresponding uncertainties coming from the determination of the parton densities and from perturbative corrections beyond the NLO. Owing to the increased centre-of-mass energy and to QCD scaling violation (see Fig. 4), the kinematic region with small uncertainties is larger at the LHC than at the Tevatron.
For most of the QCD processes at the LHC, the uncertainty from the parton densities is smaller than $`\pm 10\%`$ and, in particular, it is smaller than the uncertainty from higher-order corrections. Some relevant exceptions are the single-jet, $`W/Z`$ and top quark cross sections. In the case of the single-jet inclusive cross section at high $`E_T`$ $`(E_T>2`$ TeV), the uncertainty from the poorly known gluon density at high $`x`$ is larger than that $`(\pm 10\%)`$ from higher-order corrections. The $`W`$ and $`Z`$ production cross sections are dominated by $`q\overline{q}`$ annihilation. Since the quark densities are well known, the ensuing uncertainty on the $`W/Z`$ cross section is small $`(\pm 5\%)`$. Nonetheless, in this case the uncertainty from higher-order corrections is even smaller, since the partonic cross sections for the DY process are known at the next-to-next-to-leading order (NNLO) in perturbation theory. In the case of top-quark production at the LHC, the gluon channel dominates and leads to an uncertainty of $`\pm 10\%`$ on the total cross section. Also for this process, however, the perturbative component is known beyond the NLO. Including all-order resummation of soft-gluon contributions , the estimated uncertainty from unknown higher-order corrections is approximately $`\pm 5\%`$ .
## 3 The gluon density issue
At present, the processes<sup>\**</sup><sup>\**</sup>\**The rôle of jet production at the Tevatron has briefly been recalled in Sect. 2, and it is discussed in detail in Ref. . that are, in principle, most sensitive to the gluon density are DIS at HERA, $`b`$-quark production at the Tevatron, and prompt-photon production at fixed-target experiments. These processes constrain $`f_g`$ for $`x<10^3`$, $`x10^3`$$`10^2`$ and $`x>10^1`$, respectively. Nonetheless, the gluon density is, in practice, not well determined. The issue (or, perhaps, the puzzle) is that from a phenomenological viewpoint the standard theory, namely perturbative QCD at NLO, works pretty well for $`x<10^3`$ but not so well at larger values of $`x`$, while from theoretical arguments we should expect just the opposite to happen. This issue is discussed below mainly in its perturbative aspects. We should however keep it in mind that all these processes are dominated by hard-scattering scales $`Q`$ of the order of few GeV. Different types of non-perturbative contributions can thus be important.
From the study of DIS at HERA we can extract information on the gluon and sea-quark densities of the proton. The main steps in the QCD analysis of the structure functions at small values of the Bjorken variable $`x`$ are the following. The measurement of the proton structure function $`F_2(x,Q^2)q_S(x,Q^2)`$ directly determines the sea-quark density $`q_S=x(f_q+f_{\overline{q}})`$. Then, the DGLAP evolution equation (3) or, more precisely, the following equations (the symbol $``$ denotes the convolution integral with respect to $`x`$):
$`dF_2(x,Q^2)/d\mathrm{ln}Q^2`$ $``$ $`P_{qq}q_S+P_{qg}g,`$ (5)
$`dg(x,Q^2)/d\mathrm{ln}Q^2`$ $``$ $`P_{gq}q_S+P_{gg}g,`$ (6)
are used to extract a gluon density $`g(x,Q^2)=xf_g(x,Q^2)`$ that agrees with the measured scaling violation in $`dF_2(x,Q^2)/d\mathrm{ln}Q^2`$ (according to Eq. (5)) and fulfils the self-consistency equation (6).
The perturbative-QCD ingredients in this analysis are the AP splitting functions $`P_{ab}(\alpha _\mathrm{S},x)`$. Once they are known (and only then), the non-perturbative gluon density can be determined.
The standard perturbative-QCD framework to extract $`g(x,Q^2)`$ consists in using the truncation of the AP splitting functions at the NLO. This approach has been extensively compared with structure function data over the last few years and it gives a good description of the HERA data, down to low values of $`Q^22\mathrm{GeV}^2`$. The NLO QCD fits simply require a slightly steep input gluon density at these low momentum scales. Typically , we have $`g(x,Q_0^2)x^\lambda `$, with $`\lambda 0.2`$ at $`Q_0^22\mathrm{GeV}^2`$, and the data constrain $`g(x,Q_0^2)`$ with an uncertainty of approximately $`\pm 20\%`$.
Although it is phenomenologically successful, the NLO approach is not fully satisfactory from a theoretical viewpoint. The truncation of the splitting functions at a fixed perturbative order is equivalent to assuming that the dominant dynamical mechanism leading to scaling violations is the evolution of parton cascades with strongly-ordered transverse momenta. However, at high energy this evolution takes place over large rapidity intervals $`(\mathrm{\Delta }y\mathrm{ln}1/x)`$ and diffusion in transverse momentum becomes relevant. Formally, this implies that higher-order corrections to $`P_{ab}(\alpha _\mathrm{S},x)`$ are logarithmically enhanced:
$$P_{ab}(\alpha _\mathrm{S},x)\frac{\alpha _\mathrm{S}}{x}+\frac{\alpha _\mathrm{S}}{x}(\alpha _\mathrm{S}\mathrm{ln}x)+\mathrm{}+\frac{\alpha _\mathrm{S}}{x}(\alpha _\mathrm{S}\mathrm{ln}x)^n+\mathrm{}.$$
(7)
At asymptotically small values of $`x`$, resummation of these corrections is mandatory to obtain reliable predictions.
Small-$`x`$ resummation is, in general, accomplished by the BFKL equation . In the context of structure-function calculations, the BFKL equation provides us with improved expressions of the AP splitting functions $`P_{ab}(\alpha _\mathrm{S},x)`$, in which the leading logarithmic (LL) terms $`(\alpha _\mathrm{S}\mathrm{ln}x)^n`$, the next-to-leading logarithmic (NLL) terms $`\alpha _\mathrm{S}(\alpha _\mathrm{S}\mathrm{ln}x)^n`$, and so forth, are systematically summed to all orders $`n`$ in $`\alpha _\mathrm{S}`$. The present theoretical status of small-$`x`$ resummation is discussed in Ref. . Since in the small-$`x`$ region the gluon channel dominates, only the gluon splitting functions $`P_{gg}`$ and $`P_{gq}`$ contain LL contributions. These are known to be positive but numerically smaller than naively expected (the approach to the asymptotic regime is much delayed by cancellations of logarithmic corrections that occur at the first perturbative orders in $`P_{gg}`$ and $`P_{gq}`$). The NLL terms in the quark splitting functions $`P_{qg}`$ and $`P_{qq}`$ are known and turn out to be positive and large. A very important progress is the recent calculation of the NLL terms in $`P_{gg}`$, which are found to be negative and large. The complete NLL terms in $`P_{gq}`$ are still unknown.
The results of Refs. , the large size of the NLL terms and the alternating sign (from the LL to the NLL order and from the gluon to the quark channel) of the resummed small-$`x`$ contributions have prompted a lot of activity (see the list of references in Ref. ) on the conceptual basis and the phenomenological implications of small-$`x`$ resummation. This activity is still in progress and definite quantitative conclusions on the impact of small-$`x`$ resummation at HERA cannot be drawn yet.
At the same time, the capability of the fixed-order approach to produce a good description of the proton structure function $`F_2(x,Q^2)`$ at HERA cannot be used to conclude that the small-$`x`$ behaviour of the gluon density is certainly well determined. In fact, by comparing LO and NLO results, we could argue that the ensuing theoretical uncertainty on $`f_g`$ is sizeable . Going from LO to NLO, we can obtain stable predictions for $`F_2`$, but we have to vary the gluon density a lot. As shown in Fig. 6, the NLO gluon density sizeably differs from its LO parametrization, not only in absolute normalization but also in $`x`$-shape. For instance, at $`x=10^4`$ and $`Q^2=5\mathrm{GeV}^2`$ the NLO gluon is a factor of 2 smaller than the LO gluon. This can be understood from the fact that the scaling violation of $`F_2`$ is produced by the convolution $`P_{qg}g`$ (see the right-hand side of Eq. (5)). The quark splitting function $`P_{qg}`$ behaves as
$$P_{qg}(\alpha _\mathrm{S},x)\alpha _\mathrm{S}P_{qg}^{(LO)}(x)\left[1+2.2\frac{C_A\alpha _\mathrm{S}}{\pi }\frac{1}{x}+\mathrm{}\right],$$
(8)
where the LO term $`P_{qg}^{(LO)}(x)`$ is flat at small $`x`$, whereas the NLO correction is steep. To obtain a stable evolution of $`F_2`$, the NLO steepness of $`P_{qg}`$ has to be compensated by a gluon density that is less steep at NLO than at LO. This has to be kept in mind when concluding on the importance of small-$`x`$ resummation because the NLO steepness of $`P_{qg}`$ is the lowest-order manifestation of BFKL dynamics in the quark channel.
In the large-$`x`$ region, there is a well-known correlation between $`\alpha _\mathrm{S}`$ and $`f_g`$. At small $`x`$, there is an analogous strong correlation between the $`x`$-shapes of $`P_{qg}`$ and $`f_g`$. In the fixed-order QCD analysis of $`F_2`$, large NLO perturbative corrections at small $`x`$ can be balanced by the extreme flexibility of parton density parametrizations. It is difficult to disentangle this correlation between process-dependent perturbative contributions and non-perturbative parton densities from the study of a single quantity, as in the case of $`F_2`$ at HERA. The uncertainty on the gluon density at small $`x`$, as estimated from the NLO QCD fits of the HERA data, is evidently only a lower limit on the actual uncertainty on $`f_g`$.
The production of $`b`$ quarks at the Tevatron is also sensitive to the gluon density at relatively small values of $`x`$. The comparison between Tevatron data and perturbative-QCD predictions at NLO is shown in Fig. 7. Using standard sets of parton densities, the theoretical predictions typically underestimate the measured cross section by a factor of 2. This certainly is disappointing, although justifiable by the large theoretical uncertainty of the perturbative calculation . A lower limit on this uncertainty can be estimated by studying the scale dependence and the convergence of the perturbative expansion. Varying the factorization and renormalization scales by a factor of four around the $`b`$-quark mass $`m_b`$, the NLO cross section varies by a factor of almost 2 at the Tevatron and by a factor of 4–5 at the LHC . Similar factors are obtained by considering the ratio of the NLO and LO cross sections.
The present theoretical predictions for $`b`$-quark production at hadron colliders certainly need to be improved . Since the hard scale $`Qm_b`$ is not very large, a possible improvement regards estimates of non-perturbative contributions (for instance, effects of the fragmentation of the $`b`$-quark and of the intrinsic transverse momentum of the colliding partons). As for the evaluation of perturbative contributions at higher orders, the resummation of logarithmic terms of the type $`\alpha _\mathrm{S}^n\mathrm{ln}^n(p_t/m_b)`$ is important when the transverse momentum $`p_t`$ of the $`b`$ quark is much larger than $`m_b`$. The resummation of small-$`x`$ logarithmic contributions $`\alpha _\mathrm{S}^n\mathrm{ln}^nx`$ can also be relevant, because $`x2m_b/\sqrt{S}`$ is as small as $`10^3`$ at the Tevatron and as $`10^4`$ at the LHC. The theoretical tool to perform this resummation, namely the $`k_{}`$-factorization approach , is available. Updated phenomenological studies based on this tool and on the information from small-$`x`$ DIS at HERA would be interesting.
Prompt-photon production at fixed-target experiments is sensitive to the behaviour of the gluon density at large $`x`$ $`(x>0.1)`$. The theoretical predictions for this process, however, are not very accurate. Figure 8 shows the factorization- and renormalization-scale dependence of the perturbative cross section for the case of the E706 kinematics. If the scale is varied by a factor of 4 around the transverse energy $`E_T`$ of the prompt photon, the LO cross section varies by a factor of almost 4. Going to NLO the situation improves, but not very much, because the NLO cross section still varies by a factor of about 2.
A detailed comparison between NLO QCD calculations and data from the ISR and fixed-target experiments has recently been performed in Ref. . As shown in Fig. 9, the overall agreement with the theory is not satisfactory, even taking into account the uncertainty coming from scale variations in the theoretical predictions. Modifications of the gluon density can improve the agreement with some data sets only at the expense of having larger disagreement with other data sets. The differences between experiments at similar centre-of-mass energies (see, for instance, E706 pBe/530 at $`\sqrt{S}=31.6\mathrm{GeV}`$ and WA70 pp at $`\sqrt{S}=23\mathrm{GeV})`$ are much larger than expected from perturbative scaling violations. This can possibly suggest inconsistencies of experimental origin.
Another (not necessarily alternative) origin of the differences between data and theory could be the presence of non-perturbative effects that are not included in the NLO perturbative calculation. This explanation has been put forward in Refs. by introducing some amount of intrinsic<sup>††</sup><sup>††</sup>††To be precise, in Ref. the $`k_{}`$ of the colliding partons is not called ‘intrinsic’, but it is more generically called the $`k_{}`$ ‘from initial-state soft-gluon radiation’. transverse momentum $`k_{}`$ of the colliding partons. Owing to the steeply falling $`E_T`$ distribution $`(d\sigma /dE_T1/E_T^7)`$ of the prompt photon, even a small transverse-momentum kick<sup>‡‡</sup><sup>‡‡</sup>‡‡The $`E_T`$ distribution of the single-photon is not calculable down to $`E_T=0`$ or, in other words, $`d\sigma /dE_T`$ is not integrable in the entire kinematic range of $`E_T`$. Thus, the intrinsic $`k_{}`$ of the incoming partons does not simply produce a shift of events from the low-$`E_T`$ to the high-$`E_T`$ region. For this reason, the terminology ‘$`k_{}`$ kick’ seems to be more appropriate than ‘$`k_{}`$ smearing’. can indeed produce a large effect on the cross section, in particular, at small values of $`E_T`$. Phenomenological investigations show that this additional $`k_{}`$ kick can lead to a better agreement between calculations and data. The E706 data suggest the value $`k_{}1.2\mathrm{GeV}`$, the WA70 data prefer no $`k_{}`$, and the UA6 data in the intermediate range of centre-of-mass energy $`(\sqrt{S}=24.3\mathrm{GeV}`$) may prefer an intermediate value of $`k_{}`$. Similar conclusions are obtained in the analysis by the MRST group .
A precise physical understanding of $`k_{}`$ effects is still missing. On one side, since the amount of $`k_{}`$ suggested by prompt-photon data varies with $`\sqrt{S}`$, it is difficult to argue that the transverse momentum is really ‘intrinsic’ and has an entirely non-perturbative origin. On the other side, in the case of the inclusive production of a single photon, a similar effect cannot be justified by higher-order logarithmic corrections produced by perturbative soft-gluon radiation (see Sect. 5). A lot of model-dependent assumptions (and ensuing uncertainties) certainly enter in the present implementations of the $`k_{}`$ kick. A general framework to consistently include non-perturbative transverse-momentum effects in perturbative calculations is not yet available. Recent proposals with this aim are presented in Refs. and .
Further studies on the consistency between different prompt-photon experiments and on the issue of intrinsic-$`k_{}`$ effects in hadron–hadron collisions are necessary. Owing to the present theoretical (and, possibly, experimental) uncertainties, it is difficult to use prompt-photon data to accurately determine the gluon density at large $`x`$. Other recent theoretical improvements, such as soft-gluon resummation, of the perturbative calculations for prompt-photon production at large $`x_T=2E_T/\sqrt{S}`$ are discussed in Sect. 5.
Studies of other single-particle inclusive cross sections, such as $`\pi ^0`$ cross sections , can be valuable to constrain the parton densities and could possibly help to clarify some of the experimental and theoretical issues arisen by prompt-photon production.
## 4 Partonic cross sections: fixed-order expansions, <br>resummed calculations, parton showers
The calculation of hard-scattering cross sections according to the factorization formula (2) requires the knowledge of the partonic cross sections $`\widehat{\sigma }`$, besides that of the parton densities. The partonic cross sections are usually computed by truncating their perturbative expansion at a fixed order in $`\alpha _\mathrm{S}`$:
$`\widehat{\sigma }(p_1,p_2;Q,\{Q_1,\mathrm{}\};\mu _F^2)`$ $`=`$ $`\alpha _\mathrm{S}^k(\mu _R^2)\{\widehat{\sigma }^{(LO)}(p_1,p_2;Q,\{Q_1,\mathrm{}\})`$
$`+\alpha _\mathrm{S}(\mu _R^2)\widehat{\sigma }^{(NLO)}(p_1,p_2;Q,\{Q_1,\mathrm{}\};\mu _R^2;\mu _F^2)`$
$`+\alpha _\mathrm{S}^2(\mu _R^2)\widehat{\sigma }^{(NNLO)}(p_1,p_2;Q,\{Q_1,\mathrm{}\};\mu _R^2;\mu _F^2)+\mathrm{}\}.`$
The scale $`\mu _R`$ is the arbitrary renormalization scale introduced to define the perturbative expansion. Although the ‘exact’ partonic cross section on the left-hand side of Eq. (4) does not depend on $`\mu _R`$, each term on the right-hand side (and, hence, any fixed-order truncation) separately depends on it.
The LO (or tree-level) term $`\widehat{\sigma }^{(LO)}`$ gives only an estimate of the order of magnitude of the partonic cross section, because at this order $`\alpha _\mathrm{S}`$ is not unambiguously defined. Equivalently, we can say that since $`\widehat{\sigma }^{(LO)}`$ does not depend on $`\mu _R`$, the size of its contribution can be varied quite arbitrarily by changing $`\mu _R`$ in its coefficient $`\alpha _\mathrm{S}^k(\mu _R^2)`$. The strong coupling $`\alpha _\mathrm{S}`$ can be precisely defined only starting from NLO. A ‘reliable’ estimate of the central value of $`\widehat{\sigma }`$ thus requires the knowledge of (at least) the NLO term $`\widehat{\sigma }^{(NLO)}`$. This term explicitly depends on $`\mu _R`$ and this dependence begins to compensate that of $`\alpha _\mathrm{S}(\mu _R^2)`$.
In general, the $`n`$-th term in the curly bracket of Eq. (4) contains contributions of the type $`(\alpha _\mathrm{S}(\mu _R^2)\mathrm{ln}Q/\mu _R)^n`$. If $`\mu _R`$ is very different from the hard scale $`Q`$, these contributions become large and spoil the reliability of the truncated expansion (4). Thus, in practical applications the scale $`\mu _R`$ should be set approximately equal to the hard scale $`Q`$. As mentioned in Sect. 3, variations of $`\mu _R`$ around this central value are typically used to set a lower limit on the theoretical uncertainty of the perturbative calculation.
A better estimate of the accuracy of any perturbative expansion is obtained by considering the effect of removing the last perturbative term that has been computed. Since $`\alpha _\mathrm{S}`$ can be precisely defined only at NLO, this procedure can consistently be applied to Eq. (4) only as from its NNLO term. A ‘reliable’ estimate of the theoretical error on $`\widehat{\sigma }`$ thus requires the knowledge of the NNLO term $`\widehat{\sigma }^{(NNLO)}`$ in Eq. (4).
The LO and NLO approximations of $`\widehat{\sigma }`$ are used at present in (most of) the fixed-order QCD calculations. Prospects towards NNLO calculations of partonic cross sections and AP splitting functions are reviewed in Refs. .
The fixed-order expansion (4) provides us with a well-defined and systematic framework to compute the partonic cross section $`\widehat{\sigma }(p_1,p_2;Q,\{Q_1,\mathrm{}\};\mu _F^2)`$ of any hard-scattering process that is sufficiently inclusive or, more precisely, that is defined in an infrared- and collinear-safe manner. However, the fixed-order expansion is reliable only when all the kinematical scales $`Q,\{Q_1,\mathrm{}\}`$ are of the same order of magnitude. When the hard-scattering process involves two (or several) very different scales, say $`Q_1Q`$, the $`\mathrm{N}^n\mathrm{LO}`$ term in Eq. (4) can contain double- and single-logarithmic contributions of the type $`(\alpha _\mathrm{S}L^2)^n`$ and $`(\alpha _\mathrm{S}L)^n`$ with $`L=\mathrm{ln}(Q_1/Q)1`$. These terms spoil the reliability of the fixed-order expansion and have to be summed to all orders by systematically defining logarithmic expansions (resummed calculations).
Typical large logarithms, $`L=\mathrm{ln}Q/Q_0`$, are those related to the evolution of the parton densities from a low input scale $`Q_0`$ to the hard-scattering scale $`Q`$. These logarithms are produced by collinear radiation from the colliding partons and give single-logarithmic contributions. They never explicitly appear in the calculation of the partonic cross section, because they are systematically (LO, NLO and so forth) resummed in the evolved parton densities $`f(x,Q^2)`$ by using the DGLAP equation (3).
Different large logarithms, $`L=\mathrm{ln}Q/\sqrt{S}`$, appear when the centre-of-mass energy $`\sqrt{S}`$ of the collision is much larger than the hard scale $`Q`$. These small-$`x`$ $`(x=Q/\sqrt{S})`$ logarithms are produced by multiple radiation over the wide rapidity range that is available at large energy. They usually give single-logarithmic contributions that can be resummed by using the BFKL equation. BFKL resummation is relevant to DIS structure functions at small values of the Bjorken variable $`x`$ (see Sect. 3) and it can also be important at the LHC for the production of $`b`$ quarks and of prompt photons at relatively low $`E_T`$.
Another class of large logarithms is associated to the bremsstrahlung spectrum of soft gluons. Since soft gluons can be radiated collinearly, they give rise to double-logarithmic contributions to the partonic cross section:
$$\widehat{\sigma }\alpha _\mathrm{S}^k\widehat{\sigma }^{(LO)}\left\{1+\underset{n=1}{\overset{\mathrm{}}{}}\alpha _\mathrm{S}^n\left(C_{2n}^{(n)}L^{2n}+C_{2n1}^{(n)}L^{2n1}+C_{2n2}^{(n)}L^{2n2}+\mathrm{}\right)\right\}.$$
(10)
Soft-gluon resummation is discussed in Sect. 5.
A related approach to evaluate higher-order contributions to the partonic cross sections is based on Monte Carlo parton showers (see and the updated list of references in ). Rather than computing exactly $`\widehat{\sigma }^{(NLO)}`$, $`\widehat{\sigma }^{(NNLO)}`$ and so forth, the parton shower gives an all-order approximation of the partonic cross section in the soft and collinear regions. In this respect, the computation of the partonic cross sections performed by parton showers is somehow similar to that obtained by soft-gluon resummed calculations. There is, however, an important conceptual difference between the two approaches. This difference and the limits of applicability of the parton-shower method are briefly recalled below. Apart from these limits, parton-shower calculations can give some advantages. Multiparton kinematics can be treated exactly. The parton shower can be supplemented with models of non-perturbative effects (hadronization, intrinsic $`k_{}`$, soft underlying event) to provide a complete description of the hard-scattering process at the hadron level.
For a given cross section, resummed calculations can in principle be performed to any logarithmic accuracy. The logarithmic accuracy achievable by parton showers is instead intrinsically limited by quantum mechanics. The parton-shower algorithms are probabilistic. Starting from the LO cross section, the parton shower generates multiparton final states according to a probability distribution that approximates the square of the QCD matrix elements. The approximation is based on the universal (process-independent) factorization properties of multiparton matrix elements in the soft and collinear limits. Although the matrix element does factorize, its square contains quantum interferences, which are not positive-definite and, in general, cannot be used to define probability distributions. To leading infrared accuracy, this problem is overcome by exploiting QCD coherence (see Refs. and referencees therein): soft gluons radiated at large angle from the partons involved in the LO subprocess destructively interfere. This quantum mechanical effect can be simply implemented by enforcing an angular-ordering constraint on the phase space available for the parton shower evolution. Thus, angular-ordered parton showers can consistently compute the first two dominant towers ($`\alpha _\mathrm{S}^nL^{2n}`$ and $`\alpha _\mathrm{S}^nL^{2n1}`$) of logarithmic contributions in Eq. (10). However, parton showers contain also some subleading logarithmic contributions. For instance, they correctly compute the single-logarithmic terms $`\alpha _\mathrm{S}^nL^n`$ of purely collinear origin that lead to the LO evolution of the parton densities. Moreover, as discussed in Ref. by a comparison with resummed calculations, in the case of hard-scattering processes whose LO subprocess involves two coloured partons (e.g. DIS or DY production), angular-ordered parton showers have a higher logarithmic accuracy: they can consistently evaluate the LL and NLL terms in Eq. (15). The extension of parton-shower algorithms to higher logarithmic accuracy is not necessarily feasible and is, in any case, challenging.
Of course, because of quantum interferences and quantum fluctuations, the probabilistic parton-shower approach cannot be used to systematically perform exact calculations at NLO, NNLO and so forth. Nonetheless, important progress has been made to include matrix element corrections in parton shower algorithms . The purpose is to consider the multiparton configurations generated by parton showering from the LO matrix element and to correct them in the hard (non-soft and non-collinear) region by using the exact expressions of the higher-order matrix elements. Hard matrix element corrections to parton showers have been implemented for top quark decay and for production of $`W,Z`$ and DY lepton pairs . The same techniques could be applied to other processes, as, for instance, production of Higgs boson and vector-boson pairs .
Note also that, at present, angular-ordered parton showers cannot be considered as true ‘next-to-leading’ tools, even where their logarithmic accuracy is concerned. The consistent computation of the first two towers of logarithmic contributions in Eq. (10) is not sufficient for this purpose. For instance, to precisely introduce an NLO definition of $`\alpha _\mathrm{S}`$, we should control all the terms obtained by the replacement $`\alpha _\mathrm{S}\alpha _\mathrm{S}+c\alpha _\mathrm{S}^2+𝒪(\alpha _\mathrm{S}^3)`$. When it is introduced in the towers of double-logarithmic terms $`\alpha _\mathrm{S}^nL^{2n}`$ of Eq. (10), this replacement leads to contributions of the type $`\alpha _\mathrm{S}^{n+1}L^{2n}\alpha _\mathrm{S}^nL^{2n2}`$. Since these contributions are not fully computable at present, the parameter $`\alpha _\mathrm{S}`$ used in the parton showers corresponds to a simple LO parametrization of QCD running coupling.
## 5 Soft-gluon resummation
Double-logarithmic contributions due to soft gluons arise in all the kinematic configurations where radiation of real and virtual partons is highly unbalanced (see Ref. and referenes therein). For instance, this happens in the case of transverse-momentum distributions at low transverse momentum, in the case of hard-scattering production near threshold or when the structure of the final state is investigated with high resolution (internal jet structure, shape variables).
Soft-gluon resummation for jet shapes has been extensively studied and applied to hadronic final states produced by $`e^+e^{}`$ annihilation . Applications to hadron–hadron collisions have just begun to appear and have a large, yet uncovered, potential (from $`\alpha _\mathrm{S}`$ determinations to studies of non-perturbative dynamics).
Transverse-momentum logarithms, $`L=\mathrm{ln}Q^2/𝑸_{}^2`$, occur in the distribution of transverse momentum $`𝑸_{}`$ of systems with high mass $`Q`$ $`(QQ_{})`$ that are produced with a vanishing $`𝑸_{}`$ in the LO subprocess. Examples of such systems are DY lepton pairs, lepton pairs produced by $`W`$ and $`Z`$ decay, heavy quark–antiquark pairs, photon pairs and Higgs bosons. In these processes the LO transverse-momentum distribution is sharply peaked around $`𝑸_{}=0`$ ($`d\widehat{\sigma }/d^2𝑸_{}\delta ^{(2)}(𝑸_{}`$)). If the heavy system is produced with $`𝑸_{}^2Q`$, the emission of real radiation at higher orders is strongly suppressed and cannot balance the virtual contributions. The ensuing logarithms, $`L=\mathrm{ln}Q^2/𝑸_{}^2`$, diverge order by order when $`𝑸_{}0`$, but after all-order resummation they leads to a finite smearing of the LO distribution.
Threshold logarithms, $`L=\mathrm{ln}(1x)`$, occur when the tagged final state produced by the hard scattering is forced to carry a very large fraction $`x`$ ($`x1`$) of the available centre-of-mass energy $`\sqrt{S}`$. Also in this case, the radiative tail of real emission is stronly suppressed at higher perturbative orders. Oustanding examples of hard processes near threshold are DIS at large $`x`$ (here $`x`$ is the Bjorken variable), production of DY lepton pairs with large invariant mass $`Q`$ ($`x=Q/\sqrt{S}`$), production of heavy quark–antiquark pairs ($`x=2m_Q/\sqrt{S}`$), production of single jets and single photons at large transverse energy $`E_T`$ ($`x=2E_T/\sqrt{S}`$).
To emphasize the difference between transverse-momentum logarithms and threshold logarithms generated by soft gluons, it can be instructive to consider prompt-photon production. In the case of production of a photon pair<sup>*</sup><sup>*</sup>*The same discussion applies to the production of a DY lepton pair. with invariant mass squared $`Q^2=(p_1^{(\gamma )}+p_2^{(\gamma )})^2`$ and total transverse momentum $`𝑸_{}=𝒑_1^{(\gamma )}+𝒑_2^{(\gamma )}`$, transverse-momentum logarithms and threshold logarithms appear when $`𝑸_{}^2Q^2`$ and $`𝑸_{}^2(S/4Q^2)`$, respectively. However, in the case of production of a single photon with transverse energy (or, equivalently, transverse momentum) $`E_T`$, soft gluons can produce logarithms only in the threshold region $`x_T=2E_T/\sqrt{S}1`$. If the prompt photon has a transverse energy that is not closeEventually, when $`x_T1`$, higher-order corrections are single-logarithmically enhanced. This small-$`x`$ logarithms, $`(\alpha _\mathrm{S}\mathrm{ln}x_T)^n`$, have to be taken into account by BFKL resummation. to its threshold value, the emission of accompanying radiation is not kinematically suppressed and there are no soft logarithms analogous to those in the transverse-momentum distribution of a photon pair. In particular, there are no double-logarithmic contributions of the type $`(\alpha _\mathrm{S}\mathrm{ln}^2E_T^2/S)^n`$, and perturbative soft gluons are not distinguishable from perturbative hard gluons.
Studies of soft-gluon resummation for transverse-momentum distributions at low transverse momentum and for hard-scattering production near threshold started two decades ago . The physical bases for a systematic all-order summation of the soft-gluon contributions are dynamics and kinematics factorizations . The first factorization follows from gauge invariance and unitarity: in the soft limit, multigluon amplitudes fulfil factorization formulae given in terms of universal (process-independent) soft contributions. The second factorization regards kinematics and strongly depends on the actual cross section to be evaluated. If, in the appropriate soft limit, the multiparton phase space for this cross section can be written in a factorized way, resummation is analytically feasible in form of generalized exponentiation of the universal soft contributions that appear in the factorization formulae of the QCD amplitudes.
Note that the phase space depends in a non-trivial way on multigluon configurations and, in general, is not factorizable in single-particle contributionsIn the case of jet cross sections, for instance, phase-space factorization depends on the detailed definition of jets and it can easily be violated . Some jet algorithms, such as the $`k_{}`$-algorithm , have better factorization properties.. Moreover, even when phase-space factorization is achievable, it does not always occur in the space of the kinematic variables where the cross section is defined. Usually, it is necessary to introduce a conjugate space to overcome phase-space constraints. This is the case for transverse-momentum distributions and hard-scattering production near threshold. The relevant kinematical constraint for $`𝑸_{}`$-distributions is (two-dimensional) transverse-momentum conservation and it can be factorized by performing a Fourier transformation. Soft-gluon resummation for $`𝑸_{}`$-distributions is thus carried out in $`𝒃`$-space , where the impact parameter $`𝒃`$ is the variable conjugate to $`𝑸_{}`$ via the Fourier transformation. Analogously, the relevant kinematical constraint for hard-scattering production near threshold is (one-dimensional) energy conservation and it can be factorized by working in $`N`$-moment space , $`N`$ being the variable conjugate to the threshold variable $`x`$ (energy fraction) via a Mellin or Laplace transformation.
Using a short-hand notation, the general structure of the partonic cross section $`\widehat{\sigma }`$ after summation of soft-gluon contributions is
$$\widehat{\sigma }=\widehat{\sigma }_{\mathrm{res}.}+\widehat{\sigma }_{\mathrm{rem}.}.$$
(11)
The term $`\widehat{\sigma }_{\mathrm{res}.}`$ embodies the all-order resummation, while the remainder $`\widehat{\sigma }_{\mathrm{rem}.}`$ contains no large logarithmic contributions. The latter has the form
$$\widehat{\sigma }_{\mathrm{rem}.}=\widehat{\sigma }^{(\mathrm{f}.\mathrm{o}.)}\left[\widehat{\sigma }_{\mathrm{res}.}\right]^{(\mathrm{f}.\mathrm{o}.)},$$
(12)
and it is obtained from $`\widehat{\sigma }^{(\mathrm{f}.\mathrm{o}.)}`$, the truncation of the perturbative expansion for $`\widehat{\sigma }`$ at a given fixed order (LO, NLO, …), by subtracting the corresponding truncation $`\left[\widehat{\sigma }_{\mathrm{res}.}\right]^{(\mathrm{f}.\mathrm{o}.)}`$ of the resummed part. Thus, the expression on the right-hand side of Eq. (11) includes soft-gluon logarithms to all orders and it is matched to the exact (with no logarithmic approximation) fixed-order calculation. It represents an improved perturbative calculation that is everywhere as good as the fixed-order result, and much better in the kinematics regions where the soft-gluon logarithms become large ($`\alpha _\mathrm{S}L1`$). Eventually, when $`\alpha _\mathrm{S}L>1`$, the resummed perturbative contributions are of the same size as the non-perturbative contributions and the effect of the latter has to be implemented in the resummed calculation.
The resummed cross section has the following typical form:
$$\widehat{\sigma }_{\mathrm{res}.}=\alpha _\mathrm{S}^k_{\mathrm{inv}.}\widehat{\sigma }^{(LO)}CS,$$
(13)
where the integral $`_{\mathrm{inv}.}`$ denotes the inverse tranformation from the conjugate space where resummation is actually carried out. Methods to perform the inverse transformation are discussed in Refs. and for $`Q_{}`$-resummation and threshold resummation, respectively. The $`C`$ term has the perturbative expansion
$$C=1+C_1\alpha _\mathrm{S}+C_2\alpha _\mathrm{S}^2+\mathrm{}$$
(14)
and contains all the constant contributions in the limit $`L\mathrm{}`$ (the coefficients $`C_1,C_2,\mathrm{}`$ do not depend on the conjugate variable). The singular dependence on $`L`$ (more precisely, on the logarithm $`\stackrel{~}{L}`$ of the conjugate variable) is entirely exponentiated in the factor $`S`$:
$$S=\mathrm{exp}\left\{Lg_1(\alpha _\mathrm{S}L)+g_2(\alpha _\mathrm{S}L)+\alpha _\mathrm{S}g_3(\alpha _\mathrm{S}L)+\mathrm{}\right\}.$$
(15)
In the exponent, the function $`Lg_1`$ resums all the leading logarithmic (LL) contributions $`\alpha _\mathrm{S}^nL^{n+1}`$, while $`g_2`$ contains the next-to-leading logarithmic (NLL) terms $`\alpha _\mathrm{S}^nL^n`$ and so forth<sup>§</sup><sup>§</sup>§To compare this notation with that of Ref. , we can notice that our functions $`g_i`$ are obtained by the straightforward integration over $`\overline{\mu }`$ of the functions $`A(\alpha _\mathrm{S}(\overline{\mu }))`$ and $`B(\alpha _\mathrm{S}(\overline{\mu }))`$ of Ref. . In particular, our terms $`g_1,g_2,g_3`$ are not to be confused with the non-perturbative parameters of the same name used in Ref. . (all the functions $`g_i`$ are normalized as $`g_i(\lambda =0)=0`$). Note that the LL terms are formally suppressed by a power of $`\alpha _\mathrm{S}`$ with respect to the NLL terms, and so forth for the successive classes of logarithmic terms. Thus, this logarithmic expansion is as systematic as the fixed-order expansion in Eq. (4). In particular, using a matched NLL+NLO calculation, we can consistently $`i)`$ introduce a precise definition (say $`\overline{\mathrm{MS}}`$) of $`\alpha _\mathrm{S}(\mu )`$ and $`ii)`$ investigate the theoretical accuracy of the calculation by studying its renormalization-scale dependence.
The structure of the exponentiated resummed calculations discussed so far has to be contrasted with that obtained by organizing the logarithmic expansion on the right-hand side of Eq. (10) in terms of towers as
$$\widehat{\sigma }\alpha _\mathrm{S}^k\widehat{\sigma }^{(LO)}\left\{t_1(\alpha _\mathrm{S}L^2)+\alpha _\mathrm{S}Lt_2(\alpha _\mathrm{S}L^2)+\alpha _\mathrm{S}^2L^2t_3(\alpha _\mathrm{S}L^2)+\mathrm{}\right\},$$
(16)
where the double-logarithmic function $`t_1(\alpha _\mathrm{S}L^2)`$ and the successive functions are normalized as $`t_i(0)=\mathrm{const}.`$ While the ratio of two successive terms in the exponent of Eq. (15) is formally of the order of $`\alpha _\mathrm{S}`$, the ratio of two successive towers in Eq. (16) is formally of the order of $`\alpha _\mathrm{S}L`$. In other words, the tower expansion sums the double-logarithmic terms $`(\alpha _\mathrm{S}L^2)^n`$, then the terms $`\alpha _\mathrm{S}^nL^{2n1}\alpha _\mathrm{S}L(\alpha _\mathrm{S}L^2)^{n1}`$, and so forth; it thus assumes that the resummation procedure is carried out with respect to the large parameter $`\alpha _\mathrm{S}L^2`$ ($`\alpha _\mathrm{S}L^2<1`$). On the contrary, in Eq. (15) the large parameter is $`\alpha _\mathrm{S}L<1`$. The tower expansion allows us to formally extend the applicability of perturbative QCD to the region $`L<1/\sqrt{\alpha }_\mathrm{S}`$, and exponentiation extends it to the wider region $`L<1/\alpha _\mathrm{S}`$. This fact can also be argued by comparing the amount of information on the logarithmic terms that is included in the truncation of Eqs. (15) and (16) at some logarithmic accuracy. The reader can easily check that, after matching to the NLO (LO) calculation as in Eq. (11), the NLL (LL) result of Eq. (15) contains all the logarithms of the first four (two) towers in Eq. (16) (and many more logarithmic terms).
In the case of $`Q_{}`$-distributions, full NLL resummation has been performed for lepton pairs, $`W`$ and $`Z`$ bosons produced by the DY mechanism and for Higgs bosons produced by gluon fusion . Corresponding resummed calculations are discussed in Refs. and references therein.
Threshold logarithms in hadron collisions have been resummed to NLL accuracy for DIS and DY production and for Higgs boson production . Recent theoretical progress regards the extension of NLL resummation to processes produced by LO hard-scattering of more than two coloured partons, such as heavy-quark hadroproduction and leptoproduction , as well as prompt-photon , quarkonium and vector-boson production.
An important feature of threshold resummation is that the resummed soft-gluon contributions regard the partonic cross section rather than the hadronic cross section. This fact has two main consequences: $`i)`$ soft-gluon contributions can be sizeable long before the threshold region in the hadronic cross section is actually approached, and $`ii)`$ the resummation effects typically enhance the fixed-order perturbative calculations.
The first consequence follows from the fact that the evolution of the parton densities sizeably reduces the energy that is available in the partonic hard-scattering subprocess. Thus, the partonic cross section $`\widehat{\sigma }`$ in the factorization formula (2) is typically evaluated much closer to threshold than the hadronic cross section. In other words, the parton densities are strongly suppressed at large $`x`$ (typically, when $`x1`$, $`f(x,\mu ^2)(1x)^\eta `$ with $`\eta 3`$ and $`\eta 6`$ for valence quarks and sea-quarks or gluons, respectively); after integration over them, the dominant value of the square of the partonic centre-of-mass energy $`\widehat{s}=x_1x_2S`$ is therefore substantially smaller than the corresponding hadronic value $`S`$.
The second consequence, which depends on the actual definition of the parton densities, follows from the fact that the resummed contributions are those soft-gluon effects that are left at the partonic level after factorization of the parton densities. After having absorbed part of the full soft-gluon contributions in the customary definitions (for instance, those in the $`\overline{\mathrm{MS}}`$ or DIS factorization schemes) of the parton densities, it turns out that the residual effect in the partonic cross section is positive and tends to enhance the perturbative predictions.
A quantitative illustration of these consequences is given below by discussing top-quark and prompt-photon production. The discussion also shows another relevant feature of NLO+NLL calculations, namely, their increased stability with respect to scale variations.
The effects of soft-gluon resummation on the top-quark production cross sections at hadron colliders have been studied in Refs. . In the case of $`p\overline{p}`$ collisions, the comparison between QCD predictions at NLO and those after NLL resummation is shown in Fig. 10 . At the Tevatron the resummation effects are not very large and the NLO cross section is increased by about $`4\%`$. This had to be expected because the top quark is not produced very close to threshold ($`x=2m_t/\sqrt{S}0.2`$, at the Tevatron). Note, however, that the dependence on the factorization/renormalization scale of the theoretical cross section is reduced by a factor of almost 2 by including NLL resummation. More precisely, the scale dependence $`(\pm 5\%)`$ of the NLO+NLL calculation becomes comparable to that obtained by using different sets of parton densities . Combining linearly scale and parton density uncertainties, the NLO+NLL cross section is $`\sigma _{t\overline{t}}=5.0\pm 0.6`$, with $`m_t=175\mathrm{GeV}`$ and $`\sqrt{S}=1.8\mathrm{TeV}`$ .
At the LHC $`(x=2m_t/\sqrt{S}0.03)`$ the top quark is produced less close to the hadronic threshold than at the Tevatron. However this is compensated by the fact that the gluon channelSince $`f_g`$ is steeper than $`f_q`$ at large $`x`$, partonic cross sections in gluon subprocesses are typically closer to threshold than in quark subprocesses. Moreover, the intensity of soft-gluon radiation from gluons is larger than that from quarks by a factor of $`C_A/C_F2`$. is more important at the LHC. As a result, the effect of including soft-gluon resummation to NLL accuracy is very similar: the NLO cross section is enhanced by $`5\%`$ and its scale dependence is reduced from $`\pm 10\%`$ to $`\pm 5\%`$. Note, however, that the uncertainty $`(\pm 10\%)`$ coming from the parton (gluon) densities is larger than at the Tevatron .
Similar qualitative results are obtained when NLL resummation is applied to prompt-photon production at fixed-target experiments. The scale dependence of the theoretical calculation is highly reduced and the resummed NLL contributions lead to large corrections at high $`x_T=2E_T/\sqrt{S}`$ (and smaller corrections at lower $`x_T`$). Of course, the impact of soft-gluon resummation is quantitatively more sizeable in prompt-photon production than in top-quark production, because $`x_T`$ can be as large as 0.6, the hard scale $`E_T`$ is much smaller than $`m_t`$ (thus, $`\alpha _\mathrm{S}(E_T)>\alpha _\mathrm{S}(m_t)`$) and the gluon channel is always important. The scale dependence of the theoretical cross section for the E706 kinematics is shown in Fig. 11. Fixing $`\mu _R=\mu _F=\mu `$ and varying $`\mu `$ in the range $`E_T/2<\mu <2E_T`$ with $`E_T=10\mathrm{GeV}`$, the cross section varies by a factor of $`6`$ at LO, by a factor of $`4`$ at NLO and by a factor of $`1.3`$ after NLL resummation. The highly reduced scale dependence of the NLO+NLL cross section is also visible in Fig. 12, which, in particular, shows that when $`E_T=10\mathrm{GeV}`$ and $`E_{\mathrm{beam}}=530\mathrm{GeV}`$ the central value (i.e. with $`\mu =E_T`$) of the NLO cross section increases by a factor of $`2.5`$ after NLL resummation. As expected, the size of these effects is reduced by increasing $`\sqrt{S}`$ at fixed $`E_T`$ (see Fig. 12) or by decreasing $`E_T`$ at fixed $`\sqrt{S}`$ (see Fig. 8).
The comparison with the E706 data shown in Fig. 13 suggests that the NLO+NLL calculation can help to better understand prompt-photon production at large $`x_T`$. Note, however, that this comparison has to be regarded as preliminary in several respects . In particular, the parton densities used in Fig. 13 are those extracted from NLO fits. Owing to the soft-gluon enhancement at large $`x_T`$, refitting the parton densities may lead to a smaller $`f_g`$ at large $`x`$ and, consequently (because of the momentum sum rule), a larger $`f_g`$ at intermediate $`x`$. As a result, this procedure could somehow increase the theoretical cross section also at smaller values of $`x_T`$.
Soft-gluon resummation at NLL accuracy is now available for all the processes (namely, DIS, DY and prompt-photon production) that are typically used to perform global fits to parton densities. A detailed extraction/evolution of parton densities by consistently using NLL resummed calculations is thus nowadays feasible.
## 6 Other topics
The activity of the QCD Working Group at this Workshop has also been devoted to other topics, such as automatic computation of matrix elements and LO cross sections for multiparticle processes at high-energy colliders, definition and properties of jets algorithms, definition of isolated photons and related NLO calculations. Corresponding contributions are included in these Proceedings.
Other studies performed during this Workshop have a large overlap with the activity of the related Workshops at FERMILAB and CERN and can be found in those Proceedings .
Acknowledgements. I would like to thank the members of the Local Organizing Committee for the excellent Workshop.
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# Third order renormalization group applied to the attractive one–dimensional Fermi gas
## I Introduction
The problem of the one–dimensional Fermi gas model of a metallic conductor, in the low energy approximation, has been approached using three methods: conventional many–body techniques and, mainly, the bosonization and the renormalization group (RG) methods. In this paper we will be concerned with the latter approach. A formulation of the Gell–Mann Low multiplicative RG for this problem was introduced in Ref.. The model considered was the g–ological model, which describes a weakly interacting one–dimensional fermion system with Tomonaga–type ($`g_2`$, $`g_4`$) and backscattering ($`g_1`$) interactions. Phonons are neglected. That method provided a satisfactory understanding of the infrared behavior in the case of a weak repulsive (effective) interaction. A short list of the most relevant results in this case may be the following (for extensive reviews see e.g. Refs.): i) the RG flows toward the Luttinger liquid fixed point; ii) there is a line of nontrivial fixed points; iii) in the infrared limit the system is not asymptotically free, as in the Fermi liquid case, but is described by anomalous exponents. These results were recovered and rigorously proved also in the case of periodic potential using the Wilson RG.
Things change considerably if we consider a weak attractive interaction. Since in this case there is not a second order (one loop) finite fixed point, in Ref. the computation of the beta function was carried to third order (two loops). It was found a $`O(1)`$ third–order fixed point. This result, if reliable, would be of extreme physical interest because it would signal a behavior completely different from the Luttinger liquid paradigm. One should expect the opening of a gap in the dispersion relations, while Luttinger spectrum is gapless, and an exponential decay of the correlation functions, while in the Luttinger case there is only a power low decay with increasing distance. The problem is, of course, how seriously one should take the very existence of an attractive perturbative fixed point on the basis of the third order result. The computation of the the fourth order (three–loop) approximation of the beta function was discussed in Refs.. A smaller but still $`O(1)`$ fixed point was found. Moreover only the first two terms of the beta function are universal. The computation of the third term is useful provided there is some evidence of a perturbatively tractable phase interacting attractively. In this case a precise determination of the renormalized couplings would be important to compute the response functions.
It it useful to make a comparison with the results obtained with the bosonization method. With bosonization it is meant the bosonic representation of fermion field operators. This method is in some sense the inverse of the one used to solve exactly the Luttinger model, where bosonic degrees of freedom are expressed in terms of fermionic operators. Probably the most important result of the bosonization is the exact solution of the model with backscattering ($`g_1`$ and $`g_2`$ terms, see below) in the particular case where $`g_1=\frac{6}{5}\pi `$. Actually the decoupling between charge and spin degrees of freedom, crucial for the exact solution, is open to questions . Moreover there are problems in the limiting procedure employed and the ladder operators restoring the correct occupation numbers are not discussed. (A version of the bosonization free from this problems has been proposed. It should be noted that this version does not deal with the crucial backscattering interaction term: in Ref. only the Luttinger model is considered). Anyway, taking for granted the Luther–Emery solution, the RG method should fill the missing information for value of $`g_1`$ in the neighborhood of the exact solution. From the bosonized representation of the interaction it is not difficult to derive the third order scaling equations and the response functions calculated in Ref. are in good agreement with the results of Ref..
From these considerations one may be tempted to give a heuristic meaning to the large but finite fixed point. In this paper we want to show that this is not the case. The main point is that even the sign of the third order fixed point depends on small variations of a parameter $`\gamma `$ whose value can be arbitrarily chosen, provided $`\gamma >1`$. We will show this both using the Gell–Mann Low (GML) and the Wilson RG.
The paper is organized as follows. In section II we briefly review the multiplicative GML approach. We explain why it is useful to check the results of this approach using other methods. Recasting the multiplicative procedure into discrete steps, instead of considering the usual Lie equation, we reach our main conclusion. In section III we formulate a Callan–Symanzik (CS) approach to the problem and compute the beta function in two–loop approximation. The same computation is proposed in section IV employing the Wilson RG in the multiscale formulation. Finally in section V we come to the conclusions.
## II The Gell–Mann Low approach
We briefly recall the GML multiplicative RG for one dimensional interacting fermions. We will follow closely Refs., with the only difference that we find it convenient to adopt a Euclidean formalism. We consider the g–ological model, defined as follows. The kinetic term is taken linear around the Fermi surface defined by the two points $`k_F`$ and $`k_F`$:
$$H_0=\underset{k,\omega ,\sigma }{}(\omega kk_F)\psi _{k,\omega ,\sigma }^+\psi _{k,\omega ,\sigma }^{},$$
where $`\psi _{k,\omega ,\sigma }^\pm `$ are creation and annihilation operators for right moving $`(\omega =1)`$ and left moving $`(\omega =1)`$ fermions with momentum $`k`$ and spin $`\sigma `$ ($`\sigma =\pm 1/2`$). We choose units such that $`v_F`$=1 ($`v_F`$ is the velocity at the Fermi surface). The ultraviolet (u.v.) stability is imposed by bandwidth cutoffs: the momenta are restricted to the intervals $`(\omega k_Fk_{\text{uv}},\omega k_F+k_{\text{uv}})`$ for $`\psi _{k,\omega ,\sigma }^\pm `$. We define $`E_0=2k_{\text{uv}}`$. The interaction Hamiltonian is
$`H_{\text{int}}=`$ $`{\displaystyle \frac{1}{2L}}{\displaystyle \underset{k,p,\omega ,\sigma ,\sigma ^{}}{}}(g_1\delta _{\sigma ,\sigma ^{}}+g_1\delta _{\sigma ,\sigma ^{}})\psi _{k_1,\omega ,\sigma }^+\psi _{k_2,\omega ,\sigma ^{}}^+\psi _{k_2+2k_F+p,\omega ,\sigma ^{}}^{}\psi _{k_1,2k_Fp,\omega ,\sigma }^{}`$ (1)
$`+`$ $`{\displaystyle \frac{1}{2L}}{\displaystyle \underset{k,p,\omega ,\sigma ,\sigma ^{}}{}}(g_2\delta _{\sigma ,\sigma ^{}}+g_2\delta _{\sigma ,\sigma ^{}})\psi _{k_1,\omega ,\sigma }^+\psi _{k_2,\omega ,\sigma ^{}}^+\psi _{k_2+p,\omega ,\sigma ^{}}^{}\psi _{k_1p,\omega ,\sigma }^{}`$ (2)
$`+`$ $`{\displaystyle \frac{1}{2L}}{\displaystyle \underset{k,p,\omega ,\sigma ,\sigma ^{}}{}}(g_4\delta _{\sigma ,\sigma ^{}}+g_4\delta _{\sigma ,\sigma ^{}})(\psi _{k_1,\omega ,\sigma }^+\psi _{k_2,\omega ,\sigma ^{}}^+\psi _{k_2+p,\omega ,\sigma ^{}}^{}\psi _{k_1p,\omega ,\sigma }^{}.`$ (3)
L is the length of the line. The umklapp interaction term ($`g_3`$) is neglected since it is important only in the half–filled band case, which will be excluded. Since $`g_1=g_2`$ it is always possible to take $`g_2=g_2=g_2`$, reducing the independent couplings to $`g_1`$, $`g_1`$, $`g_2`$. For the sake of simplicity it is possible to neglect, at least as a first approximation, $`g_4`$: we know from the Mattis model that $`g_4`$ does not change the essence of the problem.
In the Euclidean formalism the free propagator in momentum space is given by
$$G_\omega (k)=\frac{1}{ik_0+\omega k_1},$$
(4)
where $`k_0`$ is the energy, $`k_1`$ the momentum (measured from the Fermi surface), $`k=(k_0,k_1)`$ and $`\omega =1(1)`$ for right (left) moving fermions. The renormalization procedure is a prescription that defines new couplings for a theory with a lowered u.v. cutoff $`E_0`$. In the limit $`E_00`$ we obtain the renormalized couplings. If $`G_\omega ^R`$ is the interacting propagator, the $`d`$ function is defined by the relation
$$G_\omega ^R(k)=d(\frac{k_1}{k_{\text{uv}}},\frac{k_0}{E_0})G_\omega (k).$$
The multiplicative constants $`z`$ and $`z_i`$ ($`i=1,1,2`$) that relate $`d`$ and the adimensional vertex functions $`\stackrel{~}{\mathrm{\Gamma }}_i`$ for different values of the cutoff are definite by:
$`d({\displaystyle \frac{k_1}{k_{\text{uv}}^{}}},{\displaystyle \frac{k_0}{E_0^{}}},g^{})=z({\displaystyle \frac{E_0^{}}{E_0}},g)d({\displaystyle \frac{k_1}{k_{\text{uv}}}},{\displaystyle \frac{k_0}{E_0}},g)`$ (5)
$`\stackrel{~}{\mathrm{\Gamma }}_i({\displaystyle \frac{k_j}{k_{\text{uv}}^{}}},{\displaystyle \frac{w_j}{E_0^{}}},g^{})=z_i^1({\displaystyle \frac{E_0^{}}{E_0}},g)\stackrel{~}{\mathrm{\Gamma }}_i({\displaystyle \frac{k_j}{k_{\text{uv}}}},{\displaystyle \frac{w_j}{E_0}},g)`$ (6)
$`g_i^{}=g_iz^2({\displaystyle \frac{E_0^{}}{E_0}},g)z_i({\displaystyle \frac{E_0^{}}{E_0}},g),`$ (7)
where $`E_0^{}<E_0`$ is the lowered cutoff, $`g`$ and $`g^{}`$ denote respectively the old and the new couplings. The invariant couplings $`g_i^R`$ are defined by:
$$g_i^R(\frac{E}{E_0},g)=g_iz^2(\frac{E}{E_0},g)z_i(\frac{E}{E_0},g).$$
The $`g_i^R`$ are invariant in the sense that
$$g_i^R(\frac{E}{E_0^{}},g^{})=g_i^R(\frac{E}{E_0},g).$$
(8)
The couplings $`g_i^{}`$ for the theory with u.v. cutoff $`E_0^{}`$ are defined by
$$g_i^{}=g_i^R(\frac{E_0^{}}{E_0},g).$$
(9)
A differential equation for $`g_i^R`$ is readily derived and is the standard Lie equation:
$$\frac{d}{dx}g_i^R(x,g)=\frac{1}{x}\frac{d}{d\xi }g_i^R(\xi ,g^R(x,g))_{\xi =1}$$
where $`x=E_0^{}/E_0`$. We are interested in the scaling limit $`x0.`$The two–loop result is
$`{\displaystyle \frac{dg_1^R}{dx}}`$ $`={\displaystyle \frac{1}{x}}\left[{\displaystyle \frac{1}{\pi }}g_1^{R\mathrm{\hspace{0.17em}2}}+{\displaystyle \frac{1}{2\pi ^2}}g_1^Rg_1^{R\mathrm{\hspace{0.17em}2}}\right]`$ (10)
$`{\displaystyle \frac{dg_1^R}{dx}}`$ $`={\displaystyle \frac{1}{x}}\left[{\displaystyle \frac{1}{\pi }}g_1^Rg_1^R+{\displaystyle \frac{1}{4\pi ^2}}(g_1^{R2}g_1^R+g_1^{R\mathrm{\hspace{0.17em}3}})\right]`$ (11)
$`{\displaystyle \frac{dg_2^R}{dx}}`$ $`={\displaystyle \frac{1}{x}}\left[{\displaystyle \frac{1}{2\pi }}g_1^{R\mathrm{\hspace{0.17em}2}}+{\displaystyle \frac{1}{4\pi ^2}}g_1^Rg_1^{R\mathrm{\hspace{0.17em}2}}\right].`$ (12)
For spin independent interaction ($`g_1=g_1=g_1`$) the nontrivial fixed point is found for $`g_1^{}=2\pi `$.
We now want to recover this result iterating by discrete steps the procedure that defines the new couplings when the cutoff is lowered: we aim to study the dependence on the scaling parameter. Let $`\gamma >1`$. In proper units we put $`E_0=\gamma ^0`$ and $`g_{i,0}=g_i(E_0)`$ for $`i=1,1,2`$. $`g_{i,1}`$ is defined as (see Eq. (9))
$$g_{i,1}=g_i^R(\frac{\gamma ^1}{\gamma ^0},g_{j,0})=g_i^R(\frac{\gamma ^1}{E_0},g)$$
where $`g_{(0)}=g(E_0)=g`$. The procedure is iterated in the following way: for $`n<0`$ we define
$$g_{i,n1}=g_i^R(\frac{\gamma ^{n1}}{\gamma ^n},g_{j,n})i=1,1,2.$$
From Eqs. (8) e (9) we have that the $`g_{i,n}`$ for $`n=1,2\mathrm{}`$ are the couplings corresponding to the cutoff sequence $`\{\gamma ^n\}`$:
$$g_i^R(\frac{\gamma ^{n1}}{\gamma ^n},g_{j,n})=g_i^R(\frac{\gamma ^{n1}}{\gamma ^0},g_{j,0})=g_i(\gamma ^{n1}).$$
In the limit $`n\mathrm{}`$ we get the renormalized couplings. We have:
$`g_{1,n1}=g_{1,n}{\displaystyle \frac{g_{1,n}^2}{\pi }}\mathrm{ln}\gamma `$ (13)
$`+{\displaystyle \frac{1}{\pi ^2}}g_{1,n}g_{1,n}^2\left(\mathrm{ln}^2\gamma {\displaystyle \frac{1}{2}}\mathrm{ln}\gamma \right)`$ (14)
$`g_{1,n1}=g_{1,n}{\displaystyle \frac{1}{\pi }}g_{1,n}g_{1,n}\mathrm{ln}\gamma `$ (15)
$`+{\displaystyle \frac{1}{2\pi ^2}}\left(g_{1,n}^2g_{1,n}+g_{1,n}^3\right)\left(\mathrm{ln}^2\gamma {\displaystyle \frac{1}{2}}\mathrm{ln}\gamma \right)`$ (16)
$`g_{2,n1}=g_{2,n}{\displaystyle \frac{1}{2\pi }}g_{1,n}^2\mathrm{ln}\gamma `$ (17)
$`+{\displaystyle \frac{1}{2\pi ^2}}g_{1,n}g_{1,n}^2\left(\mathrm{ln}^2\gamma {\displaystyle \frac{1}{2}}\mathrm{ln}\gamma \right).`$ (18)
It is easily checked that in the limit $`\gamma 1^+`$ Eqs. (10) are recovered. In general the fixed point depends on $`\gamma `$ (let’s remember that it is a third order fixed point). For $`\gamma \sqrt{e}`$, if $`g_1=g_1=g_1`$, we have:
$$g_1^{}=\frac{\pi }{(\mathrm{ln}\gamma 1/2)}.$$
(19)
When $`\gamma 1`$ we obtain the previous result $`g_1^{}=2\pi `$. The Lie equations (10) should not be fundamental and we find no reason to use the continuous RG instead of its discrete version. The dependence on $`\gamma `$ will be discussed in section V.
As a final comment on this method we note that Eqs. (5) rely on neglecting small contributes that would not allow one to set multiplicative relations where the $`z`$ factors do not depend on the external momenta. For example in the case of the one–loop approximation of the four–point vertex function, proportional to
$$\frac{1}{2\pi }\mathrm{ln}\left(\frac{k_0}{E_0}\right)+\frac{1}{4\pi }\mathrm{ln}\left(1+\frac{k_0^2}{E_0^2}\right),$$
($`k_0`$ is the external energy, see figure 1), the second term is neglected. That is to say that the vertex functions can be divided in scaling and not scaling terms. The first ones are taken into account while the second are not discussed in . For this reason we find useful to check Eqs. (10) using other methods.
## III The Callan–Symanzik approach
Within the framework of the multiplicative RG, it is not difficult to formulate a Callan–Symanzik approach for our problem. We follow a common procedure: first we renormalize the theory in the u.v. with a fixed (renormalized ) i.r. cutoff $`m`$, then we will compute the beta function and study the i.r. behavior for $`m0`$. This approach, devised for a Field Theory, in our case may be considered unnecessary. Nevertheless we consider it a way to support the GML result.
The i.r. regularized free propagator is defined by inserting a bare mass $`m_0`$ in the propagator (4):
$$G_\omega (k,m_0^2)=\frac{ik_0+\omega k_1}{k^2+m_0^2},$$
where $`k^2=k_0^2+k_1^2`$ and again $`\omega =1(1)`$ for right (left) moving fermions. We know that the Luttinger model with a local interaction is not renormalizable in the u.v. (this is also seen from the exact solution ). In order to impose the u.v. stability we choose a nonlocal interaction whose strength decreases with increasing distance. The interaction Hamiltonian of the model is
$`H_{\text{int}}=`$
$`{\displaystyle \underset{\omega ,\sigma =\sigma ^{}}{}}{\displaystyle d^2xd^2y\psi _{x,\omega ,\sigma }^+\psi _{y,\omega ,\sigma ^{}}^+V_1(xy)\psi _{y,\omega ,\sigma ^{}}^{}\psi _{x,\omega ,\sigma }^{}}`$
$`{\displaystyle \underset{\omega ,\sigma \sigma ^{}}{}}{\displaystyle d^2xd^2y\psi _{x,\omega ,\sigma }^+\psi _{y,\omega ,\sigma ^{}}^+V_1(xy)\psi _{y,\omega ,\sigma ^{}}^{}\psi _{x,\omega ,\sigma }^{}}`$
$`+`$ $`{\displaystyle \underset{\omega ,\sigma ,\sigma ^{}}{}}{\displaystyle d^2xd^2y\psi _{x,\omega ,\sigma }^+\psi _{y,\omega ,\sigma ^{}}^+V_2(xy)\psi _{y,\omega ,\sigma ^{}}^{}\psi _{x,\omega ,\sigma }^{}}`$
$`+`$ $`{\displaystyle \underset{\omega ,\sigma \sigma ^{}}{}}{\displaystyle d^2xd^2y\psi _{x,\omega ,\sigma }^+\psi _{y,\omega ,\sigma ^{}}^+V_4(xy)\psi _{y,\omega ,\sigma ^{}}^{}\psi _{x,\omega ,\sigma }^{}},`$
where $`\psi _{x,\omega ,\sigma }^\pm `$ are the fermion field operators in coordinate space. The potentials $`V_i`$ may be chosen for instance as follows:
$$V_i(x)=\frac{g_i}{4}pe^{p|x_1|}\delta (x_0),i=1,1,2,4,$$
(20)
where $`p>0`$ is fixed; $`x_0`$ and $`x_1`$ are the time and space coordinates. In momentum space the model is the same as (1) with the only difference that the bandwidth cutoffs are replaced by the nonlocal couplings
$$g_ig_i\frac{p^2}{k_1^2+p^2},$$
where $`k_1`$ is the exchanged momentum in the given interaction vertex and $`p`$ is introduced in Eq. (20). In the limit $`p\mathrm{}`$ we recover the local couplings of Eq. (1).
Fermion loops are logarithmically divergent. The theory is regularized introducing a cutoff $`\mathrm{\Lambda }`$ by the means of the standard Schwinger parametrization:
$$\frac{1}{k^2+m_0^2}=_0^{\mathrm{}}𝑑\alpha e^{\alpha (k^2+m_0^2)}_{\mathrm{\Lambda }^2}^{\mathrm{}}𝑑\alpha e^{\alpha (k^2+m_0^2)}.$$
In order to renormalize the theory we find it convenient to follow the scheme for the local ($`p=\mathrm{}`$) case, even if when $`p`$ is finite we make more subtractions than strictly necessary. It is a simple exercise of standard power counting to find the superficial degree of divergence $`D`$ for the $`n`$–point vertex functions in the local case:
$$D(\mathrm{\Gamma }_n)=2\frac{n}{2}.$$
We renormalize the couplings ($`g_ig_i^R`$), the mass ($`m_0m`$) and the wave function ($`\psi ^\pm \psi ^{\pm R}`$). The multiplicative constant $`Z`$ is formally introduced by the relation $`\psi ^\pm =Z^{1/2}\psi ^{\pm R}`$. Let $`\mathrm{\Gamma }_n^R`$ be the renormalized proper $`n`$–point vertex functions. The relation between bare and renormalized vertex functions is:
$$\mathrm{\Gamma }_n(𝐪,m_0,g,\mathrm{\Lambda })=Z^{n/2}\mathrm{\Gamma }_n^R(𝐪,m,g^R),$$
(21)
where $`𝐪`$ denotes the $`n1`$ independent external momenta of $`\mathrm{\Gamma }_n`$ and $`g=\{g_1,g_1,g_2,g_4\}`$. Of course the $`\mathrm{\Gamma }_n`$ are functions of the spin and $`\omega `$ indices attached to the external fields. For simplicity we have not indicated this explicitly in Eq. (21). For the four–point functions the different possible cases are labeled by a single index $`i=1,1,2,4`$. It should be noted that while $`\mathrm{\Gamma }_{4,i}`$ do not depend on the value of $`\omega `$, $`\mathrm{\Gamma }_2`$ does, so we need an $`\omega `$ label for this vertex function. It proves useful to introduce the reduced two–point vertex function $`\widehat{\mathrm{\Gamma }}_2(k)`$:
$$\widehat{\mathrm{\Gamma }}_2(k)=(ik_0+\omega k_1)\mathrm{\Gamma }_{2,\omega }(k).$$
(22)
In the local case $`\widehat{\mathrm{\Gamma }}_2(k)`$ does not depend on $`\omega `$. In the non local case the non vanishing terms in the infrared limit will be $`\omega `$ independent, so we will neglect the $`\omega `$ dependence of $`\widehat{\mathrm{\Gamma }}_2(k)`$. The normalization conditions, which define $`g_i^R`$, $`m`$ and the finite part (zero–loop term) of $`Z`$, are:
$`\mathrm{\Gamma }_{4,i}^R(\mathrm{𝟎})`$ $`=g_i^R`$ (24)
$`\widehat{\mathrm{\Gamma }}_2^R(0)`$ $`=m^2`$ (25)
$`{\displaystyle \frac{1}{2k_0}}{\displaystyle \frac{}{k_0}}\widehat{\mathrm{\Gamma }}_2^R(k)|_{k=(0,0)}`$ $`=1.`$ (26)
The CS equations are derived considering insertions of operators related to the derivatives of the vertex functions respect to the i.r. cutoff $`m_0`$. To this end we introduce the operator $`O`$:
$$O(z)=\underset{\omega ,\sigma }{}\frac{d^2x}{2\pi }\frac{\psi _{x,\omega ,\sigma }^+\psi _{z,\omega ,\sigma }^{}}{x_0z_0i\omega (x_1z_1)}.$$
The corresponding source term in the action has the form $`d^2xv(x)O(x)`$ with $`[v]=2`$. In momentum space the operator $`O`$ is
$$\stackrel{~}{O}(q)=\underset{\omega ,\sigma }{}\frac{d^2k}{(2\pi )^2}\frac{\psi _{k+q,\omega ,\sigma }^+\psi _{k,\omega ,\sigma }^{}}{[i(k_0+q_0)+\omega (k_1+q_1)]},$$
where $`\psi _{k,\omega ,\sigma }`$ are the field operators in momentum space, and $`q`$ is the external momentum of the inserted $`\stackrel{~}{O}`$ operators. When $`O`$ is inserted in a vertex function $`\mathrm{\Gamma }_n`$, the value of $`D`$ for $`\mathrm{\Gamma }_{n,O}`$ is
$$D(\mathrm{\Gamma }_{n,O})=2\frac{n}{2}+([O]2)=\frac{n}{2}.$$
The previous relation means that no new u.v. divergences appear due to $`O`$ insertions (we do not consider $`\mathrm{\Gamma }_{0,O}`$. We remember that the vertex functions with insertions are defined as usual by the Legendre transformation on the field source only and not on the source of the inserted operators). Since $`O`$ does not introduce new divergences, we have $`O=ZO^R`$. The insertion of $`s`$ operator $`O`$ in a vertex function $`\mathrm{\Gamma }_n`$ will be denoted with $`\mathrm{\Gamma }_{n,s}`$. In analogy with Eq. (22) we define $`\widehat{\mathrm{\Gamma }}_{2,s}(q,𝐤)`$ and $`\widehat{\mathrm{\Gamma }}_{2,s}^R(q,𝐤)`$, where $`𝐤`$ denotes the $`s`$ external momenta of the $`s`$ inserted $`\stackrel{~}{O}`$ operators. Equation (21) generalizes into
$$\mathrm{\Gamma }_{n,s}(𝐪,𝐤,m_0,g,\mathrm{\Lambda })=Z^{n/2}Z^s\mathrm{\Gamma }_{n,s}^R(𝐪,𝐤,m,g^R).$$
(27)
It is easily deduced that
$$\frac{}{m_0^2}\mathrm{\Gamma }_n(𝐪)=\mathrm{\Gamma }_{n,1}(𝐪,0),$$
(28)
where the insertion of $`\stackrel{~}{O}`$ in the r.h.s. is made at zero momentum, as indicated. From (27), (28) and (22) we have
$`\widehat{\mathrm{\Gamma }}_{2,1}^R(0)=1`$ (29)
$`m^2=Zm_0^2,`$ (30)
where the second relation follows from $`\widehat{\mathrm{\Gamma }}_2^R(0)=Z\widehat{\mathrm{\Gamma }}_2(0)`$. From Eq. (28) we have:
$$m\frac{}{m}\mathrm{\Gamma }_n(𝐪,m_0,g,\mathrm{\Lambda })|_{g,\mathrm{\Lambda }}=m\frac{m_0^2}{m}|_{g,\mathrm{\Lambda }}\mathrm{\Gamma }_{n,1}(𝐪,0,m_0,g,\mathrm{\Lambda }).$$
(31)
From equations (31) and (27), with the definitions
$`\gamma _1`$ $`={\displaystyle \frac{1}{Z}}m{\displaystyle \frac{Z}{m}}|_{g,\mathrm{\Lambda }}`$
$`\beta _i`$ $`=m{\displaystyle \frac{g_i^R}{m}}|_{g,\mathrm{\Lambda }},`$
($`\gamma _1`$, defined in the previous equations should not be confused with the RG rescaling factor $`\gamma `$) we obtain:
$`\left(m{\displaystyle \frac{}{m}}+{\displaystyle \underset{i}{}}\beta _i(g^R){\displaystyle \frac{}{g_i^R}}{\displaystyle \frac{n}{2}}\gamma _1\right)\mathrm{\Gamma }_n^R(𝐪,m,g^R)=`$ (32)
$`Z(m{\displaystyle \frac{m_0^2}{m}}|_{g_0,\mathrm{\Lambda }})\mathrm{\Gamma }_{n,1}^R(𝐪,0,m,g^R).`$ (33)
It is easy to eliminate any reference to the bare theory. From Eq. (32) written for $`n=2`$ we have:
$`\left(m{\displaystyle \frac{}{m}}+{\displaystyle \underset{i}{}}\beta _i(g^R){\displaystyle \frac{}{g_i^R}}\gamma _1\right)\widehat{\mathrm{\Gamma }}_2^R(q,m,g^R)`$
$`Z(m{\displaystyle \frac{m_0^2}{m}}|_{g_0,\mathrm{\Lambda }})\widehat{\mathrm{\Gamma }}_{2,1}^R(q,0,m,g^R).`$
From Eqs. (25) and the first of (29) we conclude
$$Z(m\frac{m_0^2}{m}|_{g_0,\mathrm{\Lambda }})=(2\gamma _1)m^2,$$
so that Eq. (32) can be written as
$`\left(m{\displaystyle \frac{}{m}}+{\displaystyle \underset{i}{}}\beta _i(g^R){\displaystyle \frac{}{g_i^R}}{\displaystyle \frac{n}{2}}\gamma _1\right)\mathrm{\Gamma }_n^R(𝐪,m,g^R)=`$
$`\left(2\gamma _1\right)m^2\mathrm{\Gamma }_{n,1}^R(𝐪,0,m,g^R).`$
The generalization of Eq. (32) to the case of $`s`$ insertions is immediate:
$`\left(m{\displaystyle \frac{}{m}}+{\displaystyle \underset{i}{}}\beta _i(g^R){\displaystyle \frac{}{g_i^R}}+\left({\displaystyle \frac{n}{2}}+s\right)\gamma _1\right)\mathrm{\Gamma }_{n,s}^R(𝐪,𝐤,m,g^R)`$
$`=\left(2\gamma _1\right)m^2\mathrm{\Gamma }_{n,s+1}^R(𝐪,𝐤,0,m,g^R).`$
Having set the general definitions and relations of the CS approach we can proceed. We will limit ourselves to the computation of the beta function, which is our problem. The normalization conditions (22) fix the zero–loop terms in the loop–wise expansion of $`g_i`$, $`m_0`$ and $`Z`$:
$`g_i^R`$ $`=g_i^{(0)}`$
$`m^2`$ $`=m_0^{2(0)}`$
$`Z^{(0)}`$ $`=1.`$
One–loop calculations are easily done. Of course $`m_0^{2(1)}=0`$ and $`Z^{(1)}=0`$. This implies that up to one loop $`\gamma _1=0`$ (that is $`\gamma _1=O(g^2)`$). It is convenient to write down the results for the couplings in terms of $`g_4`$, $`g_1`$, $`g_2`$ and $`\stackrel{~}{g}g_2g_1`$. We find:
$`g_{1,}^{(1)}`$ $`={\displaystyle \frac{1}{\pi }}\left[\mathrm{ln}m\left(\mathrm{ln}2+\mathrm{ln}p1\right)\right]g_1^Rg_2^R`$ (35)
$`+{\displaystyle \frac{1}{\pi }}\left[\mathrm{ln}m\left(\mathrm{ln}2+\mathrm{ln}p{\displaystyle \frac{1}{2}}\right)\right]g_1^R\stackrel{~}{g}^R{\displaystyle \frac{1}{2\pi }}g_1^Rg_4^R`$
$`g_2^{(1)}`$ $`={\displaystyle \frac{1}{2\pi }}\left[\mathrm{ln}m(\mathrm{ln}2+\mathrm{ln}p1)\right]g_1^{R\mathrm{\hspace{0.17em}2}}`$ (36)
$`\stackrel{~}{g}^{(1)}`$ $`={\displaystyle \frac{1}{2\pi }}\left[\mathrm{ln}m\mathrm{ln}\mathrm{\Lambda }+{\displaystyle \frac{1}{2}}(1+𝑪+\mathrm{ln}2)\right]g_1^{R\mathrm{\hspace{0.17em}2}}`$ (37)
$`g_4^{(1)}`$ $`={\displaystyle \frac{1}{4\pi }}g_1^{R\mathrm{\hspace{0.17em}2}},`$ (38)
where $`𝑪`$ is the Euler constant. In Eq. (35) we note the presence of $`p`$. However simply on the basis of dimensional analysis we can exclude that $`p`$ will appear in the final result. Two–loop calculations are tedious and we will omit the details. We calculate only the singular terms in $`m`$ since we do not plan to go beyond the two–loop approximation. We report the results for $`Z`$ and $`g_1`$:
$$Z^{(2)}=\frac{1}{2^3\pi ^2}\mathrm{ln}m\left(g_1^{R\mathrm{\hspace{0.17em}2}}+g_1^{R\mathrm{\hspace{0.17em}2}}+2g_2^{R\mathrm{\hspace{0.17em}2}}2g_1^Rg_2^R\right)$$
(39)
$`g_{1,}^{(2)}={\displaystyle \frac{1}{2\pi ^2}}\left[\mathrm{ln}^2m2\mathrm{ln}m\left(\mathrm{ln}2+\mathrm{ln}p1\right)\right]g_1^Rg_2^{R\mathrm{\hspace{0.17em}2}}`$ (40)
$`+\left\{{\displaystyle \frac{1}{2\pi ^2}}\mathrm{ln}m{\displaystyle \frac{1}{\pi ^2}}\left[\mathrm{ln}^2m2\mathrm{ln}m\left(\mathrm{ln}2+\mathrm{ln}p{\displaystyle \frac{3}{4}}\right)\right]\right\}g_1^Rg_2^R\stackrel{~}{g}^R`$ (41)
$`+{\displaystyle \frac{1}{2\pi ^2}}\mathrm{ln}mg_1^Rg_2^Rg_4^R{\displaystyle \frac{1}{2\pi ^2}}\mathrm{ln}mg_1^R\stackrel{~}{g}^Rg_4^R+{\displaystyle \frac{1}{2\pi ^2}}\left[\mathrm{ln}^2m2\mathrm{ln}m\left(\mathrm{ln}2+\mathrm{ln}p{\displaystyle \frac{3}{4}}\right)\right]g_1^{R\mathrm{\hspace{0.17em}3}}`$ (42)
$`+{\displaystyle \frac{1}{2\pi ^2}}\left[\mathrm{ln}^2m2\mathrm{ln}m\left(\mathrm{ln}2+\mathrm{ln}p{\displaystyle \frac{1}{2}}\right)\right]g_1^R\stackrel{~}{g}^{R\mathrm{\hspace{0.17em}2}}{\displaystyle \frac{1}{2^2\pi ^2}}\mathrm{ln}m\left(g_1^{R\mathrm{\hspace{0.17em}2}}+g_1^{R\mathrm{\hspace{0.17em}2}}+2g_2^{R\mathrm{\hspace{0.17em}2}}2g_1^Rg_2^R\right)`$ (43)
From Eqs. (40) and (35) we derive $`\beta _1^{(2)}(g^R)`$. The final result is
$$\beta _1(g^R)=\frac{1}{\pi }g_1^Rg_1^R+\frac{1}{4\pi ^2}(g_1^{R\mathrm{\hspace{0.17em}2}}g_1^R+g_1^{R\mathrm{\hspace{0.17em}3}})+O(g^4).$$
(44)
It can be noted that $`p`$ does not appear in Eq. (44), as expected. A crucial use in deriving Eq. (44) is made of Eq. (39), which is responsible for the cancellations expected from the exact solutions of the Luttinger and Mattis models , and for the anomalous behavior of the theory. Of course the anomalous exponent $`\eta =\gamma _1(g^{})`$ derived from Eq. (39) when $`0<g_11`$ is in agreement with the exact solution of the Luttinger model, where $`g_1=0`$, $`g_4=0`$.
Equation (44) is the same as Eq. (10) for the $`g_1`$ coupling. For $`g_2`$ and $`g_1`$ the same conclusion holds: the beta function of the GML method is recovered. The present CS approach, admittedly too involved, has perhaps the only value in that no use is made of approximate multiplicative relations.
## IV The Wilson approach
The multiscale formulation of the Wilson RG is particularly well suited to study the running of the coupling constants by discrete steps. The application of this method to interacting one–dimensional fermionic systems started with Refs. and was thoroughly developed and applied to various problems. Here we give a short and simplified account of the method and refer to the cited papers for the details.
In the coordinate space the free propagator $`G_\omega (x)`$ for $`\omega `$ particles (again $`\omega =\pm 1`$ and $`v_F=1`$) is:
$$G_\omega (x)=\frac{1}{(2\pi )^2}𝑑k_0𝑑k_1\frac{e^{i(k_0x_0+k_1x_1)}}{ik_0+\omega k_1}.$$
Actually it is not necessary to start with a kinetic term linearized around the Fermi surface: the RG can deal with realistic quadratic dispersion relations . This simplification is however inessential for our purposes. Let $`p`$ be an arbitrary momentum scale which for instance may be chosen equal to the inverse of the range of the potential. The propagator is decomposed in the sum
$$G_\omega (x)=\underset{h=\mathrm{}}{\overset{1}{}}G_\omega ^{(h)}(x),$$
with
$`G_\omega ^{(1)}(x)`$ $`={\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle 𝑑ke^{ikx}\frac{1e^{p^2(k_0^2+k_1^2)}}{ik_0+\omega k_1}},`$ (45)
$`G_\omega ^{(h)}(x)`$ $`={\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle 𝑑k\frac{e^{ikx}}{ik_0+\omega k_1}}`$ (47)
$`\times \left[e^{p^2\gamma ^{2h}(k_0^2+k_1^2)}e^{p^2\gamma ^{2h+2}(k_0^2+k_1^2)}\right],`$
where $`h0`$, $`\gamma >1`$ and $`kx=k_0x_0+k_1x_1`$. This decomposition divides the u.v. from the i.r. singularity of the propagator: $`G_\omega ^{(1)}`$ is singular in the u.v. while $`_{h=\mathrm{}}^0G_\omega ^{(h)}=G_\omega ^{i.r.}`$ is singular in the i.r. It is important to note that $`G_\omega ^{(h)}(x)`$, the propagator on scale $`h`$, for $`h0`$ has an u.v. and an i.r. cutoff: $`G_\omega ^{(h)}(x)`$ is essentially different from $`0`$ only for $`x\gamma ^h`$ ($`k\gamma ^h`$ in momentum space).
One imagines that this decomposition stems from a similar decomposition of the fields:
$$\psi _{x,\omega ,\sigma }^\pm =\underset{h=\mathrm{}}{\overset{1}{}}\psi _{x,\omega ,\sigma }^{\pm (h)}$$
such that the pairings in the Grassmannian Wick rule are
$`{\displaystyle P(d\psi _\omega ^{(h)})}`$ $`\psi _{x,\omega ,\sigma }^{+(h)}\psi _{y,\omega ^{},\sigma ^{}}^{(h^{})}\psi _{x,\omega ,\sigma }^{+(h)}\psi _{y,\omega ^{},\sigma ^{}}^{(h^{})}`$
$`\delta _{\omega ,\omega ^{}}\delta _{\sigma ,\sigma ^{}}\delta _{h,h^{}}G_\omega ^h(xy).`$
We are interested to study the i.r. effective potential $`V^{(0)}`$ arising from the integration of the u.v. component $`\psi _\omega ^{(1)}`$ from the effective potential $`V_{\text{eff}}(\phi )`$ defined by:
$$e^{V_{\text{eff}}(\phi )}=\frac{1}{𝒩}P(d\psi )e^{V(\psi +\phi )},$$
where $`𝒩`$ is a normalization constant and $`V`$ is the interaction potential. The ultraviolet integration was actually performed for the spinless model. In the following we suppose to start directly with $`V^{(0)}`$.
The core of the method consists of a procedure that, integrating out the fields from the higher to the lower scales $`h`$ ($`h\mathrm{}`$), constructs a well defined dynamical system of running coupling constants $`g_h`$, whose iteration map is the beta functional.
The operators $``$ and $`=1`$ are introduced. $``$ is the usual renormalization operator of the BPHZ scheme: its action on a given vertex $`\mathrm{\Gamma }`$, in momentum space for instance, is given by $`(\mathrm{\Gamma })=\mathrm{\Gamma }t^\mathrm{\Gamma }(\mathrm{\Gamma })`$, where $`t^\mathrm{\Gamma }`$ denotes the Taylor series with respect to the external momenta of $`\mathrm{\Gamma }`$ up to order $`D(\mathrm{\Gamma })`$, if $`D(\mathrm{\Gamma })`$ is the $`\mathrm{\Gamma }`$ superficial degree of divergence. Of course $`(\mathrm{\Gamma })=t^\mathrm{\Gamma }(\mathrm{\Gamma })`$.
The couplings $`g_h`$ on a given scale $`h`$ are defined by an inductive scheme. Let us assume we have constructed the effective potential $`V^{(h)}(\psi _\omega ^{(h)},g_{h+1},\mathrm{},g_0)`$ on scale $`h`$, where $`\psi _\omega ^{(h)}=_{nh}\psi _\omega ^{(n)}`$ and $`g_{h+1},\mathrm{},g_0`$ are the previously defined couplings on scales $`h+1,\mathrm{},0`$, . We define
$$\overline{V}^{(h)}(\psi _\omega ^{(h)},g_h)V^{(h)}(\psi _\omega ^{(h)},g_{h+1},\mathrm{},g_0).$$
The previous relation introduces the $`g_h`$ and relates them to the $`g_{h+1},\mathrm{},g_0`$ through the beta functional $`B_h`$: $`g_h=g_{h+1}+B_h(g_{h+1},\mathrm{},g_0)`$. The effective potential $`V^{(h1)}`$ on scale $`h1`$ is defined by
$`e^{V^{(h1)}(\psi ^{(h1)})}`$ (48)
$`{\displaystyle \frac{1}{𝒩^{}}}{\displaystyle P(d\psi ^{(h)})e^{V^{(h)}(\psi ^{(h)})V^{(h)}(\psi ^{(h)})}}.`$ (49)
Of course $`V^{(h1)}=V^{(h1)}(\psi _\omega ^{(h1)},g_h,\mathrm{},g_0)`$. The procedure is then iterated. The starting point is given by the couplings $`g_0`$ of $`V^{(0)}`$. The final goal is to find a region in the space of parameters $`g_0`$ where each initial value generates a trajectory $`g_h=g_{h+1}+B_h(g_{h+1},\mathrm{},g_0)`$ such that the Schwinger functions are analytic in the $`g_h`$.
Unfortunately this scheme in our problem requires emendation. From the second order result it becomes clear that $`\alpha _h`$ and $`\zeta _h`$ grow too fast independently on the initial conditions. The point is that we know that the interacting propagator has an anomalous behavior: asymptotically for large distances it decays faster than the free propagator. The wavefunction renormalization necessary to cure this problem is accomplished by an inductive procedure that redefines step by step the free measure of the functional integral and the couplings by the means of a sequence of parameters $`Z_h`$ with $`h=0,1,\mathrm{}`$. Let us assume we have introduced $`Z_h,Z_{h+1},\mathrm{},Z_0`$ and applied our procedure integrating out the scales from $`0`$ to $`h+1`$ ($`h<0`$). We get an effective potential $`\stackrel{~}{V}^{(h)}`$ (different from $`V^{(h)}`$, defined by Eq. (48)). We denote with $`P_{Z_h}(\psi _\omega ^{(h)})`$, $`P_{Z_h}(\psi _\omega ^{(h1)})`$ and $`\stackrel{~}{P}_{Z_h}(\psi _\omega ^{(h)})`$ the free measures with propagators, respectively, $`G_\omega ^{(h)}/Z_h`$, $`G_\omega ^{(h1)}/Z_h`$ and $`\stackrel{~}{G}_\omega ^{(h)}/Z_h`$, where the last one is the modified propagator on scale $`h`$ and $`G_\omega ^{(h1)}=_{ih1}G_\omega ^{(i)}`$. $`\widehat{V}^{(h1)}`$ is defined by
$`{\displaystyle P_{Z_h}(d\psi _\omega ^{(h1)})e^{\widehat{V}^{(h1)}(\sqrt{Z_h}\psi _\omega ^{(h1)})}}=`$ (50)
$`{\displaystyle P_{Z_h}(d\psi _\omega ^{(h1)})\stackrel{~}{P}_{Z_h}(d\psi _\omega ^{(h)})e^{\stackrel{~}{V}^{(h)}(\sqrt{Z_h}\psi _\omega ^{(h)})}}.`$ (51)
$`\widehat{V}^{(h1)}(Z_h^{1/2}\psi ^{(h1)})`$ has the form:
$`\widehat{V}^{(h1)}(Z_h^{1/2}\psi _\omega ^{(h1)})=(+)\widehat{V}^{(h1)}(Z_h^{1/2}\psi _\omega ^{(h1)})`$
$`=Z_h\{\nu _{h1}{\displaystyle \underset{\omega ,\sigma }{}}{\displaystyle }{\displaystyle \frac{d^2k}{(2\pi )^2}}\psi _{k,\omega ,\sigma }^{(h1)+}\psi _{k,\omega ,\sigma }^{(h1)}`$
$`+\zeta _{h1}{\displaystyle \underset{\omega ,\sigma }{}}{\displaystyle \frac{d^2k}{(2\pi )^2}\psi _{k,\omega ,\sigma }^{(h1)+}(ik_0)\psi _{k,\omega ,\sigma }^{(h1)}}`$
$`+\alpha _{h1}{\displaystyle \underset{\omega ,\sigma }{}}{\displaystyle }{\displaystyle \frac{d^2k}{(2\pi )^2}}\psi _{k,\omega ,\sigma }^{(h1)+}(\omega k_1)\psi _{k,\omega ,\sigma }^{(h1)}\}+\mathrm{}`$
Now we add and subtract from $`\widehat{V}^{(h1)}`$ the term $`Z_h\zeta _{h1}\psi _{k,\omega ,\sigma }^{+(h1)}(\omega k_1)\psi _{k,\omega ,\sigma }^{(h1)}`$ and insert the term $`Z_h\zeta _{h1}\psi _{k,\omega ,\sigma }^{+(h1)}(ik_0+\omega k_1)\psi _{k,\omega ,\sigma }^{(h1)}`$ in the free measure. Let $`P_{Z_h}^{}(\psi _\omega ^{(h1)})`$ be the measure changed this way. We define $`Z_{h1}=Z_h(1+\zeta _{h1})`$ and write:
$`{\displaystyle P_{Z_h}(d\psi _\omega ^{(h1)})e^{\widehat{V}^{(h1)}(\sqrt{Z_h}\psi _\omega ^{(h1)})}}={\displaystyle P_{Z_h}^{}(d\psi _\omega ^{(h1)})e^{V^{(h1)}(\sqrt{Z_h}\psi _\omega ^{(h1)})}}`$ (52)
$`=`$ $`{\displaystyle P_{Z_{h1}}(d\psi _\omega ^{(h2)})\stackrel{~}{P}_{Z_{h1}}(d\psi _\omega ^{(h1)})e^{V^{(h1)}(\sqrt{Z_h}\psi _\omega ^{(h1)})}}`$ (53)
$`=`$ $`{\displaystyle P_{Z_{h1}}(d\psi _\omega ^{(h2)})\stackrel{~}{P}_{Z_{h1}}(d\psi _\omega ^{(h1)})e^{\stackrel{~}{V}^{(h1)}(\sqrt{Z_{h1}}\psi _\omega ^{(h1)})}}.`$ (54)
In Eq. (52) $`V^{(h1)}`$ is obtained from $`\widehat{V}^{(h1)}`$ dropping the $`\zeta _{h1}`$ term and substituting $`\alpha _{h1}\zeta _{h1}`$ for $`\alpha _{h1}`$. In Eq. (53) $`\stackrel{~}{P}_{Z_{h1}}(d\psi _\omega ^{(h1)})`$ is the free measure with propagator $`\stackrel{~}{G}_\omega ^{(h1)}/Z_{h1}`$ defined such that the remaining part of the free measure is exactly $`P_{Z_{h1}}(d\psi _\omega ^{(h2)})`$. Finally Eq. (54) defines $`\stackrel{~}{V}^{(h1)}`$ and the r.h.s has the same structure of the r.h.s. of Eq. (50) so the procedure may be reiterated. The starting point are $`Z_0=1`$ and $`\stackrel{~}{G}_\omega ^{(0)}=G_\omega ^{(0)}`$. The relations between the couplings of $`\widehat{V}^{(h1)}(\psi _\omega ^{(h1)})`$, $`g_{i,h1}^{}`$ ($`i=1,2,4`$), $`\alpha _{h1}`$ and $`\zeta _{h1}`$, and the couplings of $`\stackrel{~}{V}^{(h1)}(\psi _\omega ^{(h1)})`$, $`g_{i,h1}`$ and $`\delta _{h1}`$, are easily found:
$`\delta _{h1}={\displaystyle \frac{Z_h}{Z_{h1}}}(a_{h1}z_{h1})`$
$`g_{i,h1}=\left({\displaystyle \frac{Z_h}{Z_{h1}}}\right)^2g_{i,h1}^{}i=1,2,4.`$
Of course $`\zeta _{h1}`$ is no more present in $`\stackrel{~}{V}^{(h1)}(\psi _\omega ^{(h1)})`$, but $`Z_{h1}`$ is introduced. The replacement $`\alpha _h\delta _h`$ drastically improves the convergence properties in the limit $`h\mathrm{}`$.
Now we have the full recipe to proceed. Needless to say a crucial use of the linked cluster theorem will be made. For brevity only the calculation for the $`g_1`$ coupling is sketched (we take $`g_1=g_1=g_1`$).
Let $`C_{i,j}`$ denote the loop of figure 2:
$$C_{i,j}=C_{ij}=\frac{d^2k}{(2\pi )^2}G_\omega ^{(i)}(k)G_\omega ^{(j)}(k),$$
(55)
where $`i,j0`$ and $`G_\omega ^{(i)}`$ is the propagator in momentum space on scale $`i`$. It is immediate to verify that the r.h.s. of (55) does not depend on $`\omega `$ nor on the scale $`p`$ (see (47)) and that it is a function only of the difference $`ij`$. In particular $`C_{i,i}=C_0`$, $`i0`$. The second order calculation gives for $`g_1`$:
$$g_{1,h1}=g_{1,h}+\frac{1}{2}K\left(C_{h,h}g_{1,h}^2+2\underset{h<j0}{}C_{h,j}g_{1,j}^2\right).$$
(56)
$`K`$ is a combinatorial factor: $`K=4`$. The previous relations give the second order approximation for the beta functional. Since we are interested in the beta function, in Eq. (56) we write the $`g_{1,j}`$ for $`j>h`$ as functions of $`g_{1,h}`$. This inversion generates a correction to the higher orders, in particular to the third order. We have (up to $`O(g_1^4)`$ terms):
$`g_{1,h1}=g_{1,h}+{\displaystyle \frac{1}{2}}K\left(C_{h,h}+2{\displaystyle \underset{h<j0}{}}C_{h,j}\right)g_{1,h}^2`$ (57)
$`K^2{\displaystyle \underset{h<i0}{}}C_{h,i}\left({\displaystyle \underset{h<ji}{}}C_{j,j}+2{\displaystyle \underset{h<ki}{}}{\displaystyle \underset{k<j0}{}}C_{k,j}\right)g_{1,h}^3.`$ (58)
Using Eq. (47) for $`G_\omega ^{(h)}`$ we find:
$`C_{h,h}+2{\displaystyle \underset{h<j0}{}}C_{h,j}={\displaystyle \frac{1}{2\pi }}[+\mathrm{ln}\gamma \mathrm{ln}(1+\gamma ^{22h})`$ (60)
$`+\mathrm{ln}(1+\gamma ^{2h})]={\displaystyle \frac{1}{2\pi }}\mathrm{ln}\gamma +O(\gamma ^{2h}).`$
From Eqs. (57) and (60) we get the second order discrete beta function, up to $`\gamma ^{2h}`$ terms (let’s remember that $`\gamma >1`$ and that we are interested in the $`h\mathrm{}`$ limit):
$$g_{1,h1}=g_{1,h}\frac{\mathrm{ln}\gamma }{\pi }g_{1,h}^2+O(g^3).$$
The second line of Eq. (57) gives the corrections to the third order result. We find:
$`{\displaystyle \underset{h<i0}{}}C_{h,i}\left({\displaystyle \underset{h<ji}{}}C_{j,j}+2{\displaystyle \underset{h<ki}{}}{\displaystyle \underset{k<j0}{}}C_{k,j}\right)=`$
$`{\displaystyle \frac{1}{2\pi }}{\displaystyle \underset{h<i0}{}}C_{h,i}\left[\mathrm{ln}\left(1+\gamma ^{2i}\right)+(ih)\mathrm{ln}\gamma \right],`$
where the r.h.s. is easily calculated:
$`\left|{\displaystyle \underset{h<i0}{}}C_{h,i}\mathrm{ln}\left(1+\gamma ^{2i}\right)\right|A\gamma ^h{\displaystyle \underset{h<i0}{}}1,`$ (61)
$`{\displaystyle \underset{h<i0}{}}(ih)C_{h,i}={\displaystyle \frac{\mathrm{ln}2}{4\pi }}.`$ (62)
$`A`$ is a constant. We neglect the r.h.s of Eq. (61). Equation (62) too is derived neglecting terms exponentially vanishing with $`h`$. Putting all together we find a correction $`𝐜`$ to the third order given by:
$$𝐜=\frac{2\mathrm{ln}2\mathrm{ln}\gamma }{\pi ^2}g_{1,h}^3.$$
The computation of the third order is tedious. We will limit ourselves to note that exists a contribution $`\mathrm{ln}2\mathrm{ln}\gamma `$ of two–loop diagrams that cancels exactly $`𝐜`$. This is important because there is no such term in Eqs. (10), (13) or (44). It comes from the diagrams $`D_1`$ and $`D_2`$ of figure 3, which are related by $`D_1=2D_2`$. We will consider the simpler $`D_1`$. The contributions to $`D_1`$ given by
$`{\displaystyle \underset{h<i0}{}}C_{h,h}2C_{h,i}+{\displaystyle \underset{h<i0}{}}{\displaystyle \underset{h<j<i}{}}2C_{h,i}2C_{h,j}`$
$`+{\displaystyle \underset{h<i0}{}}{\displaystyle \underset{h<j,k<1}{}}C_{h,i}2C_{j,k}+{\displaystyle \underset{h<i0}{}}{\displaystyle \underset{h<j<i}{}}{\displaystyle \underset{h<kj}{}}2C_{h,j}2C_{k,i}`$
amount to
$$\frac{1}{(2\pi )^2}\mathrm{ln}2\mathrm{ln}\gamma \frac{1}{2\pi }\underset{h<i<0}{}2C_{i,h}\mathrm{ln}\left(1+\gamma ^{2i}\right).$$
(63)
The second term of Eq. (63) can be neglected as for Eq. (61) and the first one gives the desired cancellation (taking into account the combinatorial factors). The final result is the same as Eq. (13) or (44) with $`g_1_{}=g_1=g_1`$. Again we find the fixed point $`g_1^{}`$ of Eq. (19) and we recover Eqs. (10) in the limit $`\gamma 1`$.
## V Conclusions
We have computed the third order (two–loop) approximation of the beta function for a one dimensional model of interacting fermions, aiming in particular to study the case of attractive interaction. An existing result, derived making use of tacitly assumed approximations, pointed out a $`O(1)`$ fixed point. We tried to support this conclusion setting a Callan–Symanzik approach and using the Wilson RG formulated as in Refs.. In each case we recovered the aforementioned result.
An attempt to pursue further the examination of the problem was made in Refs. where the fourth–order approximation, which is by no means universal, was computed. We propose a different approach focused on the study of the dependence on $`\gamma `$, the rescaling factor of the RG group. A similar idea was discussed in Ref., where the dependence of the fixed points on the parameters of the RG was analyzed. Our simple idea is that the stability of the result with respect to $`\gamma `$ should indicate how reliable one should consider the perturbative result. It was expected a third order result dependent on $`\gamma `$ but we found a too strong dependence: the fixed point happens to change sign if $`\gamma >\sqrt{e}`$, which is still $`1`$ (of course taking $`\gamma 1`$ and, at the same time, truncating at the third order would be questionable). Nothing similar happens to the nontrivial fixed points for repulsive interaction $`\{g_1^{}=0,g_2^{}\}`$, which are in some sense insensitive to the value of $`\gamma `$.
Which conclusions can be drawn? It is useful to compare the one dimensional interacting Fermi Gas with the well known Kondo problem. It can be noted that the scaling equation for $`g_1`$ (Fermi Gas) and the one for the impurity coupling $`\lambda `$ (Kondo model) have the same structure (see e.g. Ref. for a review in the modern language of Conformal Field Theory). The Kondo effect was thoroughly investigated. The ferromagnetic case corresponds to the Fermi gas with repulsive interaction ($`g_1>0`$): the RG flow is such that $`\lambda 0`$. If the coupling is antiferromagnetic the system flows toward a strong coupling phase ($`\lambda \mathrm{}`$). Moreover the infrared divergences induce a scale, the Kondo temperature $`T_\text{K}`$, characterizing the low energy physics.
In our case the particular instability of the perturbative result with respect to $`\gamma `$, a dimensionless parameter without a physical meaning, should indicate that the RG flow does not actually stop at a finite value and suggests a conclusion similar to the previous one. In our case the characteristic scale should be a gap $`\mathrm{\Delta }`$ for the spin degrees of freedom, whose expression, according to Ref. , should have for small coupling an expression of the type $`\mathrm{\Delta }\sqrt{g_1}\mathrm{exp}(1/g_1)`$.
###### Acknowledgements.
This work received support from the INFN, section of Cagliari, and from the “Dipartimento di Scienze Fisiche” of Cagliari University. We are deeply indebted with Prof. G. Gallavotti, from whom we learned the multiscale formulation of the RG, for having originated the theme of this research as well as for many discussions. We acknowledge useful discussions with Prof. E. Marinari, Dr. V. Mastropietro, Prof. G. Benfatto and Dr. M. Lissia.
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# Solid like friction of a polymer chain
## I Polymer chain interacting with traps
### A The one trap problem.
We study a Gaussian polymer chain interacting with one trap on a solid surface. For simplicity, we consider the chain as two-dimensional and lying in the surface with a Gaussian conformation (at rest). The end point of the chain $`s=N`$ moves at a velocity $`u`$. When the chain is free (not trapped), its average conformation can be obtained from the Rouse equation 1. The average position $`z_{free}(s,t)`$ of monomer $`s`$ at time $`t`$ is
$$z_{free}(s,t)=\frac{1}{2}\frac{\zeta u}{k}s^2=\frac{1}{2}uT_R\left(\frac{s}{N}\right)^2$$
(2)
The tension of the chain at monomer $`s`$ is given by the friction force on the chain end section containing $`s`$ monomers $`\tau _{ch}(s,t)=k\frac{z_{free}(s,t)}{s}=\zeta us`$. The tension on the free end of the chain vanishes, and the external force applied to pull the chain is $`N\zeta u=\frac{k}{N}uT_R`$. This can be viewed either as the friction force on the N monomers, or as the elastic tension of N springs in series extended over a distance $`uT_R`$. In the absence of trapping, the energy dissipated per unit time by the solvent viscosity is simply $`\frac{d_d}{dt}=N\zeta u^2`$.
When the end monomer (s=0) gets trapped, the chain is progressively elongated by the motion of the pulled end (s=N). The tension applied on the trapped monomer increases until it reaches the critical value $`\tau `$. The trapped chain end is then released and the chain contracts back to its free conformation.
The Rouse equation for the trapped chain allows the calculation of the chain tension (appendix A) :
$$\tau _{ch}(s,t)=k\frac{uT_R}{N}\left[\frac{t}{T_R}+\frac{s}{N}+\frac{2}{\pi ^2}\underset{n=1}{\overset{+\mathrm{}}{}}\frac{1}{n^2}\left(1\mathrm{exp}(n^2\pi ^2\frac{t}{T_R})\right)\mathrm{cos}\left(\pi \frac{n}{N}(sN)\right)\right]$$
(3)
The tensions on the trapped end $`\left(\tau _{tr}\left(t\right)=\tau _{ch}\left(s=0,t\right)\right)`$ and on the pulled end $`\left(\tau _{ch}\left(s=N,t\right)\right)`$ are plotted on figure 2. At short time after the trapping ($`tT_R`$), the tension on the pulled end remains constant (equal to the viscous friction force $`N\zeta u`$), and the tension on the trapped monomer increases as $`\tau _{tr}\zeta Nu\left(\frac{t}{T_R}\right)^{1/2}`$. At larger times, ($`tT_R`$), the tension is uniform along the chain and both tensions increase as $`\zeta Nu\frac{t}{T_R}`$. After the sticking time $`t_s`$, the tension $`\tau _{tr}`$ becomes equal to the critical value $`\tau `$.
The characteristic velocity of the problem is obtained by comparing the friction force on the chain $`\zeta Nu`$ to the maximum force $`\tau `$ that can be exerted by the trap $`u_c=\frac{\tau }{N\zeta }=\frac{N\tau }{kT_R}`$.
If $`uu_c`$, the viscous force on the chain is smaller than $`\tau `$. When the trapped end is released from the trap, the whole chain is stretched with a constant tension and its extension is $`L=\frac{N\tau }{k}`$. The tension acting on the first monomer at time $`t`$ is $`\tau _{tr}(t)\frac{k}{N}ut`$ and the sticking time $`t_s`$ is given by :
$`t_s`$ $`=`$ $`{\displaystyle \frac{u_c}{u}}T_R(uu_ct_sT_R)`$ (4)
The work performed to extend the chain is $`\frac{N\tau ^2}{2k}`$. After the release of the trapped end, the elastic energy stored in the chain $`W_d=\frac{N\tau ^2}{2k}`$ is dissipated by a viscous process. This argument is essentially identical to the classical argument of Lake and Thomas which has been used to study the adhesive properties of polymers by de Gennes .
If $`uu_c`$, the viscous force on the chain is larger than $`\tau `$, and only a small part of the chain involving $`s_0(t)\sqrt{\frac{k}{\zeta }t}`$ monomers is elongated when the release of the trapped end occurs. It contributes to the tension on the trapped monomer as $`\tau _{tr}(t)\frac{k}{s_0(t)}ut`$ and the sticking time $`t_s`$ is now given by :
$`t_s`$ $`=`$ $`\left({\displaystyle \frac{\sqrt{\pi }}{2}}{\displaystyle \frac{u_c}{u}}\right)^2T_R={\displaystyle \frac{\pi \tau ^2}{4ku^2\zeta }}`$ (5)
where the numerical prefactor is obtained from the detailed calculations of the appendix. The force acting on the pulled monomer remains constant and equal to $`N\zeta u`$ during the trapping with a good approximation. After the release of the traped end, the tension of the pulled monomer is calculated in appendix B and is given by (B6). The extra energy dissipated due to the trapping is of the order of the elastic energy stored during the trapping in the deformed part of the chain $`W_d=s_0(t_s)\frac{\tau ^2}{k}=N\frac{\tau ^2}{k}\frac{u_c}{u}`$. More complete calculations confirming these scaling law can be done using the Rouse model. They are presented in appendices A and B. After the release of the anchored monomer, the chain relaxes to the steady state, and its elastic energy is dissipated. The total dissipated energy is
$$W_d=N\zeta u^2t+\frac{4}{3\pi ^{1/2}}N\zeta u^2t_s\left(\frac{t_s}{T_R}\right)^{1/2}$$
(6)
### B The many-traps problem.
We now consider a polymer pulled at a velocity u on a two dimensional surface with a density of traps n. We refer to appendix D for a discussion of the trapping probability that gives the mean distance $`d(u)`$ between trapping events. The problem is then equivalent to a one dimensional array of traps along the velocity direction with an average density $`1/d(u)`$. For simplicity, we consider here a periodic lattice with a distance d(u) between traps. We then calculate the average friction force in a steady state as $`F_{fr}=\frac{W_d}{d}`$ where $`W_d`$ is the energy dissipated on each trap.
We first consider the limit of very low velocities $`uu_c`$. When the chain is trapped on one defect and is just on the verge of detaching, the chain has a constant tension $`\tau `$ and its size is $`L=\frac{N\tau }{k}`$. If this is smaller than the distance $`d`$ between traps, the traps are independent and we can use the results of the previous section. The energy dissipated per trap is $`W_d\frac{N\tau ^2}{k}`$ and the friction force is $`F_{fr}\frac{N\tau ^2}{kd}`$. If $`dN\tau /k`$, the chain end gets trapped before the full relaxation of the tension. The number of monomers where the tension has relaxed increases as $`t^{1/2}`$ and the position of the first monomer is given by $`z(t)\frac{\tau t^{1/2}}{(\zeta k)^{1/2}}`$. The chain gets trapped again when $`z(t)=d`$ or after a time $`t_a\frac{\pi \zeta kd^2}{4\tau ^2}`$. The tension then becomes uniform along the chain over the Rouse time $`T_R`$ and equal to $`\tau ^{}=k(Ld)/N`$. The chain stretches then back to a uniform tension $`\tau `$ over the time $`d/u`$. The dissipated energy per trap is $`W_d=\frac{N\tau ^2}{2k}\frac{N\tau ^2}{2k}`$ and the friction force is $`F_{fr}\tau `$.
This quasi-static description remains valid as long as $`d/u>T_R`$. If this is not the case, the tension does not have time to relax to the smaller value $`\tau ^{}`$ over the whole chain and remains equal to $`\tau `$ (in fact sightly larger) on the pulled monomer. As long as the viscous force is small ($`u<u_c`$) this tension balances the friction force and thus $`F_{fr}\tau `$. In this case, the dissipated energy per trap is of the order of the elastic energy of a chain spanning the distance $`d`$ i.e. containing $`s=kd/\tau `$ monomers.
If the velocity $`u`$ is larger than the critical value $`u_c`$, the total dissipation is dominated by the viscous forces. The defects are independent if $`t_s(u)\frac{d}{u}`$ i.e. if $`d>\frac{\tau ^2}{k\zeta u}`$. In this case the dissipation is due to individual traps and the friction force is $`F_{fr}\frac{\tau ^3}{k\zeta ud}\frac{N\tau ^2}{kd}\frac{u_c}{u}`$. It decreases with the velocity. If $`d\frac{\tau ^2}{k\zeta u}`$, as above, the dissipated energy per trap is of the order of the elastic energy of the chain section spanning the length $`d`$ with a tension $`\tau `$ and the friction force is $`F_{fr}=\tau `$.
All these results are summarized on figure 3 in a velocity-distance between traps $`(u,d)`$ diagram.
## II Thermal fluctuations.
In the previous section, we have considered point-like traps and we have ignored thermal fluctuations describing the chain statistics by the average position of each monomer. In a more realistic model, these two assumptions are relaxed. Each trap is an attractive potential well $`V(zz_0)`$ for the chain end with an extension $`b`$ around the trap position $`z_0`$ and a depth $`U_0`$. The maximum tension that can be exerted by the trap $`\tau `$ is the slope of this potential at the inflection point $`z_1`$. We want to study here strong trapping and we only focus on the case where $`\frac{\tau b}{k_BT}1`$. This can be achieved for a reasonable tension if the size of the trap is larger than the monomer size $`a`$ (typically a few times) but is much smaller than the chain size and the distance between traps $`d`$. A typical trapping potential is sketched on figure 4.a.
Thermal fluctuations have two effects; they can induce a new trapping of a chain end that has just escaped from the trap; they can also cause the release of the end monomer even if the tension is lower than the critical value $`\tau `$.
In references , Wittmer et al and Clement et al have shown that the readsorption of the chain can occur until the terminal monomer leaves the first Pincus blob of size $`\xi =k_BT/\tau _{tr}`$. The size of this Pincus blob in the strong trapping limit is smaller than the size of the trap $`b`$ when the tension is equal to the critical tension. The retrapping effect can then be neglected.
The escape of a particle submitted to an external force from a potential well can be studied using Kramers rate theory. If the constant external force is $`\tau _{tr}`$, particle is submitted to the effective potential $`U(z)=V(z)\tau _{tr}z`$ (choosing the origin at the center of the trap). This potential has a minimum value at a position $`z_{min}`$ corresponding to the equilibrium position of the particle and a maximum value at a position $`z_{max}`$ ($`z_{min}<z_1<z_{max}`$). The escape time due to thermal fluctuations from the potential is the first passage time of the particle at the position $`z_{max}`$. Within Kramers rate theory , this is given by
$$t_e(\tau _{tr})=t_1\mathrm{exp}\left[\left\{U(z_{max})U(z_{min})\right\}/k_BT\right]$$
(7)
where $`t_1`$ is given by the curvature of the potential $`U(z)`$ at the positions $`z_{min}`$ and $`z_{max}`$
$$t_1(\tau )=\frac{2\pi \zeta }{\left(\frac{d^2U}{dz_{max}^2}\frac{d^2U}{dz_{min}^2}\right)^{1/2}}$$
(8)
Our aim here is to use these result to discuss the escape of the trapped end under the action of the time-dependent tension and to calculate the chain tension $`\tau _e`$ when the trapped end escapes from the potential well. The results of the previous section can then be used with the effective critical tension $`\tau _e`$ instead of $`\tau `$.
In reference , two approximations (so-called adiabatic and self-consistent force approximation) are proposed to discuss the escape of a particle from a potential well under the action of a non constant force. We have checked that both lead to very similar results. For a sake of clarity, we present here the simpler self-consistent force approximation. In this approximation, the tension is obtained self-consistently by calculating the escape time from equation 7 and 8 for a tension $`\tau _{tr}`$ increasing with time and by imposing that $`\tau _e=\tau _{tr}(t_e)`$.
In the limit of small velocities $`uu_c`$, the chain tension at the end point increases as $`\tau _{tr}=\frac{kut}{N}`$. The self-consistency equation for the effective critical tension can then be written as :
$$\frac{N\tau _e}{ku}=t_1\mathrm{exp}\left[\left\{U(z_{max})U(z_{min})\right\}/k_BT\right]$$
(9)
where the potential is calculated with the tension $`\tau _e`$. If the velocity is very small, the critical tension $`\tau _e`$ is very small and the escape time can be expanded in powers of the tension. For a vanishing tension, the escape time is $`t_e^0=t_1exp(U_0/k_BT)`$ where $`t_1\zeta b/\tau `$ (assuming that the potential has only one energy scale $`U_0\tau b`$ and one length scale $`bz_1`$). The effective critical tension is then :
$$\tau _e=\frac{kut_e^0}{N}\frac{1}{1+\alpha (\tau b/k_BT)(t_e^0/t_s(u))}$$
(10)
where $`t_s(u)`$ is the sticking time introduced in the previous section and $`\alpha `$ a numerical constant. This result is valid as long as the thermal escape time is small enough $`t_e^0\frac{k_BT}{\tau b}t_s(u)`$ leading to $`\frac{u}{u_c}\frac{u_1}{u_c}(Na/b)^2\mathrm{exp}(U_0/k_BT)`$.
We assume for simplicity that the thermal escape time is larger than the Rouse time of the whole chain. The tension is thus uniform along the whole chain when the escape from the trap arises; the velocity $`u_1`$ is then smaller than $`u_c`$. In this limit of vanishing velocities, if the traps are independents, the friction force is $`F_{fr}=\frac{N\tau _e^2}{2kd}`$ and vanishes as $`u^2`$. If the escape time is smaller than the Rouse time only a fraction of the chain containing $`s(kt_e^0/\zeta )^{1/2}`$ monomers is stretched. The critical tension is $`\tau _e=(k\zeta t_e^0)^{1/2}u`$ and the friction force is $`F_{fr}=\frac{s\tau _e^2}{2kd}`$; it also vanishes as $`u^2`$.
If the velocity is larger than $`u_1`$, the effective tension increases between a value of order $`k_BT/b`$ and $`\tau `$. The variation of the effective critical tension explicitly depends on the shape of the potential and cannot be expressed in simple analytical form. A simple expression can only be given if the velocity is large enough. In this case, the two positions $`z_{max}`$ and $`z_{min}`$ are close to the potential inflection point $`z_1`$. One can then expand the potential around $`z=z_1`$
$$U(z)=U_1+(\tau \tau _e)(zz_1)\frac{1}{3}\gamma \tau z_1\left(\frac{zz_1}{z_1}\right)^3$$
(11)
where $`\gamma `$ is a number of order one. The effective tension can be calculated from this expansion :
$$\tau _e=\tau \left(1\frac{1}{\beta ^{2/3}}\mathrm{log}^{2/3}\left(t_s(u)\frac{\tau _e}{\tau }(1\frac{\tau _e}{\tau })^{1/2}/t_1^0\right)\right)$$
(12)
where we define the dimensionless number $`\beta =(4\tau z_1)/(3k_BT\gamma ^{1/2})`$ and the microscopic time $`t_1^0=(\zeta z_1)/(2\gamma ^{1/2}\tau )`$. This expansion is valid if $`\beta `$ is larger than $`\mathrm{log}(t_s/t_1)`$. Qualitatively the effective tension is equal to the maximal force exerted by the potential if the sticking time $`t_s`$ is smaller than a time of the order of the thermal escape time. The effective tension increases smoothly (with a logarithmic dependence in u) between $`k_BT/b`$ and $`\tau `$. At first order we can use the power law found in the previous section in the framework of the zero temperature theory.
A similar calculation can be made at larger velocities $`uu_c`$. If the thermal escape time is larger than the Rouse time, the effective critical tension is close to $`\tau `$ and a law similar to eq. 12 is found (only the argument of the logarithm is slightly different). One can thus also use the zero temperature result for the friction force.
Finally, if the traps are not independent one should compare the thermal escape time to the sticking time that, in this case, is of order $`d/u`$. The escape of the trap is dominated by thermal fluctuations if the thermal escape time is smaller.
## III Concluding Remarks
In this paper, we have studied a very specific model for the friction of a polymer chain on a solid surface. We have considered a two-dimensional gaussian polymer chain interacting with traps on a solid substrate. All the dissipation in this model is due to viscous friction in the solvent surrounding the chain and there is no dissipation associated to the substrate. Our main result is that, neglecting thermal fluctuations, the traps exert a finite friction force on the chain that does not vanish when the velocity goes to zero. This friction force originates in the elastic energy stored in the chain which is then dissipated by viscous friction in the liquid. At higher velocity the friction force decreases with the velocity due to the fact that the portion of chain that is stretched gets smaller. If thermal fluctuations are taken into account, the friction force formally vanishes at zero velocity. Nevertheless, for most practical cases the traps are strong enough and the effect of thermal fluctuations will appear only at infinitely small velocities. However, a very weak logarithmic dependence of the friction force on the velocity still remains, due to the fact that the trapping cross section depends logarithmically on the velocity.
For a macroscopic sample consisting of a statistical ensemble of independent chains, the static friction coefficient is proportional to the average number of trapped chains. Furthermore, our model gives an interesting result concerning the dependence of this coefficient with the sample history. During a friction experiment at velocity $`u`$, the average fraction of trapped chains is given by $`f\frac{t_s(u)u}{d}`$. For $`uu_c`$ this fraction is constant and equal to $`f_0\frac{N\tau }{kd}`$. For $`uu_c`$ it decreases with the velocity as $`f_u\frac{N\tau }{kd}\frac{u_c}{u}`$. In this case, if the motion is stopped fast enough, the static friction coefficient is lower than $`f_0`$ and increases with time until it reaches $`f_0`$.
Our model may look very simplistic, but we believe that many of the assumptions Πmade here do not have a strong influence on the main results although they may change the quantitative value of the friction force. For example, we have assumed that only the last monomer of the chain (the chain end that is not pulled) interacts with the traps. It is easily checked that identical results are obtained (within a factor $`2`$) if any monomer not too close to the chain end interact with the traps. We have explicitly performed the Rouse calculations in this case, they are slightly more complicated and we do not present them in this paper. Another extension that we would like to make in the future is the case where several monomers of the same chain interact with the traps. There is of course a possibility that at the same time several monomers are trapped; this then leads to correlated release (cascades). We believe however that the same qualitative results would be obtained and that the friction force would remain finite at low velocities. We also have assumed a periodic array of traps. A more realistic approach would be to introduce disorder both in the positions of the traps and in the distance between traps. A weak disorder is clearly not important. If the traps are independent, then the averaging of the dissipation is straightforward. If the average distance between traps is small, the tension of the end monomer remains roughly constant as the tension fluctuations due to the trappings do not propagate to the pulled monomer. Disorder could play an important role if it introduces strong fluctuations (such as rare trapping events with a strong critical tension). It would then be interesting to study the fluctuations in the friction force.
A last limitation is that we have only studied the stationary friction state. Transient effects are obviously very important and must be studied independently.
It is also important to discuss the orders of magnitude of the predicted effects. Our more quantitative discussion of the release from the trap clearly shows that the solid-like friction can only be observed if the energy barrier is large $`\tau b/k_BT1`$ otherwise the whole escape process is thermally activated and there is no solid friction. This could be achieved for example with a trap size of the order of a few monomer sizes $`b10a`$ and $`\tau b/k_BT10`$. The critical velocity for a chain of $`1000`$ monomers is then of the order of $`0.1mm/s`$ i.e. in the accessible range for atomic force microscopy. One must however notice that the chain tensions can then be rather large and that it is not clear that the chain elasticity remains gaussian. This would of course change the quantitative value of the friction force but not the qualitative physics.
In the future, we would like to use our very simple model to study the solid friction of a polymer gel moving on a solid substrate in order to give a quantitative interpretation of the recent experiments of Gong and Osada
## A Trapped chain.
It is convenient to perform all calculations using the Laplace transformation, defined by :
$$\stackrel{~}{z}(s,p)=_0^+\mathrm{}z(s,t)e^{pt}𝑑t$$
(A1)
For the trapped chain, the following Rouse equation has to be solved:
$$\zeta \left(p\stackrel{~}{z}(s,p)z_{free}\right)=k\frac{^2}{s^2}\stackrel{~}{z}(s,p)$$
(A2)
with the boundary conditions :
$`\stackrel{~}{z}(s=0,p)=0`$ (A3)
$`\stackrel{~}{z}(s=N,p)={\displaystyle \frac{u}{p^2}}`$ (A4)
The solution of this equation is given by:
$`\stackrel{~}{z}_d(s,p)`$ $`=`$ $`{\displaystyle \frac{u}{p^2}}\left[1+{\displaystyle \frac{1}{2}}pT_R\left({\displaystyle \frac{s}{N}}\right)^2+{\displaystyle \frac{\mathrm{sinh}\sqrt{T_Rp}\frac{\left(sN\right)}{N}}{\mathrm{sinh}\sqrt{T_Rp}}}\right]`$ (A5)
$`z_d(s,t)`$ $`=`$ $`uT_R\left[{\displaystyle \frac{t}{T_R}}{\displaystyle \frac{s}{N}}+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{s}{N}}\right)^2+{\displaystyle \frac{2}{\pi ^3}}{\displaystyle \underset{n=1}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{(1)^n}{n^3}}\left(1e^{n^2\pi ^2\frac{t}{T_R}}\right)\mathrm{sin}\left(\pi {\displaystyle \frac{n}{N}}(sN)\right)\right]`$ (A6)
and the tension along the chain is given by:
$$\tau _{ch}(s,t)=k\frac{uT_R}{N}\left[\frac{t}{T_R}+\frac{s}{N}+\frac{2}{\pi ^2}\underset{n=1}{\overset{+\mathrm{}}{}}\frac{1}{n^2}\left(1e^{n^2\pi ^2\frac{t}{T_R}}\right)\mathrm{cos}\left(\pi \frac{n}{N}(sN)\right)\right]$$
(A7)
At large time $`tT_R`$, the tension on the trapped monomer behaves as :
$$\tau _{tr}(s=0,t)\frac{k}{N}ut$$
(A8)
which corresponds to the tension of N springs under the extension $`ut`$.
For short times $`tT_R`$, the tension on the trapped monomer behaves as :
$$\tau _{tr}(s=0,t)=\frac{2}{\sqrt{\pi }}\frac{k}{N}uT_R\sqrt{\frac{t}{T_R}}$$
(A9)
As given in the main text, the sticking time is obtained when this tension is equal to the critical value $`\tau `$.
In this case, the number of extended springs is $`s_0(t)N\sqrt{\frac{t}{T_R}}`$.
## B Detachment
After the detachment of the trapped monomer, the chain relaxes to the steady state, and its elastic energy is dissipated. The Rouse equation becomes:
$$\zeta \left(p\stackrel{~}{z}(s,p)z_d(s)\right)=k\frac{^2}{s^2}\stackrel{~}{z}(s,p)$$
(B1)
where $`z_d(s)uT_R\left[\frac{1}{2}\left(\frac{s}{N}\right)^2+\frac{t_s}{T_R}4\frac{t_s}{T_R}i^2erfc\left(\frac{1}{2}\sqrt{\frac{T_R}{t_s}\frac{s}{N}}\right)\right]`$ ($`i^nerfc(x)`$ is a repeated integral of the error function defined in ) gives the conformation of the chain at the time where the chain detaches form the trap . The boundary conditions are:
$`{\displaystyle \frac{\stackrel{~}{z}}{s}}(s=0,p)`$ $`=`$ $`0`$ (B2)
$`\stackrel{~}{z}(s=N,p)`$ $`=`$ $`uT_R\left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{t_s}{T_R}}\right){\displaystyle \frac{1}{p}}+{\displaystyle \frac{u}{p^2}}`$ (B3)
The solution of (B1) is given by:
$$\stackrel{~}{z}(s,p)=Au_+(s,p)+Bu_{}(s,p)\frac{1}{2N}\sqrt{\frac{T_R}{p}}\left(u_{}(s,p)_0^sz_d(s^{})e^{\sqrt{pT_R}\frac{s^{}}{N}}𝑑s^{}u_+(s,p)_0^sz_d(s^{})e^{\sqrt{pT_R}\frac{s^{}}{N}}𝑑s^{}\right)$$
(B4)
where $`u_\pm =e^{\pm \sqrt{pT_R}\frac{s}{N}}`$, $`A`$ and $`B`$ are two constants determined by (B3). The first boundary condition gives $`A=B=\frac{\alpha }{2}`$. In the limit of large times, the second condition gives $`\alpha \frac{\stackrel{~}{z}(N)}{\mathrm{cosh}(\sqrt{pT_R})}`$. At large times, the tension on the pulled monomer is given by:
$`\stackrel{~}{\tau _d}(N,p)`$ $``$ $`N\zeta {\displaystyle \frac{u}{p}}+8i^3erfc(0)N\zeta u\sqrt{{\displaystyle \frac{t_s}{T_R}}}t_s`$ (B5)
$`\tau _d(N,t)`$ $``$ $`N\zeta u+8i^3erfc(0)N\zeta u\sqrt{{\displaystyle \frac{t_s}{T_R}}}t_s\delta (t)`$ (B6)
By integrating this tension over time one finds the dissipated energy given in the main text (equation 6).
## C Hydrodynamic interactions.
The calculations presented in the main text of this paper have been made using the Rouse dynamics for the polymer chains on the solid surface. This ignores two effects : the direct friction of the chains on the solid surface and the hydrodynamic interactions between monomers. The hydrodynamic interactions close to a surface are difficult to calculate explicitly. However if we remain at a scaling level, it is sufficient to introduce the local stretching blobs . If the local tension of the chain is $`\tau _{ch}`$, the local blob has a size $`\xi _tk_BT/\tau _{ch}`$ and contains $`g_t(\xi _t/a)^2`$ monomers. The friction force on a blob is $`6\pi \eta \xi _t`$ and the hydrodynamic interactions are screened at the scale of the blob so that the dynamics of the chain of blobs can by studied using the Rouse model; the total friction is the sum of the frictions over all the blobs. As in the Rouse model, a free chain pulled by one end with a velocity $`u`$ has on average a trumpet like shape studied in details by Brochard and her coworkers . The Zimm equation of motion is obtained by writing a force balance on the blobs
$$k\frac{^2z}{s^2}=6\pi \eta u\frac{z}{s}$$
(C1)
The average position of monomer $`s`$ is then
$$z_{free}(s)=\left(\frac{k_BT}{6\pi \eta u}\right)^{1/2}\left[\mathrm{exp}(\frac{6\pi \eta us}{k})1\right]$$
(C2)
The force balance on a section of chain containing $`s`$ monomers close to the free end point is obtained by integration of the equation of motion $`k\frac{z}{s}=6\pi \eta uz`$. Note that the tension increases very rapidly from the free end point and that the nonlinear elasticity certainly becomes relevant. We do not discuss this effect here and refer to reference .
For a chain trapped by one end on the surface, the sticking time does not depend on the dynamical model at low velocities (if it is larger than the characteristic relaxation time of the chain which is the Rouse time of the chain of blobs) $`t_s=N\tau /(uk)`$. The dissipated energy is still the elastic energy of a chain with uniform tension $`\tau `$.
At high velocities, the sticking time is obtained by a scaling argument similar to that used in the Rouse model. When the end point escapes from the trap after a time $`t_s`$, a chain section of $`s(t_s)`$ monomers is stretched with a tension $`\tau `$. The characteristic relaxation time of this chain section is the Rouse time of the corresponding chain of blobs $`t_s=\frac{6\pi \eta s^2\tau }{k^2}`$. Using the integrated equation of motion, we find the sticking time $`t_s=\frac{\tau }{6\pi \eta u^2}`$. It has the same dependence on the velocity as in the Rouse model but a different dependence on the tension which is due to the fact that one must now consider the chain as a Rouse chain of Pincus blobs and not of monomers.
The crossover velocity is obtained by comparing the low and high velocity results $`u_c=\frac{k}{6\pi \eta N}`$ and the energy dissipated on a trap at high velocity is the elastic energy of a chain section containing $`s(t_s)`$ monomers: $`W_d=\frac{\tau ^2}{6\pi \eta u}`$. Here also we find the same dependence on velocity but a different tension dependence if we compare to the Rouse model.
## D Trapping cross section
In this appendix, we study the trapping cross section for the free chain end by the potential well of a trap on the surface. We focus on the undisturbed part of chain end of size $`\delta r`$. At very low velocity $`uu_0=R/T_R`$ ($`T_R`$ is the Rouse time) and $`\delta r`$ is the chain radius $`R`$. At higher velocity the chain is stretched by the hydrodynamic friction and $`\delta r`$ is equal to the local stretching blob size ($`\delta ru^{1/3}`$). If we take the origin at the center of the first blob containing the free chain end, the surface monomer concentration at a distance $`r`$ is :
$$c(r)\delta r^2\mathrm{exp}(r/\delta r)^2$$
(D1)
The trap acts as an absorbing site and we consider that the chain end gets trapped when it diffuses to the surface of the trap of radius $`b`$ (figure 1). If the trap is at a distance $`r`$ from the average position of the end point, the concentration gradient close to the trap is $`c(r)/b`$. The total flux of end point towards the trap is thus given by :
$$q=2\pi bDc(r)/b$$
(D2)
The relevant diffusion constant D is equal to the Rouse diffusion constant of the whole chain $`D=k_BT/(N\zeta )`$ at very low velocity and to the Rouse diffusion coefficient of the stretching blob $`D=k_BT/(s\zeta )`$ ($`s`$ is the number of monomers in the blob) at high velocity. Note that in two dimensions there are logarithmic corrections to this result as there is no stationary diffusion state. If the impact parameter (the minimum distance between the trap and the average chain end position) is $`c`$ and if the distance is minimum at time $`t=0`$ the distance $`r`$ varies with time as $`r^2=c^2+u^2t^2`$. The probability $`p(c)`$ that the chain with an impact parameter $`c`$ gets trapped is obtained by integration over time of the flux $`q`$, $`p(c)\frac{u_0}{u}\mathrm{exp}(c/R)^2`$. This probability is of course smaller than one. The trapping probability is thus equal to one if $`c<R\mathrm{log}^{\frac{1}{2}}(u_0/u)`$ and decays as a gaussian for larger values of $`c`$. The trapping cross section (a length in two dimensions) is obtained by integration over the impact parameter
$$\sigma R\mathrm{log}^{\frac{1}{2}}(u_0/u)$$
(D3)
At higher velocity $`u>u_0`$, the trapping probability is always smaller than $`1`$ and one can integrate the flux $`q`$ to find the trapping probability and then the cross-section. One must however take into account that $`\delta r`$ is now the size of the stretching blob and that the diffusion constant is not the diffusion constant of the whole chain but the one of a chain of size $`\delta r`$. The result is
$$\sigma \left(\frac{a^2k_BT}{\zeta u}\right)^{1/3}$$
(D4)
The average distance between trapping events can then be calculated from the surface density of traps $`n`$, $`d=1/\sigma n`$.
Figures caption
* Figure 1: Polymer chain pulled at constant velocity $`u`$ on a periodic 2D lattice of traps.
* Figure 2: Force (normalised by N$`\zeta u`$) applied on the trapped monomer (s=0, black) and on the pulled one (s=N, grey) versus time (normalised by T<sub>R</sub>), according to equation 3.
* Figure 3: Friction force in a ($`u`$, $`d`$) diagram for the many traps problem at zero temperature.
* Figure 4: Trapping potential a) $`V(z)`$ and b) $`U(z)=V(z)\tau _{tr}z`$.
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# Green’s Function Monte Carlo study of SU(3) lattice gauge theory in (3+1)D
## I Introduction
Classical Monte Carlo simulations provide a very powerful and accurate method for the study of Euclidean lattice gauge theories. In the Hamiltonian formulation , on the other hand, the corresponding quantum Monte Carlo methods have been somewhat neglected. Here we present a study of SU(3) Yang-Mills theory in (3+1) dimensions, using the Green’s Function Monte Carlo approach , adapted to lattice gauge theory by Chin et al .
Quantum Monte Carlo methods in Hamiltonian lattice gauge theory have a somewhat chequered history. The first calculations used a strong-coupling basis involving discrete “electric field” link variables, and a “Projector Monte Carlo” approach , which used the Hamiltonian itself to project out the ground state. A later version of this was the “stochastic truncation” approach of Allton et al . Using this approach one can successfully compute string tensions and mass gaps for Abelian models . For non-Abelian models, however, some technical problems arose . The use of a Robson-Webber recoupling scheme at the lattice vertices requires the use of Clebsch-Gordan coefficients or 6$`j`$-symbols, which are not known to high order for $`SU(3)`$; and furthermore, the ‘minus sign’ problem rears its head, in that destructive interference occurs between different paths to the same final state. It may well be that a better choice of strong-coupling basis, such as the ‘loop representation’, might avoid these problems; but this has not yet been demonstrated.
In the meantime, Heys and Stump and Chin et al. pioneered the use of “Green’s Function Monte Carlo” (GFMC) or “Diffusion Monte Carlo” techniques in Hamiltonian LGT, in conjunction with a weak-coupling representation involving continuous gauge field link variables. This was successfully adapted to non-Abelian Yang-Mills theories , with no minus sign problem arising. In this representation, however, one is simulating the wave function in gauge field configuration space by a discrete ensemble or density of random walkers: it is not possible to determine the derivatives of the gauge fields for each configuration, or to enforce Gauss’s law explicitly, and the ensemble always relaxes back to the ground state sector. Hence one cannot compute the string tensions and mass gaps directly as Hamiltonian eigenvalues corresponding to ground states in different sectors, as one does in the strong-coupling representation. Instead, one is forced back to the more laborious approach used in Euclidean calculations: namely, to measure an appropriate correlation function, and estimate the mass gap as the inverse of the correlation length. We have introduced the ‘forward-walking’ technique, well-known in many-body theory , to measure the expectation values and correlation functions. The technique has been demonstrated for the cases of the transverse Ising model in (1+1)D , and the U(1) LGT in (2+1)D .
Here we apply the technique for the first time to a non-Abelian model, namely SU(3) Yang-Mills theory in (3+1)D. The ground state energy and Wilson loop values are calculated, and approximate values are extracted for the string tension in the weak-coupling regime. Comparisons are made with earlier calculations, where they are available .
Our conclusions are that the method is a viable one, but requires the use of an improved “guiding wave function” to achieve better accuracy. There are certain drawbacks intrinsic to the method , such as the necessity to use a branching algorithm and a guiding wave function, which tend to introduce substantial errors into the results, both statistical and systematic. For these reasons, it may be preferable to employ a Path Integral Monte Carlo approach to these models, which avoids the problems above.
Our methods are presented in Section II, the results are outlined in Section III, and our conclusions are discussed in Section IV.
## II METHOD
### A Lattice Hamiltonian
The Green’s Function Monte Carlo formalism has been adapted to SU(2) Yang-Mills theory by Chin, van Roosmalen, Umland and Koonin , and sketched for the SU(3) case by Chin, Long and Robson . Here we provide a slightly fuller discussion of the SU(3) case, following the earlier treatment of Chin, van Roosmalen et al .
The SU(3) lattice Hamiltonian is given by
$$H=\frac{g^2}{2a}\{\underset{l}{}E_l^aE_l^a\frac{\lambda }{3}\underset{p}{}Tr(U_p+U_p^{})\}$$
(1)
where $`E_l^a`$ is a component of the electric field at link l, $`\lambda =6/g^4`$, the index $`a`$ runs over the 8 generators of SU(3), and $`U_p`$ denotes the product of four link operators around an elementary plaquette.
Commutation relations between electric field and link operators are
$$[E_l^a,U_l^{}]=\frac{1}{2}\lambda ^aU_l\delta _{ll^{}},$$
(2)
choosing the $`E_l^a`$ as left generators of SU(3), where the $`\{\lambda ^a,a=1,8\}`$ are the Gell-Mann matrices for SU(3). We will work with the dimensionless operator
$$H=\frac{1}{2}\underset{l}{}E_l^aE_l^a\frac{\lambda }{6}\underset{p}{}Tr(U_p+U_p^{})$$
(3)
The link variables are elements of the group SU(3) in the fundamental representation
$$U=\mathrm{exp}(i\frac{1}{2}\lambda ^aA^a)$$
(4)
There is no simple equivalent of the quaternion representation for SU(2). Following Beg and Ruegg , we can represent
$`U`$ $`=`$ $`\left(\begin{array}{ccc}z_1& z_2& z_3\\ u_1& u_2& u_3\\ w_1& w_2& w_3\end{array}\right)`$ (8)
where z,u,w are three-dimensional complex vectors; then if U is to be unitary, we require z,u,w to be orthonormal; and if U is to have determinant unity, we require
$$ϵ_{ijk}z_iu_jw_k=1$$
(9)
which is satisfied if
$$w_i=ϵ_{ijk}z_i^{}u_k^{}$$
(10)
One possible representation which satisfies these conditions is
$`z_1`$ $`=`$ $`(1i(x_3+x_8/\sqrt{3}))/N_1,`$ (11)
$`z_2`$ $`=`$ $`(x_2ix_1)/N_1,`$ (12)
$`z_3`$ $`=`$ $`(x_5ix_4)/N_1`$ (13)
where
$$N_1=[1+(x_3+x_8/\sqrt{3})^2+x_1^2+x_2^2+x_4^2+x_5^2]^{1/2};$$
(14)
and
$$u_3=(x_7ix_6)/N_2\stackrel{~}{u_3}/N_2,$$
(15)
$$u_2=(1i(x_8/\sqrt{3}x_3))/N_1\stackrel{~}{u_2}/N_2,$$
(16)
$$u_1=\stackrel{~}{u_1}/n_2,\stackrel{~}{u_1}=[\stackrel{~}{u_2}z_2^{}+\stackrel{~}{u_3}z_3^{}]/z_1^{}$$
(17)
where
$$N_2=[|\stackrel{~}{u_1}|^2+|\stackrel{~}{u_2}|^2+|\stackrel{~}{u_3}|^2]^{1/2}$$
(18)
and
$$w_i=ϵ_{ijk}z_j^{}u_k^{}.$$
(19)
This involves 8 unrestricted parameters $`\{x_a,a=\mathrm{1..8}\}`$, as expected. For small x,
$$UIi\lambda ^ax^a$$
(20)
i.e.
$$x^aA^a/2,$$
(21)
where $`A^a`$ is the gauge field on that link. The product of two link variables can be found by simple matrix multiplication.
### B Green’s Function Monte Carlo method
The Green’s Function Monte Carlo method employs the operator $`\mathrm{exp}(\tau (HE))`$, i.e. the time evolution operator in imaginary time, as a projector onto the ground state $`|\psi _0`$:
$`|\psi _0`$ $``$ $`\underset{\tau \mathrm{}}{lim}\{e^{\tau (HE)}|\mathrm{\Phi }\}`$ (22)
$`=`$ $`\underset{\mathrm{\Delta }\tau 0,N\mathrm{\Delta }\tau \mathrm{}}{lim}e^{N\mathrm{\Delta }\tau (HE)}|\mathrm{\Phi }`$ (23)
where $`|\mathrm{\Phi }`$ is any suitable trial state. To procure some variational guidance, one performs a “similarity transformation” with the trial wave function $`\mathrm{\Phi }`$, and evolves the product $`\mathrm{\Phi }|\psi _0`$ in imaginary time. The heart of the procedure is the calculation of the matrix element corrersponding to a single small time step $`\mathrm{\Delta }\tau `$. Chin et al show that
$`𝐱^{}|\mathrm{\Phi }e^{\mathrm{\Delta }\tau (HE)}\mathrm{\Phi }^1|𝐱`$ $`=`$ $`{\displaystyle \underset{l}{}}U_l^{}|N\{\mathrm{exp}({\displaystyle \frac{1}{2}}\mathrm{\Delta }\tau E_l^aE_l^a)\mathrm{exp}[\mathrm{\Delta }\tau E_l^a(E_l^a\mathrm{ln}\mathrm{\Phi })]\}|U_l`$ (25)
$`\mathrm{exp}\{\mathrm{\Delta }\tau [E\mathrm{\Phi }^1H\mathrm{\Phi }(𝐱)]\}+O(\mathrm{\Delta }\tau ^2)`$
$``$ $`p(𝐱^{},𝐱)w(𝐱)+O(\mathrm{\Delta }\tau ^2)`$ (26)
where $`𝐱=\{U_l\}`$ denotes an entire lattice configuration of link fields.
The product $`\mathrm{\Phi }|\psi `$ is simulated by the density of an ensemble of random walkers, as in the SU(2) case. At the kth. step, the ‘weight’ of each walker at $`𝐱_k`$ is multiplied by $`w(𝐱_k)`$ and the next ensemble $`\{𝐱_{k+1}\}`$ is evolved from $`\{𝐱_k\}`$ according to the matrix element $`p(𝐱_{𝐤+\mathrm{𝟏}},𝐱_𝐤)`$. The effect of $`p(𝐱_{𝐤+\mathrm{𝟏}},𝐱_𝐤)`$ is to alter each link variable $`U_l`$ in $`\{𝐱_𝐤\}`$ to $`U_l^{}`$ by a Gaussian random walk plus a “drift step” guided by the trial wave function:
$$U^{}=\mathrm{\Delta }UU_dU$$
(27)
where $`U_d=\mathrm{exp}[i\frac{1}{2}\lambda ^a(i\mathrm{\Delta }\tau E^a\mathrm{ln}\mathrm{\Phi })]`$ is the drift step, and $`\mathrm{\Delta }U`$ is an SU(3) group element randomly chosen from a Gaussian distribution around the identity, with variance $`\mathrm{\Delta }s^2=8\mathrm{\Delta }\tau `$ (i.e. $`\mathrm{\Delta }\tau `$ for each index a), where
$$\mathrm{\Delta }s^2\underset{a}{}A^aA^a=8\mathrm{\Delta }\tau ,\text{(small }A^a)$$
(28)
or
$$x^ax^a=\frac{1}{4}A^aA^a\frac{\mathrm{\Delta }\tau }{4},\text{each}a.$$
(29)
The simulation is carried out for a large number of iterations $`\mathrm{\Delta }\tau `$, until an equilibrium distribution $`\mathrm{\Phi }|\psi _0`$ is reached. The energy E in (26) is adjusted after each iteration so as to maintain the total ensemble weight constant. The average value of E can then be taken as an estimate of $`E_0`$, the ground-state energy.
As time evolves, the weights of some walkers grow larger, while others grow smaller, which would produce an increased statistical error. To avoid this, a “branching” process is employed, whereby a walker with weight larger than some threshold is split into two independent walkers, while others with weights lower than another threshold are amalgamated. We use Runge’s technique for this purpose.
### C Trial Wave Function
The trial wave function is chosen to be the one-parameter form
$$\mathrm{\Phi }=\mathrm{exp}[\alpha \underset{p}{}Tr(U_p+U_p^{})]$$
(30)
Then the drift step is
$`U_d`$ $`=`$ $`\mathrm{exp}[i{\displaystyle \frac{\lambda ^a}{2}}(i\mathrm{\Delta }\tau E^a\mathrm{ln}\mathrm{\Phi })]`$ (31)
$``$ $`\mathrm{exp}[i{\displaystyle \frac{\lambda ^a}{2}}A_l^a]`$ (32)
for each link, where
$$A_l^a=i\mathrm{\Delta }\tau \frac{\alpha }{2}\underset{pl}{}Tr[\lambda ^aU_l..U_4^{}h.c.]$$
(33)
i.e.
$$x_l^a\frac{A_l^a}{2}=i\frac{\alpha \mathrm{\Delta }\tau }{4}\underset{pl}{}Tr[\lambda ^aU_l..U_4^{}h.c.]$$
(34)
(note: the effect of $`E_l^a`$ on a plaquette operator is to ‘insert’ a $`\lambda ^a/2`$ in front of the appropriate link operator $`U_l`$, to be followed by the remaining link operators in the plaquette, taken in the direction of the link l).
Finally, the trial energy factor is
$`\mathrm{\Phi }^1H\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \underset{l}{}}\{{\displaystyle \frac{\alpha ^2}{8}}({\displaystyle \underset{pl}{}}Tr[\lambda ^aU_l..U_4^{}h.c.])^2+({\displaystyle \frac{2\alpha }{3}}{\displaystyle \frac{\lambda }{24}}){\displaystyle \underset{pl}{}}Tr(U_p+U_p^{})\}.`$ (35)
Therefore the weight factor is
$`w(𝐱)`$ $`=`$ $`\mathrm{exp}\{\mathrm{\Delta }\tau (E\mathrm{\Phi }^1H\mathrm{\Phi })\}`$ (36)
$`=`$ $`\mathrm{exp}\{\mathrm{\Delta }\tau (E_{trial}({\displaystyle \frac{2\alpha }{3}}{\displaystyle \frac{\lambda }{24}}){\displaystyle \underset{l}{}}{\displaystyle \underset{pl}{}}2Re\{TrU_p\}`$ (38)
$`+{\displaystyle \frac{\alpha ^2}{8}}{\displaystyle \underset{l}{}}({\displaystyle \underset{pl}{}}2Im\{Tr[\lambda ^aU_l..U_4^{}]\})^2))\}.`$
### D Forward Walking estimates
The “forward walking” technique is used to estimate expectation values . Its application to the U(1) lattice gauge theory in (2+1)D was discussed by Hamer et al . It is based on the following equation: for an operator Q,
$`<Q>_0`$ $`=`$ $`{\displaystyle \frac{\psi _0|Q|\psi _0}{\psi _0|\psi _0}}`$ (39)
$`\stackrel{}{J\mathrm{}}`$ $`{\displaystyle \frac{\mathrm{\Phi }|K^JQ|\psi _0}{\mathrm{\Phi }|K^J|\psi _0}}`$ (40)
$`=`$ $`{\displaystyle \frac{\stackrel{~}{K}(𝐱_J,𝐱_{J1})\mathrm{}\stackrel{~}{K}(𝐱_2,𝐱_1)Q(𝐱_1)\stackrel{~}{\psi }_0(𝐱_1)}{\stackrel{~}{K}(𝐱_J,𝐱_{J1})\mathrm{}\stackrel{~}{K}(𝐱_2,𝐱_1)\stackrel{~}{\psi }_0(𝐱_1)}}`$ (41)
where $`K(𝐱_J,𝐱_{J1})`$ is the evolution operator for time $`\mathrm{\Delta }\tau `$, and $`\stackrel{~}{K}(𝐱_J,𝐱_{J1})`$ is the same operator in the similarity transformed basis. Again we have assumed that the operator $`Q`$ is diagonal in the basis of plaquette variables $`𝐱`$.
This equation is implemented by the following procedure:
* Record the value $`Q(𝐱_i)`$ for each “ancestor” walker at the beginning of a measurement;
* Propagate the ensemble as normal for $`J`$ iterations, keeping a record of the “ancestor” of each walker in the current population;
* Take the weighted average of the $`Q(𝐱_i)`$ with respect to the weights of the descendants of $`𝐱_i`$ after the $`J`$ iterations, using sufficient iterations $`J`$ that the estimate reaches a ‘plateau’.
## III Results
Simulations were carried out for LxLxL lattices up to L=8 sites, using runs of typically 4000 iterations and an ensemble size of 250 to 1000 depending (inversely) on lattice size. These statistics are approximately 100 times less than those used in the U(1) calculations , but about 100 times greater than those used in the previous SU(3) calculations of Chin et al . Time steps $`\mathrm{\Delta }\tau `$ of 0.01 and 0.05 “seconds” were used, with each iteration consisting of 5 sweeps and 1 sweep through the lattice, respectively,followed by a branching process. The first 400 iterations were discarded to allow for equilibration. The data were block averaged over blocks of up to 256 iterations, to minimize the effect of correlations on the error estimates.
The results taken at $`\mathrm{\Delta }\tau =0.01`$ and $`\mathrm{\Delta }\tau =0.05`$ were extrapolated linearly to $`\mathrm{\Delta }\tau =0`$. Figure 1 demonstrates that the dependence of the ground-state energy on $`\mathrm{\Delta }\tau `$ is approximately linear.
The variational parameter $`c`$ was given values as shown in Table I. These are essentially the values used by Chin et al , obtained from a variational Monte Carlo calculation. We checked that these were approximately the optimum values for small lattices.
Forward-walking measurements were taken over $`J`$ iterations, where $`J`$ ranged from 20 to 100, depending on the coupling $`\lambda `$. Ten separate measurements were taken over this time interval, in order to check whether the value measured by forward-walking had reached equilibrium. A new measurement was started soon after the previous one had finished.
### A Ground-state Energy
The dependence of the ground-state energy on the variational parameter $`c`$ is illustrated in Figure 2. It can be seen that the energy reaches a broad minimum at about the expected value ($`c=0.33`$ at this coupling).
Our estimates of the ground-state energy are listed in Table I, as a function of the coupling $`\lambda `$ and lattice size $`L`$. The dependence on lattice size is illustrated in Figure 3, at two fixed couplings $`\lambda =3.0`$ and $`\lambda =5.0`$. In the “strong-coupling” case, $`\lambda =3.0`$, it can be seen that the results converge exponentially fast in $`L`$, whereas in the “weak-coupling” regime, $`\lambda =5.0`$, the convergence is more like $`1/L^4`$ at these lattice sizes. This behaviour merits some further explanation.
A similar phenomenon occurs in the case of the U(1) theory in (2+1)D . In the strong-coupling regime, where the mass gap is large, the usual exponential convergence occurs. In the weak-coupling regime, however, where the mass gap M is very small, the finite-size scaling behaviour for small lattice sizes is that of a massless theory, and it is only at much larger lattice sizes $`L1/M`$ that a crossover to exponential convergence occurs. In the U(1) case, it has been shown that the finite-size scaling behaviour at small $`L`$ is well described by an “effective Lagrangian” approach, using the Lagrangian for free, massless photons that the model was originally constructed to simulate. In the same way, an “effective Lagrangian” corresponding to free, massless gluons (non-interacting QCD) should describe the finite-size behaviour in the present case, in line with the idea of asymptotic freedom. By analogy with the (2+1)D case, we expect a $`1/L^4`$ dependence for the corrections to the ground-state energy per site. We hope to pursue this analysis further at a later date.
An anomalous feature in Figure 3b) is that the $`L=8`$ point lies well out of line with the others. This occurs at other couplings also. We suspect that the results for $`L=8`$ are not reliable, and that the trial wave function will have to be further improved to give reliable results for such large lattices. Supporting evidence for this will be presented below.
We have made estimates of the bulk limit, extrapolating mainly from the smaller L values where possible, and the results are listed in Table II. Our present estimates generally lie a little below those of Chin et al , and we believe them to be more accurate in view of our greater statistics. The estimates for the bulk ground-state energy per site are graphed as a function of coupling in Figure 4, where they are compared with previous estimates obtained by an ‘Exact Linked Cluster Expansion’ (ELCE) procedure, and with the asymptotic weak-coupling series
$$ϵ_03\lambda +7.798\lambda ^{1/2},\lambda \mathrm{}.$$
(42)
The Monte Carlo results agree very well with the ELCE estimates, and appear to match nicely onto the expected weak-coupling behaviour for $`\lambda 6`$.
### B Wilson Loops
The forward-walking method was used to estimate values for the $`mxn`$ Wilson loops, $`W(m,n)`$. Figure 5 shows an example, namely the estimate of the mean plaquette $`W(1,1)`$ as a function of $`J`$ for the case $`L=6`$, $`\lambda =1.5`$. It can be seen that the estimate relaxes exponentially towards a plateau value as the number of iterations $`J`$ increases: an exponential fit is performed to estimate the asymptotic value. It can also be seen that the statistical error on each point is much larger than the point-to-point variation: the estimates are highly correlated, and all the points tend to move up and down together as one goes from one sample to the next.
A problem that arises in these measurements is the loss in statistical accuracy at large couplings. At large couplings the weights of the random walkers vary rapidly with time, and it can easily happen that during a measurement the descendants of all the ‘ancestor’ walkers but one die out from the ensemble, at which point the result ‘freezes’, and the number of members of the ensemble has effectively been reduced to one. This inevitably means a severe loss of statistical accuracy.
A graph of the ‘mean plaquette’ $`W(1,1)`$ versus the variational parameter $`c`$ is shown in Figure 6. Another problem is immediately apparent. The estimate for $`W(1,1)`$ is not independent of $`c`$, in fact it depends linearly on $`c`$ over this range, and the size of the variation is such that the probable systematic error due to the choice of $`c`$ is an order of magnitude larger than the random statistical error in the results. Thus it would be advatageous in future studies to put more effort into improving the trial wave function, rather than merely improving the statistics.
Figure 7 shows examples of the dependence of the results on lattice size $`L`$. Once again, the results at strong coupling $`\lambda =1.5`$ converge exponentially fast, while those at the weak coupling value $`\lambda =5.0`$ can be approximately fitted by a $`1/L^4`$ dependence. The $`L=8`$ value and even the $`L=6`$ value again lie off the trend of the smaller lattices, and are probably not very reliable.
Estimates of the bulk limit are listed in Table II. The estimates for the mean plaquette are graphed as a function of coupling $`\lambda `$ in Figure 8, and compared with series estimates at strong and weak coupling . The agreement is quite good.
### C String Tension
Having obtained estimates for the Wilson loop values on the bulk lattice, one can extract estimates for the ‘spacelike’ string tension using the Creutz ratios:
$$Ka^2R_n=\mathrm{ln}\left[\frac{W(n,n)W(n1,n1)}{W(n,n1)^2}\right]$$
(43)
or the cruder 2-point estimates
$$R_n^{}=\frac{1}{n}\mathrm{ln}\left[\frac{W(n,n)}{W(n,n1)}\right]$$
(44)
The results are shown in Figure 9. Also shown in Figure 9 are some previous estimates derived from the ‘axial’ string tension, obtained using an ‘Exact Linked Cluster Expansion’ (ELCE) method. The axial string tension $`aT`$ is calculated as an energy per link, and must be converted to a dimensionless, ‘spacelike’ tension by dividing by the ‘speed of light’ c,
$$Ka^2=\frac{aT}{c},$$
(45)
where
$$c\frac{2}{g^2}[10.1671g^2]=\sqrt{\frac{2\lambda }{3}}\left[10.1671\sqrt{\frac{6}{\lambda }}\right],\lambda \mathrm{}.$$
(46)
We have also used the weak-coupling relationship between the scales of Euclidean and Hamiltonian lattice Yang-Mills theory calculated by Hasenfratz et al to plot the results against the Euclidean coupling $`\beta =6/g_E^2`$, where
$$\beta =\sqrt{6\lambda }0.01308$$
(47)
It can be seen that the present GFMC results are in rough agreement with the axial string tension results in the region $`4\beta 5`$, which is also the region where the ‘roughening’ transition occurs in the string tension . For $`\beta >5`$, however, the Creutz ratio $`R_2`$ runs above the ELCE estimate, and shows no sign of the expected crossover to an exponentially decreasing scaling behaviour at $`\beta 6`$. We presume that this is a finite-size effect, and that the Creutz ratios $`R_n`$ for larger $`n`$ will show a substantial decrease in the ‘weak-coupling’ regime $`\beta 6`$. That is certainly the pattern seen in the Euclidean calculations , or in the U$`(1)_{2+1}`$ model . Unfortunately, however, our present results for the larger Wilson loops are not of sufficient accuracy to allow worthwhile estimates of $`R_n`$ for $`n2`$.
## IV Summary and Conclusions
We have presented the results of a new Green’s Function Monte Carlo study of the SU(3) Yang-Mills theory in the (3+1)D Hamiltonian formulation. A forward-walking method has been used to estimate values for the Wilson loops as well as the ground-state energy, and hence some rather crude estimates of the string tension have been extracted. Comparisons have been made with an earlier Hamiltonian calculation of the axial string tension . The two sets of results agree in the ‘roughening’ region; but our Monte Carlo results do not extend to the large Wilson loops that would be required to demonstrate ‘scaling’ behaviour in the weak-coupling regime.
Some significant problems with the GFMC method have emerged from this study. The ‘forward-walking’ technique was introduced specifically to avoid any variational bias from the trial wave function . As it turns out, however, the results for the Wilson loops show a substantial dependence on the trial wave function parameter $`c`$. The systematic error due to this dependence is an order of magnitude larger than the statistical error, so it would pay to put more effort in future studies into improving the trial wave function, rather than simply increasing the statistics. Furthermore, the effective ensemble size decreases during each measurement as the descendants of each ‘ancestor’ state die out, and this produces a substantial loss in statistical accuracy at weak coupling, as well.
It would be preferable if one were able to do away entirely with all the paraphernalia of trial wave function, weights, branching algorithms, etc, and just rely on some sort of Metropolis-style accept/reject algorithm to produce a correct distribution of walkers. Within a quantum Hamiltonian framework, a way is known to do this, namely the Path Integral Monte Carlo (PIMC) approach . We conclude that the PIMC approach may be better suited than GFMC to the study of large and complicated lattice Hamiltonian systems.
## Acknowledgments
This work is supported by the Australian Research Council. Calculations were performed on the SGI Power Challenge Facility at the New South Wales Centre for Parallel Computing and the Fujitsu VPP300 vector machine at the Australian National Universtiy Supercomputing Facility: we are grateful for the use of these facilities.
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# Higher Type Adjunction Inequalities in Seiberg-Witten Theory
## 1. Introduction
In this paper, we prove certain adjunction inequalities, which give relations between the Seiberg-Witten invariants of a four-manifold $`X`$ and the genus of embedded surfaces in $`X`$. These results are generalizations of results from , , , see also .
The investigations center on a construction of an appropriate Seiberg-Witten-Floer functor for manifolds which bound circle bundles $`Y`$ over Riemann surfaces (with sufficiently large Euler number), which relies on the calculations of . Special cases of this theory were studied in , where the authors used similar techniques to prove the symplectic Thom conjecture. That problem requires an analysis of those $`\mathrm{Spin}_{}`$ structures over $`Y`$ for which the Seiberg-Witten moduli space contains only reducible solutions, which simplifies the corresponding Floer homology. In this paper, we work out the theory in the other, more complicated cases. We will give more applications of these techniques in .
Before stating the results, we set up some notation. Let $`X`$ be a closed, connected, smooth four-manifold equipped with an orientation for which $`b_2^+(X)>0`$ (where $`b_2^+(X)`$ is the dimension of a maximal positive-definite linear subspace $`H^+(X;)`$ of the intersection pairing on $`H^2(X;)`$) and an orientation for $`H^1(X;)H^+(X;)`$. Given such a four-manifold, together with a $`\mathrm{Spin}_{}`$ structure $`𝔰`$, the Seiberg-Witten invariants (see , , ) form an integer-valued function
$$SW_{X,𝔰}:𝔸(X),$$
where $`𝔸(X)`$ denotes the graded algebra obtained by tensoring the exterior algebra on $`H_1(X)`$ (graded so that $`H_1(X)`$ has grading one) with the polynomial algebra $`[U]`$ on a single two-dimensional generator. The invariants are constructed via intersection theory on the moduli space $`_X(𝔰)`$ of solutions $`(A,\mathrm{\Phi })`$ modulo gauge to the Seiberg-Witten equations in $`𝔰`$:
(1) $`\rho (F_A^+)`$ $`=`$ $`i\{\mathrm{\Phi },\mathrm{\Phi }\}_0\rho (i\eta )`$
(2) $`\overline{)}D_A\mathrm{\Phi }`$ $`=`$ $`0,`$
where $`\mathrm{\Phi }`$ is a section of $`W^+`$, $`A`$ is a spin-connection in the spinor bundle $`W^+`$ of $`𝔰`$, $`\overline{)}D_A`$ denotes the associated Dirac operator, $`\rho `$ denotes Clifford multiplication, $`\eta `$ is some fixed self-dual two-form, and $`\{\mathrm{\Phi },\mathrm{\Phi }\}_0`$ is the usual quadratic map (see ). Note that the invariants are zero on homogeneous elements whose degree is not $`d(𝔰)`$, where
$$d(𝔰)=\frac{c_1(𝔰)^2(2\chi (X)+3\sigma (X))}{4}$$
denotes the formal dimension of the moduli space $`_X(𝔰)`$. When $`b_2^+(X)>1`$, $`SW_{X,𝔰}`$ is a diffeomorphism invariant of the four-manifold; when $`b_2^+(X)=1`$, the invariants depend on a chamber structure (see , ). There are two distinguished chambers corresponding to the two components of $`𝒦(X)=\{\omega H^2(X;)0|\omega ^20\}`$. Given a component $`𝒦_0`$ of $`𝒦(X)`$, the corresponding invariant (still denoted $`SW_{X,𝔰}`$) is calculated using the moduli space of solutions to the Seiberg-Witten equations perturbed by any generic self-dual two-form $`\eta `$, provided that the sign of $`2\pi c_1(𝔰)\omega _g+_X\eta \omega _g`$ agrees with the sign of $`\gamma \omega _g`$, where $`\gamma `$ is any class in $`𝒦_0`$, and $`\omega _g0`$ is a harmonic (with respect to the metric $`g`$), self-dual two-form over $`X`$. Note that $`SW_{X,𝔰}`$ is a diffeomorphism invariant of $`X`$ (and the component $`𝒦_0`$).
Those $`\mathrm{Spin}_{}`$ structures $`𝔰`$ for which the invariant $`SW_{X,𝔰}`$ is non-trivial are called basic classes.
Our results are easiest to state when $`b_1(X)=0`$, where we have the following.
###### Theorem 1.1.
Let $`X`$ be a smooth, closed, connected, oriented four-manifold with $`b_2^+(X)>0`$ and $`b_1(X)=0`$, and let $`\mathrm{\Sigma }X`$ be a smoothly-embedded surface with genus $`g(\mathrm{\Sigma })>0`$ representing a non-torsion homology class with self-intersection number $`[\mathrm{\Sigma }][\mathrm{\Sigma }]0`$. If $`b_2^+(X)>1`$, then we have the following adjunction inequality
$$|c_1(𝔰),[\mathrm{\Sigma }]|+[\mathrm{\Sigma }][\mathrm{\Sigma }]+2d(𝔰)2g(\mathrm{\Sigma })2,$$
for each basic class $`𝔰\mathrm{Spin}_{}(X)`$. Furthermore, when $`b_2^+(X)=1`$, for each basic class $`𝔰`$ of $`X`$ for the component of $`𝒦(X)`$ which contains $`\mathrm{PD}[\mathrm{\Sigma }]`$ with
$$c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]0,$$
we have an inequality
$$c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]+2d(𝔰)2g(\mathrm{\Sigma })2.$$
###### Remark 1.2.
The above theorem should be seen as a refinement of the adjunction inequality proved by Kronheimer-Mrowka and Morgan-Szabó-Taubes (see , , ). Analogous results for immersed spheres were obtained by Fintushel and Stern, see .
In fact, Theorem 1.1 follows from a more general version. To state this, note first that an inclusion $`i:\mathrm{\Sigma }X`$ induces a map
$$i_{}:𝔸(\mathrm{\Sigma })𝔸(X).$$
###### Theorem 1.3.
Let $`X`$ be a smooth, closed, connected, oriented four-manifold with $`b_2^+(X)>0`$. Let $`\mathrm{\Sigma }X`$ be a surface with genus $`g(\mathrm{\Sigma })>0`$ representing a non-torsion homology class with self-intersection number $`[\mathrm{\Sigma }][\mathrm{\Sigma }]0`$. Let $`\mathrm{}`$ be an integer so that there is a symplectic basis $`\{A_j,B_j\}_{j=1}^g`$ for $`H_1(\mathrm{\Sigma })`$ so that $`i_{}(A_j)=0`$ in $`H_1(X;)`$ for $`i=1,\mathrm{},\mathrm{}`$. Let $`a𝔸(X)`$ and $`b𝔸(\mathrm{\Sigma })`$ be an element of degree $`d(b)\mathrm{}`$. If $`b_2^+(X)>1`$ then for each $`\mathrm{Spin}_{}`$ structure $`𝔰`$ so that $`SW_{X,𝔰}(ai_{}(b))`$ is non-zero, we have
$$|c_1(𝔰),[\mathrm{\Sigma }]|+[\mathrm{\Sigma }][\mathrm{\Sigma }]+2d(b)2g(\mathrm{\Sigma })2.$$
Furthermore, when $`b_2^+(X)=1`$ then for each $`\mathrm{Spin}_{}`$ structure $`𝔰`$ of $`X`$ with
$$c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]0,$$
for which $`SW_{X,𝔰}(ai_{}(b))`$ is non-zero, when calculated in the component of $`𝒦(X)`$ containing $`\mathrm{PD}[\mathrm{\Sigma }]`$, we have an inequality
(3)
$$c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]+2d(b)2g(\mathrm{\Sigma })2.$$
The Adjunction Inequality (3) does not hold without homological restrictions on $`X`$, as we can see by looking at the ruled surface $`X=S^2\times \mathrm{\Sigma }`$. In general, one can obtain only a weaker inequality (losing the factor of $`2`$ on the dimension $`d(b)`$), as follows.
###### Theorem 1.4.
Let $`X`$ be a smooth, closed, connected, oriented four-manifold with $`b_2^+(X)>1`$. Let $`\mathrm{\Sigma }X`$ be a surface with genus $`g(\mathrm{\Sigma })>0`$ representing a non-torsion homology class with self-intersection number $`[\mathrm{\Sigma }][\mathrm{\Sigma }]0`$. Let $`a𝔸(X)`$ and $`b𝔸(\mathrm{\Sigma })`$. If $`b_2^+(X)>1`$ and if $`SW_{X,𝔰}(ai_{}(b))`$ is non-zero for some $`b𝔸(\mathrm{\Sigma })`$ of degree $`d(b)`$, then we have
$$|c_1(𝔰),[\mathrm{\Sigma }]|+[\mathrm{\Sigma }][\mathrm{\Sigma }]+d(b)2g(\mathrm{\Sigma })2.$$
If $`b_2^+(X)=1`$ and $`𝔰`$ is a $`\mathrm{Spin}_{}`$ structure with
$$c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]0,$$
for which $`SW_{X,𝔰}(ai_{}(b))`$ is non-zero, when calculated in the component of $`𝒦(X)`$ containing $`\mathrm{PD}[\mathrm{\Sigma }]`$, then we have
(4)
$$c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]+d(b)2g(\mathrm{\Sigma })2.$$
###### Remark 1.5.
Adjunction inequalities for surfaces of positive square in Donaldson’s theory were first obtained in the influential paper of Kronheimer and Mrowka (see ). These inequalities were strengthened under similar, but more restrictive, hypotheses in their preprint ; see also . The conjectured relationship between the Donaldson and Seiberg-Witten invariants gives a correspondence between the adjunction inequalities arising in these two theories. For more on this correspondence, see , , , , , and .
Theorem 1.3 follows from a relation which holds for embedded surfaces with arbitrary self-intersection number. This relation can be viewed as a generalization of the relation appearing in . Once again, we begin by stating the case when $`b_1(X)=0`$, in the interest of exposition.
###### Theorem 1.6.
Let $`X`$ be a smooth, closed, connected, oriented four-manifold with $`b_1(X)=0`$, and let $`\mathrm{\Sigma }X`$ be a smoothly embedded surface with genus $`g(\mathrm{\Sigma })>0`$. Then, for each $`\mathrm{Spin}_{}`$ structure $`𝔰`$ with
$$c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]0$$
and
$$c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]+2d(𝔰)>2g(\mathrm{\Sigma })2,$$
we have
$$SW_{X,𝔰}(U^d)=SW_{X,𝔰\mathrm{PD}[\mathrm{\Sigma }]}(U^d^{}),$$
where $`d`$ and $`d^{}`$ denote the dimensions of $`𝔰`$ and $`𝔰\mathrm{PD}[\mathrm{\Sigma }]`$ respectively. In the case where $`b_2^+(X)=1`$, both invariants are to be calculated in the same component of $`𝒦(X)`$.
More generally, we have the following.
###### Theorem 1.7.
Let $`X`$ be a smooth, closed, connected, oriented four-manifold with $`b_2^+(X)>0`$. Let $`\mathrm{\Sigma }X`$ be a surface with genus $`g(\mathrm{\Sigma })>0`$. Let $`\mathrm{}`$ be an integer so that there is a symplectic basis $`\{A_j,B_j\}_{j=1}^g`$ for $`H_1(\mathrm{\Sigma })`$ so that $`i_{}(A_j)=0`$ in $`H_1(X;)`$ for $`i=1,\mathrm{},\mathrm{}`$. For each $`\mathrm{Spin}_{}`$ structure $`𝔰`$ with
$$c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]0$$
and each $`b𝔸(\mathrm{\Sigma })`$ of degree $`d(b)\mathrm{}`$ with
(5)
$$c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]+2d(b)>2g(\mathrm{\Sigma })2,$$
there is an element $`b^{}𝔸(\mathrm{\Sigma })`$ with $`d(b^{})d(b)`$ so that for any $`a𝔸(X)`$, we have
(6)
$$SW_{X,𝔰}(ai_{}(b))=SW_{X,𝔰\mathrm{PD}[\mathrm{\Sigma }]}(ai_{}(b^{})).$$
Furthermore, if $`b=U^{d/2}`$, then $`b^{}U^{d^{}/2}`$ lies in the ideal generated by $`H_1(\mathrm{\Sigma })`$ in $`𝔸(\mathrm{\Sigma })`$. Once again, in the case where $`b_2^+(X)=1`$, both invariants are to be calculated in the same component of $`𝒦(X)`$.
Theorem 1.3 is a simple consequence of Theorem 1.7, as the following proof shows.
Theorem 1.7 $``$ Theorem 1.3. Suppose Theorem 1.3 were false; i.e. suppose there were $`X`$, $`\mathrm{\Sigma }`$, $`𝔰`$, $`a`$, and $`b`$ which satisfy the hypotheses of the theorem, but which violate Adjunction Inequality (3). We can assume without loss of generality that
$$c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]0,$$
by reversing the orientation of $`\mathrm{\Sigma }`$ if necessary (when $`b_2^+(X)>1`$). Thus, Theorem 1.7 applies. Let $`b^{}`$ be the element which satisfies Relation (6), so we have that $`SW_{X,𝔰\mathrm{PD}[\mathrm{\Sigma }]}(ai_{}(b^{}))0`$. Since $`d(b)`$ and $`d(b^{})`$ are homogeneous elements with the same degree modulo two, and $`d(b^{})d(b)`$, it follows that we can find elements $`a^{}𝔸(X)`$ and $`b^{\prime \prime }𝔸(\mathrm{\Sigma })`$ with $`d(b^{\prime \prime })=d(b)`$, and $`SW_{X,𝔰\mathrm{PD}[\mathrm{\Sigma }]}(a^{}i_{}(b^{\prime \prime }))0`$. Now, since $`[\mathrm{\Sigma }][\mathrm{\Sigma }]0`$ and $`d(b^{\prime \prime })=d(b)`$, we see that $`\mathrm{\Sigma }`$ also violates the adjunction inequality for $`𝔰\mathrm{PD}[\mathrm{\Sigma }]`$, $`a^{}𝔸(X)`$, and $`b^{\prime \prime }𝔸(\mathrm{\Sigma })`$. Proceeding in this way, we see that $`𝔰n\mathrm{PD}[\mathrm{\Sigma }]`$ is a Seiberg-Witten basic class for all $`n0`$. If $`b_2^+(X)>1`$, then there are only finitely many basic classes of $`X`$, so since $`\mathrm{\Sigma }`$ is not a torsion class, we get a contradiction, proving Theorem 1.3 in this case.
The above argument works also when $`b_2^+(X)=1`$, since there are still only finitely many basic classes of the form $`𝔰n\mathrm{PD}[\mathrm{\Sigma }]`$ in the chamber corresponding to $`\mathrm{PD}[\mathrm{\Sigma }]`$. We see this as follows. Fix a metric $`g`$ on $`X`$ and a generic self-dual two-form $`\eta `$. Clearly, if $`𝔰`$ is fixed and $`n`$ is sufficiently large, the sign of $`\mathrm{PD}[\mathrm{\Sigma }]\omega _g`$ agrees with the sign of $`2\pi c_1(𝔰n\mathrm{PD}[\mathrm{\Sigma }])\omega _g+\eta \omega _g`$; i.e. for all large $`n`$, the $`\eta `$-perturbed moduli spaces for $`𝔰n\mathrm{PD}[\mathrm{\Sigma }]`$ can be used calculate the invariant in the component which contains $`\mathrm{PD}[\mathrm{\Sigma }]`$. But the usual compactness argument shows that all but finitely many of these moduli spaces are empty. Again, we have the contradiction completing the proof of Theorem 1.3. $`\mathrm{}`$
By blowing up, Theorem 1.7 is reduced to the case where the self-intersection number of $`\mathrm{\Sigma }`$ is sufficiently negative. The theorem is then proved by expressing the Seiberg-Witten invariants of a four-manifold with such an embedded surface $`\mathrm{\Sigma }`$ in terms of relative invariants, which take values in a Seiberg-Witten-Floer homology associated to non-trivial circle bundles over $`\mathrm{\Sigma }`$. In the presence of the topological hypotheses on the inclusion of $`H_1(\mathrm{\Sigma })`$ in $`H_1(X)`$, the above relation then follows from properties of this Floer homology.
The outline of this paper is as follows. In Section 2, we give examples which show that the adjunction inequalities are sharp. Our examples include four-manifolds with $`b_2^+(X)=1`$, and also examples where both $`b_2^+(X)>1`$ and $`b_1(X)>0`$. In Section 3, we show how Theorem 1.7 can be deduced from properties of a product formula, which relates the Seiberg-Witten invariants of a four-manifold containing an embedded surface with sufficiently negative self-intersection number with certain relative invariants associated to $`X\mathrm{\Sigma }`$. For completeness, we also show how a modified version of Theorem 1.7 implies Theorem 1.4. In Section 4 we review the gauge theory for circle bundles over Riemann surfaces as developed in . There is one $`\mathrm{Spin}_{}`$ structure in which the moduli space of reducibles has singularities (to which we return in a later section). In Section 5, we prove the product formula introduced in Section 3, assuming technical facts about the moduli spaces over $`N`$, the tubular neighborhood of $`\mathrm{\Sigma }`$. In Section 6, we define an invariant with irreducible boundary values and use properties of this relative invariant to analyze the terms appearing in the product formula, completing the proof of Theorem 1.7. In Section 7, we prove the technical facts about the moduli spaces over $`N`$ which were used in earlier sections. In Section 8, we show how to extend the results of Sections 4 and 7 to deal with the remaining $`\mathrm{Spin}_{}`$ structure. Finally, in Section 9, which should be viewed as an appendix, we discuss representatives for the cohomology classes used throughout the paper.
Acknowledgements. The authors wish to thank Vicente Muñoz for his very helpful comments on an early version of this paper.
## 2. Examples
We give some examples now of four-manifolds $`X`$ which admit basic classes of non-zero dimension. We begin by giving examples where $`b_2^+(X)>1`$ and $`b_1(X)>0`$, to show that the adjunction inequality in Theorem 1.3 is sharp. (It is an open problem whether manifolds with $`b_2^+(X)>1`$ and $`b_1(X)=0`$ can admit basic classes of non-zero dimension.)
### 2.1. Examples of Theorem 1.3 with $`b_2^+(X)>1`$
To construct these examples, we use the following construction.
###### Definition 2.1.
Let $`X`$ be smooth four-manifold and let $`SX`$ be an embedded two-sphere with zero self-intersection number. Let $`X^{}`$ denote a manifold obtained as surgery on $`S`$; i.e.
$$X^{}=(X\mathrm{nd}(S))_\varphi S^1\times D^3,$$
where $`\mathrm{nd}(S)`$ is an open tubular neighborhood of $`S`$ and $`\varphi :(X\mathrm{nd}(S))S^1\times S^2`$ is a orientation-reversing diffeomorphism. Note that up to isotopy there are two possible choices for $`\varphi `$. Let $`CX^{}`$ denote the closed curve which is the core of the added $`S^1\times D^3`$. Note that there is a diffeomorphism $`XSX^{}C`$.
###### Proposition 2.2.
Let $`X`$ be a closed, smooth, oriented four-manifold with $`b_2^+(X)>1`$, and let $`SX`$ be a homologically trivial embedded two-sphere. For each $`\mathrm{Spin}_{}`$ structure $`𝔰`$ on $`X`$, there is a unique induced $`\mathrm{Spin}_{}`$ structure $`𝔰^{}`$ on $`X^{}`$ with the property that
$$𝔰|_{XS}=𝔰^{}|_{X^{}C}.$$
Then, $`d(𝔰^{})=d(𝔰)+1`$; and for all $`a𝔸(X)`$
$$SW_{X^{},𝔰^{}}(a\mu (C))=SW_{X,𝔰}(a),$$
for some homology orientation on $`X^{}`$.
Proof. The dimension statement is straightforward.
To prove the relation, we pull $`X`$ apart along $`S^1\times S^2=\mathrm{nd}(S)`$, and study the corresponding moduli spaces (see Section 5 for more discussion on such matters). Let $`X_0`$ denote the complement $`XS`$, given a cylindrical-end metric modeled on the product metric $`[0,\mathrm{})\times S^1\times S^2`$, where $`S^2`$ is given its standard, round metric. Note that this metric can be extended over both $`S^1\times D^3`$ and $`D^2\times S^2`$ to give metrics with non-negative scalar curvature. Consequently, the moduli spaces of solutions over $`S^1\times S^2`$, $`S^1\times D^3`$, and $`D^2\times S^2`$ consist entirely of smooth reducibles (i.e. the moduli spaces are identified with $`S^1`$, $`S^1`$, and a point respectively).
Let $`_{X_0}(𝔰_0)`$ denote the moduli space of finite energy solutions to the Seiberg-Witten equations over $`X_0`$ in the $`\mathrm{Spin}_{}`$ structure $`𝔰_0=𝔰|_{X_0}`$. Thus, we can think of the boundary map as a map
$$\rho :_{X_0}(𝔰_0)S^1.$$
Gluing theory gives a diffeomorphism for all sufficiently large $`T>0`$:
$$_{X(T)}(𝔰)\rho ^1(x_0),$$
where $`X(T)`$ denotes the metric on $`X`$ with neck-length $`T`$ and $`x_0S^1`$ corresponds to the unique reducible on $`S^1\times S^2`$ which extends to $`D^2\times S^2`$. Consequently,
(7)
$$SW_{X,𝔰}(a)=_{X_0}(𝔰_0),\mu (a)\mu (C),$$
since $`\mu (C)`$ is represented by the holonomy class around $`C`$ (see Proposition 9.1).
Similarly, gluing gives a diffeomorphism of
$$_{X_0}(𝔰_0)_{X^{}(T)}(𝔰^{}),$$
and consequently
(8)
$$SW_{X^{},𝔰^{}}(aC)=_{X_0}(𝔰_0),\mu (aC).$$
Together, Equations (7) and (8) prove the proposition.
###### Remark 2.3.
Of course, the above result also holds when $`b_2^+(X)=1`$, provided that both invariants are evaluated in the same chamber.
Now we construct our examples. Fix natural numbers $`n`$, $`k`$, and $`m`$ with $`2kn>1`$, and let $`X`$ be the four-manifold $`E(n)\mathrm{\#}m(S^3\times S^1)`$, where $`E(n)`$ is a simply-connected elliptic surface with no multiple fibers and with geometric genus $`n1`$. Let $`\mathrm{\Sigma }_0E(n)`$ denote a symplectic submanifold representing the homology class $`S+kF`$, where $`S`$ and $`F`$ denote the homology classes of a section and a fiber respectively of the elliptic fibration. Let $`T_iX`$ denote a fiber in the elliptic fibration of the $`i^{th}`$ summand $`S^3\times S^1`$. Let $`\mathrm{\Sigma }X`$ denote the internal connected sum of $`\mathrm{\Sigma }_0\mathrm{\#}F_1\mathrm{\#}\mathrm{}\mathrm{\#}F_m`$. Note that $`g(\mathrm{\Sigma })=k+m`$ and $`\mathrm{\Sigma }\mathrm{\Sigma }=2kn0`$. Let $`𝔰`$ be the $`\mathrm{Spin}_{}`$ structure over $`X`$ induced from the canonical $`\mathrm{Spin}_{}`$ structure on $`E(n)`$, and let $`b=𝔸(\mathrm{\Sigma })`$ be the product $`B_1\mathrm{}B_m`$ where $`B_iH_1(X)`$ generates $`H_1`$ of the $`i^{th}`$ copy of $`S^1\times S^3`$. Note that $`d(b)=m`$ and $`\mathrm{\Sigma }`$ has a symplectic basis $`\{A_i,B_i\}_{i=1}^{k+m}`$ for which $`A_1,\mathrm{},A_m`$ are homologically trivial in $`X`$. By Proposition 2.2,
$$SW_{X,𝔰}(B_1\mathrm{}B_m)=1,$$
so the data $`X`$, $`b`$, $`\mathrm{\Sigma }`$, $`𝔰`$ satisfy the hypotheses of Theorem 1.3. In fact, we see that
$$\mathrm{\Sigma }\mathrm{\Sigma }+c_1(𝔰),[\mathrm{\Sigma }]+2d(b)=2g(\mathrm{\Sigma })2,$$
which shows that the inequality of the theorem is sharp, for all choices of $`g(\mathrm{\Sigma })>0`$, $`\mathrm{\Sigma }\mathrm{\Sigma }0`$, and $`d(b)`$.
### 2.2. Ruled Surfaces: The Homological Hypotheses on $`H_1(\mathrm{\Sigma })`$
By looking at ruled surfaces, we show that Inequality (4) is sharp, and hence that some homological hypotheses are necessary for the stronger inequality (which appears in Theorem 1.3) to hold.
As mentioned before, one cannot hope for the adjunction inequality of Theorem 1.3 to be valid without additional topological hypotheses on the inclusion of $`\mathrm{\Sigma }`$ in $`X`$. Indeed, fix $`n0`$ and $`g>0`$, and let $`X`$ be the two-sphere bundle over a surface $`\mathrm{\Sigma }`$ of genus $`g`$, associated to the circle bundle with Euler number $`n`$. In particular, $`X`$ contains an embedded copy of $`\mathrm{\Sigma }`$ with $`\mathrm{\Sigma }\mathrm{\Sigma }=n`$. In the chamber corresponding to $`\mathrm{PD}[\mathrm{\Sigma }]`$, there is a zero-dimensional basic class $`𝔰_0`$ with $`c_1(𝔰_0)=K_X`$, where $`K_X`$ is the canonical class of $`X`$ viewed as Kähler manifold. Moreover, letting $`F`$ be the class of the two-sphere fiber in $`X`$, we see that the moduli space associated to $`𝔰_0+d\mathrm{PD}[F]`$ is identified with $`\mathrm{Sym}^d(\mathrm{\Sigma })`$, and $`U`$ is the symmetric product of the volume form of $`\mathrm{\Sigma }`$ (see Proposition 6.10 for a related discussion). Thus, if $`𝔰=𝔰_0+d\mathrm{PD}[F]`$, then $`SW_𝔰(U^d)0`$, and
$$c_1(𝔰),[\mathrm{\Sigma }]=2d.$$
Clearly, Adjunction Inequality 1.4 is sharp for all values of $`k`$, $`d`$, $`n`$, and $`g`$ provided that $`nk`$, where $`k=c_1(𝔰),[\mathrm{\Sigma }]`$, $`2d=d(b)`$, $`n=\mathrm{\Sigma }\mathrm{\Sigma }`$, and $`g=g(\mathrm{\Sigma })`$. (This construction, strictly speaking, only gives us even values of $`d(b)`$. For odd values, one can attach an $`S^1\times S^3`$.) In particular, we see that some homological criterion on the embedding of $`\mathrm{\Sigma }X`$ is necessary for the stronger Inequality (3) to hold.
## 3. From Product Formulas to Relations
The aim of this section is to outline the proof of Theorem 1.7. By employing the blowup formula in a manner analogous to , we reduce to the case where the self-intersection number of $`\mathrm{\Sigma }`$ is very negative (Proposition 3.1). The invariants in this latter case are studied via a product formula, which we state (and prove in Section 5), whose terms are then related with other Seiberg-Witten invariants of $`X`$. In the end of the section, we discuss the modifications which are needed to prove Theorem 1.4.
We reduce Theorem 1.7 to the following special case.
###### Proposition 3.1.
Theorem 1.7 holds, under the additional hypothesis that
$`0c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]2g(\mathrm{\Sigma })2`$ and $`[\mathrm{\Sigma }][\mathrm{\Sigma }]>2g2.`$
The reduction involves the following basic result of Fintushel and Stern.
###### Theorem 3.2.
(Blowup Formula) and . Let $`X`$ be a smooth, closed four-manifold, and let $`\widehat{X}=X\mathrm{\#}\overline{}^2`$ denote its blow-up, with exceptional class $`EH^2(\widehat{X};)`$. If $`b_2^+(X)>1`$, then for each $`\mathrm{Spin}_{}`$ structure $`\widehat{𝔰}`$ on $`\widehat{X}`$ with $`d(\widehat{𝔰})0`$, and each $`a𝔸(X)𝔸(\widehat{X})`$, we have
$$SW_{\widehat{X},\widehat{𝔰}}(a)=SW_{X,𝔰}(U^ma),$$
where $`𝔰`$ is the $`\mathrm{Spin}_{}`$ structure induced on $`X`$ obtained by restricting $`\widehat{𝔰}`$, and $`2m=d(𝔰)d(\widehat{𝔰})`$. If $`b_2^+(X)=1`$, there is a one-to-one correspondence between components of $`\mathrm{\Omega }^+(X)`$ and $`\mathrm{\Omega }^+(\widehat{X})`$, and the above relation holds provided both invariants are calculated in chambers associated to corresponding components.
Before showing how to reduce Theorem 1.7 to the special case, we point out that another special case of Theorem 1.7 was already proved in Theorem 1.3 . More specifically, the following was shown:
###### Theorem 3.3.
Let $`X`$ be a smooth, closed, connected, oriented four-manifold with $`b_2^+(X)>0`$. Let $`\mathrm{\Sigma }X`$ be a surface with genus $`g(\mathrm{\Sigma })>0`$ and negative self-intersection. For each $`\mathrm{Spin}_{}`$ structure $`𝔰`$ with
$$c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]>2g(\mathrm{\Sigma })2,$$
there is an element $`b^{}𝔸(\mathrm{\Sigma })`$ so that that for any $`a𝔸(X)`$, we have
$$SW_{X,𝔰}(a)=SW_{X,𝔰\mathrm{PD}[\mathrm{\Sigma }]}(ai_{}(b^{})).$$
Furthermore, $`b^{}U^{d^{}/2}`$ lies in the ideal generated by $`H_1(\mathrm{\Sigma })`$ in $`𝔸(\mathrm{\Sigma })`$.
###### Remark 3.4.
In the language of Theorem 1.7, this case corresponds to $`\mathrm{}=0`$ and $`b=1`$.
Proposition 3.1 $``$ Theorem 1.7. Let $`g=g(\mathrm{\Sigma })`$, fix an integer $`m`$ with
$$m>[\mathrm{\Sigma }][\mathrm{\Sigma }]+2g2,$$
let $`\widehat{X}=X\mathrm{\#}m\overline{}^2`$, and let $`\widehat{\mathrm{\Sigma }}`$ be the “proper transform” of $`\mathrm{\Sigma }`$, the embedded surface obtained by internal connected sum of $`\mathrm{\Sigma }`$ with the $`m`$ exceptional spheres in the $`\overline{}^2`$ summands; i.e.
$$\mathrm{PD}[\widehat{\mathrm{\Sigma }}]=\mathrm{PD}[\mathrm{\Sigma }]E_1\mathrm{}E_m.$$
Finally, let $`\widehat{𝔰}`$ denote the $`\mathrm{Spin}_{}`$ structure on $`\widehat{X}`$ which agrees with $`𝔰`$ in the complement of the exceptional spheres, whose Chern class satisfies
$$c_1(\widehat{𝔰})=c_1(𝔰)E_1\mathrm{}E_m.$$
It is easy to check that:
$`[\widehat{\mathrm{\Sigma }}][\widehat{\mathrm{\Sigma }}]=m[\mathrm{\Sigma }][\mathrm{\Sigma }]`$ $`>`$ $`2g2,`$
$`c_1(\widehat{𝔰}),[\widehat{\mathrm{\Sigma }}]+[\widehat{\mathrm{\Sigma }}][\widehat{\mathrm{\Sigma }}]`$ $`=`$ $`c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }];`$
$`d(𝔰)`$ $`=`$ $`d(\widehat{𝔰}).`$
Now, if
$$c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]2g,$$
the hypotheses of Theorem 3.3 are satisfied; and otherwise, the hypotheses of Proposition 3.1 are. In either case, for each $`b𝔸(\mathrm{\Sigma })`$ of degree $`d(b)\mathrm{}`$, we can find $`b^{}𝔸(\mathrm{\Sigma })`$ with
(9)
$$SW_{\widehat{X},\widehat{𝔰}}(ai_{}(b))=SW_{\widehat{X},\widehat{𝔰}\mathrm{PD}[\widehat{\mathrm{\Sigma }}]}(ai_{}(b^{})).$$
According to the blow-up formula,
(10)
$$SW_{\widehat{X},\widehat{𝔰}}(ai_{}(b))=SW_{X,𝔰}(ai_{}(b));$$
and, since $`\widehat{𝔰}\mathrm{PD}[\widehat{\mathrm{\Sigma }}]`$ agrees with $`𝔰\mathrm{PD}[\mathrm{\Sigma }]`$ away from the exceptional spheres and
$$c_1(\widehat{𝔰}\mathrm{PD}[\widehat{\mathrm{\Sigma }}])=c_1(𝔰\mathrm{PD}[\mathrm{\Sigma }])E_1\mathrm{}E_m,$$
we see from another application of the blowup formula that
(11)
$$SW_{\widehat{X},\widehat{𝔰}\mathrm{PD}[\widehat{\mathrm{\Sigma }}]}(ai_{}(b^{}))=SW_{X,𝔰\mathrm{PD}[\mathrm{\Sigma }]}(ai_{}(b^{})).$$
Theorem 1.7 then follows by combining Equations (9), (10) and (11). $`\mathrm{}`$
We now turn to the special case considered in Proposition 3.1. We will study the Seiberg-Witten invariant of $`X`$ by decomposing it into two pieces
$$X=N_Y(XN),$$
where $`Y`$ a circle bundle over $`\mathrm{\Sigma }`$ (as in the proposition), and $`N`$ is the associated disk bundle. Following , the moduli space of Seiberg-Witten monopoles over $`Y`$ decomposes into an irreducible and a reducible component. (Actually, there is one $`\mathrm{Spin}_{}`$ structure over $`Y`$, where it is necessary to perturb the equations for this decomposition to occur; this perturbation is studied Section 8.) Correspondingly, we construct relative invariants of $`X\mathrm{\Sigma }`$, denoted $`SW_𝔰^{irr}`$ and $`SW_𝔰^{red}`$, arising from the $`L^2`$ moduli spaces on $`X\mathrm{\Sigma }`$ with irreducible and reducible boundary values. In Section 5 (see Lemma 5.6, and the discussion following it), we prove the following:
###### Proposition 3.5.
Suppose
$`0c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]2g(\mathrm{\Sigma })2`$ and $`[\mathrm{\Sigma }][\mathrm{\Sigma }]>2g2.`$
Then,
$$SW_{X,𝔰}=SW_𝔰^{irr}+SW_𝔰^{red}.$$
We can interpret the latter invariant in terms of the closed manifold as follows.
###### Definition 3.6.
Let $`\mathrm{\Sigma }`$ be a surface of genus $`g`$, and let $`\{A_i,B_i\}_{i=1}^g`$ be a standard symplectic basis for $`H_1(\mathrm{\Sigma };)`$. For $`j=0,\mathrm{},g`$, let $`\xi _j([\mathrm{\Sigma }])𝔸(\mathrm{\Sigma })`$ be the degree $`2j`$ component of
$$\underset{i=1}{\overset{g}{}}\left(1+U+A_iB_i\right)𝔸(\mathrm{\Sigma });$$
i.e. $`\xi _0=1`$, $`\xi _1(\mathrm{\Sigma })=gU+A_iB_i`$, …, $`\xi _g(\mathrm{\Sigma })=_{i=1}^g\left(U+A_iB_i\right)`$.
###### Proposition 3.7.
Suppose
$`0c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]2g(\mathrm{\Sigma })2`$ and $`[\mathrm{\Sigma }][\mathrm{\Sigma }]>2g2.`$
Then, letting
$$e=g1+\frac{c_1(𝔰),[\mathrm{\Sigma }][\mathrm{\Sigma }][\mathrm{\Sigma }]}{2},$$
we have that
$$SW_𝔰^{red}(a)=SW_{X,𝔰\mathrm{PD}[\mathrm{\Sigma }]}(a\xi _{g1e}([\mathrm{\Sigma }]))$$
for all $`a𝔸(X)`$.
Furthermore, under the homological condition of Theorem 1.7, we will express $`SW_𝔰^{irr}`$ in terms of $`SW_{X,𝔰\mathrm{PD}[\mathrm{\Sigma }]}`$, as follows.
###### Proposition 3.8.
Suppose
$`0c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]2g(\mathrm{\Sigma })2`$ and $`[\mathrm{\Sigma }][\mathrm{\Sigma }]>2g2,`$
and let $`\mathrm{}`$ be an integer so that there is a symplectic basis $`\{A_i,B_i\}_{i=1}^g`$ for $`H_1(\mathrm{\Sigma })`$ so that $`i_{}(A_i)=0`$ in $`H_1(X;)`$ for $`i=1,\mathrm{},\mathrm{}`$. Then, for each $`b𝔸(\mathrm{\Sigma })`$ of degree $`e<d(b)\mathrm{}`$, there is an element $`b^{\prime \prime }𝔸(\mathrm{\Sigma })`$ so that
$$SW_𝔰^{irr}(ai_{}(b))=SW_{𝔰\mathrm{PD}[\mathrm{\Sigma }]}(ai_{}(b^{\prime \prime })).$$
Furthermore, $`b_2`$ lies in the ideal generated by $`H_1(\mathrm{\Sigma })`$ in $`𝔸(\mathrm{\Sigma })`$.
The proof of Proposition 3.8 is given in the end of Section 6.
Proposition 3.1 follows immediately from Propositions 3.53.8. In the proof of these latter propositions, we will construct a natural Seiberg-Witten-Floer functor for four-manifolds which bound $`Y`$.
Before proceeding, we pause to tie up one more loose end: Theorem 1.4. That result can be reduced to a relation which replaces Theorem 1.7, using the same argument given in the proof of Theorem 1.3. The relevant relation in this case is:
###### Theorem 3.9.
Let $`X`$ be a smooth, closed, connected, oriented four-manifold with $`b_2^+(X)>0`$. Let $`\mathrm{\Sigma }X`$ be a surface with genus $`g(\mathrm{\Sigma })>0`$. For each $`\mathrm{Spin}_{}`$ structure $`𝔰`$ with
$$c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]0$$
and each $`b𝔸(\mathrm{\Sigma })`$ of degree $`d(b)`$ with
(12)
$$c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]+d(b)>2g(\mathrm{\Sigma })2,$$
there is an element $`b^{}𝔸(\mathrm{\Sigma })`$ with $`d(b^{})d(b)`$ so that for any $`a𝔸(X)`$, we have
(13)
$$SW_{X,𝔰}(ai_{}(b))=SW_{X,𝔰\mathrm{PD}[\mathrm{\Sigma }]}(ai_{}(b^{})).$$
Furthermore, if $`b=U^d`$, then $`b^{}U^d^{}`$ lies in the ideal generated by $`H_1(\mathrm{\Sigma })`$ in $`𝔸(\mathrm{\Sigma })`$.
Once again, via the blowup formula, this relation can be reduced to the case where the self-intersection number $`\mathrm{\Sigma }`$ is very negative; i.e.
$`0c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]2g(\mathrm{\Sigma })2`$ and $`[\mathrm{\Sigma }][\mathrm{\Sigma }]>2g2.`$
(compare Proposition 3.1). Like Proposition 3.1, this special case also follows from the product formula in Proposition 3.5, the relation in Proposition 3.7, together with the following analogue of Proposition 3.8 (whose proof is also given in the end of Section 6):
###### Proposition 3.10.
Suppose
$`0c_1(𝔰),[\mathrm{\Sigma }]+[\mathrm{\Sigma }][\mathrm{\Sigma }]2g(\mathrm{\Sigma })2`$ and $`[\mathrm{\Sigma }][\mathrm{\Sigma }]>2g2.`$
Then, for each $`b𝔸(\mathrm{\Sigma })`$ of degree $`2e<d(b)`$, there is an $`b^{\prime \prime }𝔸(\mathrm{\Sigma })`$ so that
$$SW_𝔰^{irr}(ai_{}(b))=SW_{𝔰\mathrm{PD}[\mathrm{\Sigma }]}(ai_{}(b^{\prime \prime })).$$
Furthermore, $`b^{\prime \prime }`$ lies in the ideal generated by $`H_1(\mathrm{\Sigma })`$.
## 4. Gauge theory on $`\times Y`$
The Seiberg-Witten moduli spaces over $`Y`$ and $`\times Y`$ were studied for Seifert fibered three-manifolds $`Y`$ in . We summarize these results here, for $`Y`$ a circle-bundle over a Riemann surface $`\mathrm{\Sigma }`$ with $`g(\mathrm{\Sigma })>0`$ and Euler number $`n`$, where $`n>2g2`$.
$`Y`$ admits a canonical $`\mathrm{Spin}_{}`$ structure whose bundle of spinors is $`\pi ^{}(K_{\mathrm{\Sigma }}^{}{}_{}{}^{1})`$, which we use to identify the $`\mathrm{Spin}_{}`$ structures on $`Y`$ with $`H^2(Y;)^{2g}/n`$.
Let $`𝒩_Y(𝔱)`$ denote the moduli space of solutions to the Seiberg-Witten equations over $`Y`$ in the $`\mathrm{Spin}_{}`$ structure $`𝔱`$. Here, we use the metric $`g_Y`$ and $`SO(3)`$-connection over $`TY`$ of . Given a pair of components $`C_1,C_2`$ in $`𝒩_Y(𝔱)`$, let $`(C_1,C_2)`$ denote the moduli space of solutions $`[A,\mathrm{\Phi }]`$ to the Seiberg-Witten equations on $`\times Y`$ for which
$`\underset{t\mathrm{}}{lim}[A,\mathrm{\Phi }]|_{\{t\}\times Y}C_1,`$ and $`\underset{t\mathrm{}}{lim}[A,\mathrm{\Phi }]|_{\{t\}\times Y}C_2.`$
This moduli space admits a translation action by $``$. Let $`\widehat{}(C_1,C_2)`$ denote the quotient of $`(C_1,C_2)`$ by the translation action.
In general, these spaces admit a Morse-theoretic interpretation. If $`c_1(𝔱)`$ is a torsion class, there is a real-valued functional
$$\mathrm{CSD}:(Y,𝔱)$$
defined over the configuration space $`(Y,𝔱)`$ of pairs $`(B,\mathrm{\Psi })`$ of spin-connections $`B`$ in $`𝔱`$ and spinors $`\mathrm{\Psi }`$ modulo gauge. The critical manifolds are the moduli spaces $`𝒩(Y;𝔱)`$. When $`c_1(𝔱)`$ is not torsion, the functional is circle-valued. The Seiberg-Witten equations on $`\times Y`$ are the upward gradient-flow equations for this functional. In keeping with this interpretation, we call $`\widehat{}(C_1,C_2)`$ the space of unparameterized flows from $`C_1`$ to $`C_2`$.
###### Theorem 4.1.
() Let $`Y`$ be a circle-bundle over a Riemann surface with genus $`g>0`$ and Euler number $`n<22g`$. The moduli space $`𝒩_Y(𝔱)`$ is empty unless $`𝔱`$ corresponds to a torsion class in $`H^2(Y;)`$. So, suppose $`𝔱`$ corresponds to $`e/nH^2(Y;)`$.
* If $`0e<g1`$ then $`𝒩_Y(𝔱)`$ contains two components, a reducible one $`𝒥`$, identified with the Jacobian torus $`H^1(\mathrm{\Sigma };/)`$, and a smooth irreducible component $`C`$ diffeomorphic to $`\mathrm{Sym}^e(\mathrm{\Sigma })`$. Both of these components are non-degenerate in the sense of Morse-Bott. There is an inequality $`\mathrm{CSD}(𝒥)>\mathrm{CSD}(C)`$, so the space $`\widehat{}(𝒥,C)`$ is empty. The space $`\widehat{}(C,𝒥)`$ is smooth of expected dimension $`2e`$; indeed it is diffeomorphic to $`\mathrm{Sym}^e(\mathrm{\Sigma })`$.
* If $`g1<e2g2`$, the Seiberg-Witten moduli spaces over both $`Y`$ and $`\times Y`$ in this $`\mathrm{Spin}_{}`$ structure are naturally identified with the corresponding moduli spaces in the $`\mathrm{Spin}_{}`$ structure $`2g2e`$, which we just described.
* For all other $`eg1`$, $`𝒩_Y(𝔱)`$ contains only reducibles. Furthermore, it is smoothly identified with the Jacobian torus.
In the $`\mathrm{Spin}_{}`$ structure corresponding to $`g1/n`$, the unperturbed Seiberg-Witten equations used in Theorem 4.1 are inconvenient, since the corresponding reducible manifold is not smooth in the sense of Morse-Bott. To overcome this difficulty, when working in this $`\mathrm{Spin}_{}`$ structure, we use a perturbation of the equations where the theory resembles the case where $`0e<g1`$ (and, in particular, the reducibles are smooth). A thorough discussion of the perturbation is given in Section 8.
## 5. The Product Formula
In this section, we define two quantities, $`SW^{irr}`$ and $`SW^{red}`$, and prove that the Seiberg-Witten invariant decomposes into a sum of these (Propostion 3.5). Furthermore, we express $`SW^{red}`$ in terms of another Seiberg-Witten invariant of $`X`$ (Proposition 3.7).
Decompose $`X`$ as
$$X=N_YX_0,$$
where $`Y`$ is unit circle bundle over $`\mathrm{\Sigma }`$ with Euler number $`n`$, with $`n>2g2`$. $`N`$ is a tubular neighborhood of the surface $`\mathrm{\Sigma }`$ (which is diffeomorphic to the disk bundle associated to $`Y`$), and $`X_0`$ is the complement in $`X`$ of the interior of $`N`$. Fix metrics $`g_{X_0}`$, $`g_N`$, and $`g_Y`$ for which $`g_{X_0}`$ and $`g_N`$ are isometric to
$$dt^2+g_Y^2$$
in a collar neighborhood of their boundaries (where $`t`$ is a normal coordinate to the boundary). Let $`X(T)`$ denote the Riemannian manifold which is diffeomorphic to $`X`$ and whose metric $`g_T`$ is obtained from the description
$$X(T)=N_{N=\{T\}\times Y}[T,T]\times Y_{\{T\}\times Y=X_0}X_0;$$
i.e. $`g_T|_N=g_N`$, $`g_T|_{[T,T]\times Y}=dt^2+g_Y^2`$, and $`g_T|_{X_0}=g_{X_0}`$. Our goal here is to provide, for all sufficiently large $`T`$, a description of the moduli space $`_{X(T)}(𝔰)`$ on $`X(T)`$ in terms of the moduli spaces for $`Y`$, $`𝒩_Y(𝔰|_Y)`$, and the finite-energy, cylindrical-end moduli spaces associated to $`X_0`$ and $`N`$, denoted $`_{X_0}(𝔰|_{X_0})`$, and $`_N(𝔰|_N)`$ respectively. In this context, finite energy means that the total variation of the Chern-Simons-Dirac functional over the infinite cylinder is bounded. Henceforth, $`X_0`$ and $`N`$ will denote the cylindrical-end manifolds obtained by attaching $`[0,\mathrm{})\times Y`$ (with appropriate orientations) to the corresponding subsets of $`X`$.
In the case where $`b_2^+(X)=1`$, we choose the perturbing form $`\eta `$ to be compactly supported in $`X_0`$ in such a way that
$$2\pi c_1(𝔰)\omega _{\mathrm{}}+_{X_0}\eta \omega _{\mathrm{}}$$
has the same sign as $`\gamma \omega _{\mathrm{}}`$, where $`\gamma `$ is a compactly supported representative for a class in the chosen component $`𝒦_0𝒦(X)`$, and $`\omega _{\mathrm{}}`$ is a self-dual harmonic two-form over $`X_0`$ with $`_{X_0}\omega _{\mathrm{}}\omega _{\mathrm{}}=1`$. Note that such a $`\gamma `$ and $`\omega _{\mathrm{}}`$ can be found since $`\mathrm{\Sigma }\mathrm{\Sigma }<0`$, forcing $`b_2^+(X_0)=1`$ (see ). Now, the moduli spaces of the $`\eta `$-perturbed Seiberg-Witten equations over $`X(T)`$ calculate the invariant in the chosen chamber for all sufficiently large $`T`$.
We collect useful facts about the moduli spaces $`_N(𝔰|_N)`$, most of which we defer to Section 7 (see also ), but first we introduce some notation. The map
$$\mathrm{Spin}_{}(N)$$
given by
$$𝔰c_1(𝔰),[\mathrm{\Sigma }]$$
induces a one-to-one correspondence between $`\mathrm{Spin}_{}`$ structures and integers which are congruent to $`n`$ modulo $`2`$. Note that the $`\mathrm{Spin}_{}`$ structure over $`Y`$ $`𝔰|_Y`$ corresponds to the mod $`n`$ reduction of
$$e=g1+\frac{c_1(𝔰),[\mathrm{\Sigma }]+n}{2}$$
appearing in Theorem 4.1.
By taking limits at the end of the tube, one can define maps
$`\rho :_N(𝔰)𝒩_Y(𝔰|_Y)`$ and $`\rho :_{X_0}(𝔰)𝒩_Y(𝔰|_Y)`$
(see ). If $`C`$ is a connected manifold of $`𝒩_Y(𝔰|_Y)`$, then $`_N(𝔰,C)`$ and $`_{X_0}(𝔰,C)`$ denotes the pre-image of $`C`$ under $`\rho `$. Throughout the following discussion, we will use the perturbation discussed in Section 8 over $`X_0`$, $`N`$, and $`Y`$, when $`𝔰|_Y`$ corresponds to $`e=g1`$ (in the notation of Section 4); i.e. in this case, $`𝒩_Y(𝔰|_Y)`$, $`_N(𝔰)`$ and $`_{X_0}(𝔰|_{X_0})`$ will denote the perturbed versions of these moduli spaces, with perturbation parameter $`u`$ in the range $`0<u<2`$, in the notation of Section 8. (We will show in Section 8 that this is an allowable perturbation to use when $`b_2^+(X)=1`$; i.e. we are computing the Seiberg-Witten invariants in the correct chamber.) When they are clear from the context, we leave the $`\mathrm{Spin}_{}`$ structures out of the notation. Note that on the cylinders, the analogous boundary value maps factor through the unparameterized spaces, defining
$`\rho _𝒥:\widehat{}(C,𝒥)𝒥`$ and $`\rho _C:\widehat{}(C,𝒥)C,`$
where $`𝒥`$ and $`C`$ are the critical manifolds of Theorem 4.1.
###### Proposition 5.1.
Suppose that $`n2g+2c_1(𝔰),[\mathrm{\Sigma }]n`$, and let
$$e=g1+\frac{c_1(𝔰),[\mathrm{\Sigma }]+n}{2},$$
Then according to Theorem 4.1, and Theorem 8.1 when $`e=g1`$, $`𝒩_Y(𝔰|_Y)`$ has two components, $`𝒥`$ and $`C`$, where $`C`$ is diffeomorphic to $`\mathrm{Sym}^e(\mathrm{\Sigma })`$. Furthermore, the expected dimensions of the moduli spaces over $`N`$ and $`X_0`$ are given by:
(14) $`\text{e-dim}_N(𝒥)`$ $`=`$ $`2e+1`$
(15) $`\text{e-dim}_N(C)`$ $`=`$ $`2e`$
(16) $`\text{e-dim}_{X_0}(𝒥)`$ $`=`$ $`2d+2g2e2`$
(17) $`\text{e-dim}_{X_0}(C)`$ $`=`$ $`2d,`$
where $`d=d(𝔰)`$ and $`g=g(\mathrm{\Sigma })`$. Moreover, $`_N^{}(𝒥)`$, $`_N(C)`$, $`_{X_0}(𝒥)`$, and $`_{X_0}(C)`$ are transversally cut out by the Seiberg-Witten equations (in particular, they are manifolds of the expected dimension).
Proof. This is a combination of Proposition 7.9 and 7.10 when $`c_1(𝔰),[\mathrm{\Sigma }]n`$, and Proposition 8.3 in the remaining case.
When studying the deformation theory of reducibles inside $`_N(𝒥)`$, the kernel and the cokernel of the Dirac operator play a central role. These spaces can be concretely understood, thanks to the holomorphic interpretation of the Dirac operator (see also ).
###### Proposition 5.2.
Suppose that $`n2g+2c_1(𝔰),[\mathrm{\Sigma }]n`$, then there is a natural correspondence between reducibles $`[(A,0)]_N(𝒥)`$ with holomorphic line bundles $``$ of degree $`e`$ over $`\mathrm{\Sigma }`$ which identifies
$`\mathrm{Ker}\overline{)}D_A=H^0(\mathrm{\Sigma },)`$ and $`\mathrm{Coker}\overline{)}D_A=H^1(\mathrm{\Sigma },).`$
Proof. This follows from Theorem 7.4 and Proposition 7.5 (see also the proof of Theorem 8.1 in the perturbed case).
The above proposition allows us to understand an important class of reducibles.
###### Definition 5.3.
The jumping locus $`\mathrm{\Theta }_N(𝒥)`$ is the locus of reducible solutions $`[(A,0)]_N(𝒥)`$ for which $`\mathrm{Ker}\overline{)}D_A`$ is non-trivial.
###### Corollary 5.4.
Suppose that $`n2g+2c_1(𝔰),[\mathrm{\Sigma }]n`$, then the jumping locus $`\mathrm{\Theta }𝒥=_N^{\mathrm{red}}(𝒥)`$ is the image of a smooth map $`\mathrm{Sym}^e(\mathrm{\Sigma })𝒥`$.
Proof. According to Proposition 5.2, the space $`\mathrm{\Theta }𝒥`$ is identified with the space of degree $`e`$ line bundles over $`\mathrm{\Sigma }`$ with non-trivial $`H^0`$. The forgetful map $`\mathrm{Sym}^e(\mathrm{\Sigma })𝒥`$ which takes a degree $`e`$ divisor, thought of as a complex line bundle with section, to the underlying complex line bundle gives the surjection to this locus.
We will also need to understand those $`\mathrm{Spin}_{}`$ structures $`𝔰\mathrm{Spin}_{}(N)`$ for which $`n<c_1(𝔰),[\mathrm{\Sigma }]n`$.
###### Proposition 5.5.
If
$$n<c_1(𝔰),[\mathrm{\Sigma }]n,$$
then the moduli space $`_N(𝒥)`$ contains only reducibles. Moreover, the space of reducibles is smoothly identified with the Jacobian torus $`𝒥`$ (i.e. the kernel and the cokernel of the Dirac operator coupled to any reducible vanishes). Furthermore, $`_N(C)`$ is empty.
Proof. When $`|c_1(𝔰),[\mathrm{\Sigma }]|<n`$, this is proved in Section 7, where it appears as Proposition 7.6. The remaining case is covered by Proposition 8.2.
With these preliminaries in place, we turn to the Seiberg-Witten invariants of $`X`$, by investigating the moduli spaces over $`X(T)`$. Specifically, choose some $`a𝔸(X)`$ of degree $`d(𝔰)`$, and indeed choose representatives for the corresponding homology classes which are compactly supported in $`X_0`$. Let $`V(a)`$ denote the corresponding representatives for $`\mu (a)`$ in the configuration spaces for $`X_0`$ and $`X(T)`$ as appropriate (see Section 9 for a discussion of such representatives). Recall that $`SW_{X,𝔰}(a)`$ is the number of points in $`_{X(T)}(𝔰)V(a)`$, counted with appropriate sign.
###### Lemma 5.6.
Suppose that $`n2g+2c_1(𝔰),[\mathrm{\Sigma }]n`$, then for each $`ϵ>0`$, there is a $`T_0>0`$ so that for all $`TT_0`$ the restriction of $`[(A,\mathrm{\Phi })]_{X(T)}(𝔰)V(a)`$ to any slice $`\{t\}\times Y`$ with $`t[T_0,T_0]`$ lies within $`ϵ`$ (in the $`C^{\mathrm{}}`$ topology) from either $`𝒥`$ or $`C`$. Accordingly, if $`ϵ`$ is sufficiently small, then $`[(A,\mathrm{\Phi })]`$ satisfies exactly one of the following two conditions:
* $`[(A,\mathrm{\Phi })]|_N`$ is $`C^{\mathrm{}}`$ close to smooth reducible and $`[(A,\mathrm{\Phi })]|_{X_0}`$ is $`C^{\mathrm{}}`$ close to (the restriction to $`X_0`$) of a configuration in $`_{X_0}(𝒥)V(a)`$;
* $`[(A,\mathrm{\Phi })]|_N`$ is $`C^{\mathrm{}}`$ close to a configuration in $`_N(C)`$, and $`[(A,\mathrm{\Phi })]|_{X_0}`$ is $`C^{\mathrm{}}`$ close to a configuration in the cut-down moduli space $`_{X_0}(C)V(a)`$
Proof. This is a dimension-counting argument. Suppose we have a sequence $`[A_i,\mathrm{\Phi }_i]_{X(T_i)}(𝔰)V(a)`$, for some increasing, unbounded sequence $`\{T_i\}_{i=1}^{\mathrm{}}`$ of real numbers. By local compactness, there is a subsequence which converges in $`C_{\mathrm{loc}}^{\mathrm{}}`$ to a pair of configurations $`(A_N,\mathrm{\Phi }_N)`$ and $`(A_{X_0},\mathrm{\Phi }_{X_0})`$ over $`N`$ and $`X_0`$ respectively. By the usual compactness arguments (see ), the total variation of the Chern-Simons-Dirac functional of $`(A_i,\mathrm{\Phi }_i)`$ over the cylinder $`[T_i,T_i]\times Y`$ remains globally bounded (independent of $`i`$), so $`(A_N,\mathrm{\Phi }_N)`$ and $`(A_{X_0},\mathrm{\Phi }_{X_0})`$ both have finite energy.
First, we prove that either Hypothesis (H-1) or (H-2) is satisfied. There are a priori four cases, according to which critical manifolds $`\rho [A_{X_0},\mathrm{\Phi }_{X_0}]`$ and $`\rho [A_N,\mathrm{\Phi }_N]`$ lie in.
* The case where $`\rho [A_N,\mathrm{\Phi }_N]𝒥`$ while $`\rho (A_{X_0},\mathrm{\Phi }_{X_0})C`$ is excluded because $`\mathrm{CSD}(C)>\mathrm{CSD}(𝒥)`$.
* The case where $`\rho [A_N,\mathrm{\Phi }_N]C`$ while $`\rho (A_{X_0},\mathrm{\Phi }_{X_0})𝒥`$ is excluded by a dimension count, as follows. In this case, we see that $`\rho [A_{X_0},\mathrm{\Phi }_{X_0}]\rho _𝒥((C,𝒥))\rho (_{X_0}(𝒥)V(a))`$. But
$$\rho _𝒥((C,𝒥))=\rho _𝒥(\widehat{}(C,𝒥)),$$
so
$`\text{e-dim}(\rho _𝒥((C,𝒥))\rho (_{X_0}(𝒥)V(a))`$ $`=`$ $`2.`$
It follows from Theorems 4.1 and 8.1 that $`(C,𝒥)`$ is smooth of the expected dimension, so from the usual transversality results, the above intersection is generically empty.
* Suppose that $`\rho (A_N,\mathrm{\Phi }_N)𝒥`$ and $`\rho (A_{X_0},\mathrm{\Phi }_{X_0})𝒥`$. Then we see that
$$\rho [A_N,\mathrm{\Phi }_N]=\rho [A_{X_0},\mathrm{\Phi }_{X_0}]\rho (_N(𝒥))\rho (_{X_0}(𝒥)V(a));$$
but, according to Proposition 5.1
$`\text{e-dim}\rho (_N^{}(𝒥))\rho (_{X_0}(𝒥)V(a))`$
$`=`$ $`\text{e-dim}_N(𝒥)+_{X_0}(𝒥)2d2g`$
$`=`$ $`1,`$
which is generically empty. Thus, it follows that $`[A_N,\mathrm{\Phi }_N]`$ must be reducible. Moreover, according to Proposition 5.4,
$`\text{e-dim}\rho (\mathrm{\Theta })\rho (_{X_0}(𝒥)V(a))`$ $`=`$ $`2e+\text{e-dim}_{X_0}(𝒥)2d2g`$
$`=`$ $`2,`$
which is also generically empty. Hence, $`[A_N,\mathrm{\Phi }_N]`$ and $`[A_{X_0},\mathrm{\Phi }_{X_0}]`$ satisfy Hypotheses (H-1).
* If $`\rho [A_N,\mathrm{\Phi }_N]`$ and $`\rho [A_{X_0},\mathrm{\Phi }_{X_0}]`$ both lie in $`C`$, then the Hypotheses (H-2) are satisfied.
The assertion at the beginning of the proposition follows easily.
The above proposition says that we can partition the points in the cut-down moduli space (which is an oriented, zero-dimensional manifold) for sufficiently large $`T`$ into two disjoint sets, the subsets of configurations which satisfy (H-1) and (H-2) respectively. Thus, if we let $`SW_𝔰^{red}(a)`$ and $`SW_𝔰^{irr}(a)`$ be the signed number of points satisfying (H-1) and (H-2) respectively, then
(18)
$$SW_{X,𝔰}(a)=SW_𝔰^{red}(a)+SW_𝔰^{irr}(a).$$
As we shall see, gluing theory allows us to compute both of these quantities in terms of cylindrical-end moduli spaces. So, in the next step, we study these cylindrical-end moduli spaces.
###### Lemma 5.7.
For all $`\mathrm{Spin}_{}`$ structures $`𝔰`$ on $`X`$ the corresponding moduli spaces $`_N(C)`$, $`_{X_0}(𝒥)`$, and $`_{X_0}(C)V(a)`$ are all compact manifolds.
Proof. The compactness of $`_{X_0}(𝒥)`$ and $`_N(C)`$ follows from the usual compactness arguments, together with the facts that the Chern-Simons-Dirac functional is real-valued, $`\mathrm{CSD}(𝒥)>\mathrm{CSD}(C)`$, and there are no other critical manifolds. Compactness of $`_{X_0}(C)V(a)`$ follows from this, together with a straightforward dimension count (see the discussion above in the proof of Lemma 5.6, part (P-2)).
Compactness of $`_{X_0}(𝒥)`$ allows us to define a relative invariant with reducible boundary values. We pause to discuss some relevant properties of this invariant.
###### Definition 5.8.
Let $`𝔰_0`$ be a $`\mathrm{Spin}_{}`$ structure on $`X_0`$ which extends over $`X`$. Since the moduli space $`_{X_0,𝔰_0}(𝒥)`$ is compact, there is a relative Seiberg-Witten invariant
$$SW_{(X_0,𝔰_0,𝒥)}:𝔸(X_0),$$
defined by the pairing $`SW_{(X_0,𝔰_0,𝒥)}(a)=[_{X_0,𝔰_0}(𝒥)],\mu (a).`$
This relative invariant is related to an absolute invariant, according to the following.
###### Proposition 5.9.
If $`𝔰`$ satisfies $`n<c_1(𝔰),[\mathrm{\Sigma }]n`$, then for all $`a𝔸(X)`$,
$$SW_{X,𝔰}(a)=SW_{(X_0,𝔰_0,𝒥)}(a),$$
where $`𝔰_0=𝔰|_{X_0}`$.
Proof. Recall that $`(𝔰|_N)`$ consists entirely of reducibles all of which are smooth, according to Proposition 5.5; thus, gluing theory identifies the moduli spaces $`_{X(T)}(𝔰)`$ for large $`T`$ with $`_{X_0,𝔰_0}(𝒥)`$. (See also , where this result appears as Proposition 2.7.)
We now return to the discussion of $`SW^{red}`$ and $`SW^{irr}`$. Although the definitions of both terms implicitly use $`T`$, we show now that if $`T`$ is sufficiently large, then the terms can be computed from absolute invariants (and hence are independent of the parameter).
###### Proposition 5.10.
Suppose that $`𝔰`$ satisfies
$$n2g+2c_1(𝔰),[\mathrm{\Sigma }]n,$$
where $`\mathrm{\Sigma }`$ has self-intersection number $`n`$, and let
$$e=g1+\frac{c_1(𝔰),[\mathrm{\Sigma }]+n}{2}.$$
Then, for all sufficiently large $`T`$,
$$SW_𝔰^{red}(a)=SW_{X,𝔰\mathrm{PD}[\mathrm{\Sigma }]}(a\xi _{g1e}(\mathrm{\Sigma })),$$
where $`\xi _{g1e}(\mathrm{\Sigma })𝔸(\mathrm{\Sigma })`$ is the element defined in Definition 3.6.
Proof. The moduli space $`_N^{\mathrm{red}}(𝒥)\mathrm{\Theta }`$ comes equipped with an obstruction bundle $`\mathrm{\Xi }_N^{\mathrm{red}}(𝒥)\mathrm{\Theta }`$, defined by $`\mathrm{\Xi }_{[(A,0)]}=\mathrm{Coker}\overline{)}D_A`$. (whose $`K`$-theory class canonically extends over all of $`_N^{\mathrm{red}}(𝒥)`$). The dimension count in Lemma 5.6 guarantees that each solution in $`_{X_0}(𝒥)V(a)`$ extends uniquely to a smooth reducible over $`N`$. Thus, gluing theory gives that
$`SW_𝔰^{red}(a)`$ $`=`$ $`[_{X_0}(𝒥)V(a)],𝐞(\rho ^{}(\mathrm{\Xi }))`$
$`=`$ $`[_{X_0}(𝒥)],\mu (a)𝐞(\rho ^{}(\mathrm{\Xi })),`$
where $`𝐞`$ denotes the Euler class of a bundle (or $`K`$-theory element).The Riemann-Roch formula says $`dim(\mathrm{\Xi })=2g22e`$. Using the index theorem for families, together with the holomorphic interpretation of the obstruction bundle $`\mathrm{\Xi }`$ given in Proposition 5.2, it is a straightforward computation that the total Chern class of $`\mathrm{\Xi }`$ is
$$\underset{i=1}{\overset{g}{}}(1+\mu (A_i)\mu (B_i))$$
(see also Proposition 2.6); thus,
$$𝐞(\rho ^{}(\mathrm{\Xi }))=c_{g1e}(\rho ^{}(\mathrm{\Xi }))=\xi _{g1e}([\mathrm{\Sigma }]).$$
Putting all this together, we have that
(19)
$$SW_𝔰^{red}(a)=SW_{(X_0,𝔰_0,𝒥)}(a\xi _{g1e}(\mathrm{\Sigma })),$$
where $`𝔰_0=𝔰|_{X_0}`$. Since $`n2g+2c_1(𝔰\mathrm{PD}[\mathrm{\Sigma }]),[\mathrm{\Sigma }]n`$ and $`𝔰\mathrm{PD}[\mathrm{\Sigma }]|_{X_0}=𝔰_0`$, the proposition then follows from Proposition 5.9.
###### Proposition 5.11.
For sufficiently large $`T`$,
$$SW_𝔰^{irr}(a)=\mathrm{\#}_{X_0}(C)V(a).$$
Proof. Gluing shows that
$$SW_𝔰^{irr}(a)=\left(\mathrm{\#}_{X_0}(C)V(a)\right)\left(\mathrm{deg}(\rho :_N(C)C)\right).$$
According to Propositions 7.9 and 8.3, $`\rho :_N(C)C`$ either has degree $`+1`$, or $`_N(C)`$ is empty. The latter case would force $`SW_𝔰^{irr}(a)0`$ (for the given genus and self-intersection number).
To rule out this latter case, we need only look at an example where the irreducible term is non-zero. Let $`X`$ be a ruled surface $`X`$ over $`\mathrm{\Sigma }`$ associated to the line bundle with Euler number $`n`$. Let $`\mathrm{\Sigma }X`$ denote the section with self-intersection number $`n`$, and fix any $`0eg1`$. Let $`𝔰`$ denote the $`\mathrm{Spin}_{}`$ structure over $`X`$ given by $`𝔰=𝔰_0+e\mathrm{PD}[F]`$, where $`𝔰_0`$ is the canonical $`\mathrm{Spin}_{}`$ structure on $`X`$ associated to the Kähler structure, and $`F`$ denotes a fiber in the ruling. It is easy to see that $`SW_{X,𝔰\mathrm{PD}[\mathrm{\Sigma }]}0`$, as the corresponding space of divisors is empty (see Proposition 7.5). Moreover, we know that $`SW_{X,𝔰}0`$ (compare Example 2.2). Thus, in light of Equation (18) and Proposition 5.10, we have examples where $`SW^{irr}0`$, forcing the degree to be non-zero.
We will give the seemingly ad hoc quantity $`\mathrm{\#}_{X_0}(C)V(a)`$ a more intrinsic formulation in Section 6. With the help of this formulation, we can then prove a vanishing result for this term under suitable algebro-topological hypotheses on the embedding of $`\mathrm{\Sigma }X`$ (Proposition 3.8).
## 6. Relative Invariants
Let $`X_0`$ be a smooth, oriented manifold-with-boundary with $`b_2^+(X_0)>0`$, whose boundary is identified with $`X_0=Y`$, a circle bundle over a Riemann surface $`\mathrm{\Sigma }`$ of genus $`g>0`$ with Euler number $`n`$, where $`n>2g2`$.
In Section 5, we studied the moduli space $`_{X_0}(C)`$, and used it to define a relative invariant
$$SW_𝔰^{irr}:𝔸(X_0),$$
by cutting down the moduli space $`_{X_0}(C)`$ by submanifolds representing $`\mu (a)`$ which are induced from compactly supported representatives for homology in $`X_0`$ (see Proposition 5.11). When $`a𝔸(Y)`$, there are alternate representatives which are supported “at infinity.” The advantage of these representatives is that the corresponding relative invariant inherits relations arising from the cohomology ring of $`C`$. In view of the non-compactness of $`_{X_0}(C)`$, the two types of representatives do not necessarily give rise to the same invariant. However, the difference can be explicitly computed in terms of other Seiberg-Witten invariants. In this section, we recast this discussion in a more algebraic setting, defining an invariant
$$SW_{(X_0,C)}:𝔸(X_0)H^{}(C).$$
which simultaneously captures both types of representatives; in particular,
$$SW_𝔰^{irr}(a)=SW_{(X_0,C)}(a1).$$
Proposition 3.8 then follows from properties of this invariant.
A subtlety arises in the definition of $`SW_{(X_0,C)}`$, since the moduli space $`_{X_0}(C)`$ is not compact. However, we have the following weak compactness theorem.
###### Definition 6.1.
A sequence of configurations $`\{[A_i,\mathrm{\Phi }_i]\}_{i=1}^{\mathrm{}}`$ is said to converge weakly to a configuration
$$[B,\mathrm{\Psi }]\times [A,\mathrm{\Phi }]\widehat{}(C,𝒥)\times _𝒥_{X_0}(𝒥)$$
if $`[A_i,\mathrm{\Phi }_i]`$ converges to $`[A,\mathrm{\Phi }]`$ in $`C_{\mathrm{loc}}^{\mathrm{}}`$, and there is an increasing, unbounded sequence of real numbers $`\{T_i\}_{i=1}^{\mathrm{}}`$ with $`T_i>i`$, so that the translates of $`\{[A_i,\mathrm{\Phi }_i]|_{[0,2T_i]\times Y}\}_{i=1}^{\mathrm{}}`$, viewed as a sequence of configurations on $`[T_i,T_i]\times Y`$, converge in $`C_{\mathrm{loc}}^{\mathrm{}}`$ to a configuration which is equivalent (under translations) to $`[B,\mathrm{\Psi }]`$.
###### Proposition 6.2.
Weak convergence gives the space
$$\overline{}_{X_0}(C)=_{X_0}(C)\widehat{}(C,𝒥)\times _𝒥_{X_0}(𝒥)$$
the structure of a compact Hausdorff space.
Proof. This a standard argument from Morse-Floer theory. A general discussion of compactness results for the anti-self-duality equation can be found in (see especially Theorem 6.3.3 of ); so we sketch the argument here only briefly.
A sequence $`[A_i,\mathrm{\Phi }_i]_{X_0}(C)`$ converges in $`C_{\mathrm{loc}}^{\mathrm{}}`$, after passing to a subsequence, to some solution $`[A,\mathrm{\Phi }]`$ to the Seiberg-Witten equations on $`X_0`$. Since each of the $`[A_i,\mathrm{\Phi }_i]`$ have finite energy, so does $`[A,\mathrm{\Phi }]`$; thus, it has a boundary value. If $`\rho [A,\mathrm{\Phi }]C`$, then the length-energy estimates of L. Simon ( ) can be used to show that the convergeance is $`C^{\mathrm{}}`$ as in .
If, on the other hand, $`\rho [A,\mathrm{\Phi }]C`$, it must be the case that $`\rho [A,\mathrm{\Phi }]𝒥`$. Now, let $`T_i`$ be the number so that
$$\mathrm{CSD}[A_i,\mathrm{\Phi }_i]_{\{T_i\}\times Y}=\frac{\mathrm{CSD}(𝒥)+\mathrm{CSD}(C)}{2}.$$
Clearly, $`T_i\mathrm{}`$. After passing to a subsequence, we can find a configuration $`[B,\mathrm{\Psi }]`$ so that the sequence $`[A_i,\mathrm{\Phi }_i]|_{[0,2T_i]\times Y}`$, viewed as a sequence of configurations over $`[T_i,T_i]`$, converges in $`C_{\mathrm{loc}}^{\mathrm{}}`$ to $`[B,\mathrm{\Psi }]`$. In fact, $`[B,\mathrm{\Psi }]`$ must solve the Seiberg-Witten equations and it must have finite energy, so $`[B,\mathrm{\Psi }](C,𝒥)`$. The usual length-energy estimates then guarantee that the boundary values match up.
The topological space $`\overline{}_{X_0}(C)`$ defined in Proposition 6.2 is called the compactified moduli space. The following result follows immediately from its definition.
###### Proposition 6.3.
The inclusion maps
$$i:_{X_0}(C)^{}(X_0(0,\mathrm{})\times Y)$$
and
$$i\mathrm{\Pi }_2:\widehat{}(C,𝒥)\times _𝒥_{X_0}(𝒥)^{}(X_0(0,\mathrm{})\times Y)$$
fit together to give a continuous map
$$\overline{i}:\overline{}_{X_0}(C)^{}(X_0(0,\mathrm{})\times Y),$$
where $`^{}`$ denotes the irreducible configurations.
Similarly, we can extend the restriction map over the compactified moduli space, as follows.
###### Proposition 6.4.
The restriction maps
$$\rho _C:_{X_0}(C)C$$
and
$$\rho _C\mathrm{\Pi }_1:\widehat{}(C,𝒥)\times _𝒥_{X_0}(𝒥)C$$
fit together to give a continuous map
$$\overline{\rho }_C:\overline{}_{X_0}(C)C.$$
Proof. If a sequence $`[A_n,\mathrm{\Phi }_n]_{X_0}(C)`$ converges to an ideal point
$$[B,\mathrm{\Psi }]\times [A,\mathrm{\Phi }]\widehat{}(C,𝒥)\times _𝒥_{X_0}(𝒥),$$
then there is a divergent sequence $`\{T_n\}_{n=1}^{\mathrm{}}`$ of real numbers so that
$$\underset{n\mathrm{}}{lim}\tau _n^{}[A_n,\mathrm{\Phi }_n]|_{[T_n,\mathrm{})\times Y}=[B,\mathrm{\Psi }]|_{[0,\mathrm{})\times Y},$$
where
$$\tau _n:[0,\mathrm{})\times Y[T_n,\mathrm{})\times Y$$
is the map induced by translation by $`T_n`$ on the first coordinate. Since each path has finite energy, continuity of the restriction maps (see ) guarantees that
$$\underset{n\mathrm{}}{lim}\rho [A_n,\mathrm{\Phi }_n]|_{\{t\}\times Y}=\underset{n\mathrm{}}{lim}\rho \tau _n^{}[A_n,\mathrm{\Phi }_n]=\rho [B,\mathrm{\Psi }].$$
Gluing gives this space more structure.
###### Proposition 6.5.
Gluing endows $`\overline{}_{X_0}(C)`$ with the structure of a manifold. The space of ideal solutions
$$\widehat{}(C,𝒥)\times _𝒥_{X_0}(𝒥)$$
has the structure of a smooth submanifold of codimension two. In particular, a fundamental class for $`_{X_0}(C)`$ gives rise to a unique fundamental class for $`\overline{}_{X_0}(C)`$.
Proof. Gluing describes the end of $`(X,C)`$ as a fibered product
$$\left(\widehat{}^0(C,𝒥)\times _𝒥_{X_0}^0(𝒥)\times (0,\mathrm{})\right)/S^1,$$
where the superscript denotes based versions of the moduli spaces. This gives the space of ideal solutions a disk-bundle neighborhood in $`\overline{}_{X_0}(C)`$.
In light of the above result, we can define the relative Seiberg-Witten invariant $`SW_{(X_0,C)}`$, as follows.
###### Definition 6.6.
The relative Seiberg-Witten invariant
$$SW_{(X_0,C)}:𝔸(X_0)H^{}(C)$$
is defined by
$$SW_{(X_0,C)}(a\omega )=[\overline{}_{X_0}(C)],\overline{i}^{}(\mu (a))\overline{\rho }_C^{}(\omega ).$$
We now spell out the strategy for proving Proposition 3.8. First, it is shown that for $`b𝔸(Y)`$, $`SW_{(X_0,C)}(ab\omega )`$ can be expressed in terms of $`SW_{(X_0,C)}(ab\omega )`$ and $`SW_{(X_0,𝒥)}`$ (Lemma 6.7 and Proposition 6.9). Here, $`b\omega `$ denotes the action of $`𝔸(Y)`$ on $`H^{}(C)`$ induced from the inclusion of $`C`$ in $`^{}(Y)`$. (Note the cohomology classes over $`\overline{}_{X_0}(C)`$ induced from $`𝔸(Y)`$ through the action on $`H^{}(C)`$, and pulled back via $`\overline{\rho }`$, correspond to divisor representatives over $`X_0`$ which are supported “at infinity.”) Then, it is shown that $`SW_{(X_0,C)}(ab\omega )`$ vanishes, when $`b`$ has sufficiently high degree. This follows from algebraic considerations, according to which $`b\omega =b^{}\omega `$, where $`b^{}𝔸(Y)`$ lies in the ideal generated by the cycles in $`Y`$ which bound in $`X_0`$ (Proposition 6.12). It is then easy to see that $`SW(ab^{}\omega )`$ vanishes (Corollary 6.13).
Now, we express the “commutator” $`SW(ab\omega )SW(ab\omega )`$. First note that if $`b`$ is induced from $`H_1(Y)`$, the commutator vanishes, as follows.
###### Lemma 6.7.
Let $`[\gamma ]H_1(Y)`$, then for all $`a𝔸(X)`$ and $`\omega H^{}(C)`$,
$$SW_{(X_0,C)}(a\mu [\gamma ]\omega )=SW_{(X_0,C)}(a\mu [\gamma ]\omega ).$$
Proof. We must show that $`\overline{\rho }_C^{}(\mu [\gamma ])`$ is homologous to $`\overline{i}^{}(\mu [\gamma ])`$. It suffices to verify this over the subset $`(X_0,C)\overline{}_{X_0}(C)`$, since the complement has codimension two, and the classes in question are one-dimensional. Over the subset, now, the claim is easy to verify. On $`_{X_0}(C)`$, $`\rho ^{}(\mu [\gamma ])`$ is represented by $`(\mathrm{Hol}_\gamma \rho _C)^{}(d\theta )`$, the holonomy around a representative of $`\gamma `$ “at infinity” (see Proposition 9.1); while $`i^{}(\mu [\gamma ])`$ is represented by $`\mathrm{Hol}_{\gamma _0}^{}(d\theta )`$, where $`\gamma _0=0\times \gamma 0\times Y[0,\mathrm{})\times YX`$. Now, the cylinder $`[0,\mathrm{})\times \gamma `$ provides a homotopy between $`\mathrm{Hol}_\gamma \rho _C`$ and $`\mathrm{Hol}_{\gamma _0}`$.
It remains to see how the point class commutes. For this class, we can express the commutator in terms of $`SW_{(X,𝒥)}`$ and another Seiberg-Witten invariant, defined below.
###### Definition 6.8.
There is a Seiberg-Witten invariant of the tube
$$\widehat{S}W_{(C,𝒥)}:H^{}(C)H^{}(𝒥)𝔸(X_0)$$
which raises degree by dimension $`2gdim\widehat{}(C,𝒥)=2gdimC`$, defined by
$$\widehat{S}W_{(C,𝒥)}(\omega )=(P_2)_{}\left((\rho _𝒥\times \mathrm{Id})^{}\mathrm{PD}[\mathrm{\Delta }](\rho _CP_1)^{}\omega \right),$$
where $`P_1`$ and $`P_2`$ are the projection maps
$`P_1:\widehat{}(C,𝒥)\times 𝒥\widehat{}(C,𝒥)`$ and $`P_2:\widehat{}(C,𝒥)\times 𝒥𝒥,`$
and $`\mathrm{PD}[\mathrm{\Delta }]`$ denotes the Poincaré dual of the diagonal $`\mathrm{\Delta }𝒥\times 𝒥`$. Thus, $`\widehat{S}W_{(C,𝒥)}`$ satisfies:
$$\widehat{}(C,𝒥)\times _𝒥_{X_0}(𝒥),\overline{i}^{}(\mu (a))\overline{\rho }_C^{}(\omega )=SW_{(X_0,𝒥)}(a\widehat{S}W_{(C,𝒥)}(\omega )).$$
We can now calculate the commutator, which involves comparing the cohomology classes $`\overline{\rho }_C^{}\mu (y)`$ and $`\overline{i}^{}\mu (x)`$ over $`\overline{}_{X_0}(C)`$, where $`y`$ is a point in $`Y`$ and $`x`$ is a point in $`X_0`$.
###### Proposition 6.9.
Choose points $`xX_0`$ and $`yY`$. In $`\overline{}_{X_0}(C)`$, we have
$$\overline{\rho }_C^{}(\mu (y))\overline{i}^{}(\mu (x))=\mathrm{PD}[\widehat{}(C,𝒥)\times _𝒥_{X_0}(𝒥)].$$
Consequently, there is a relation between Seiberg-Witten invariants:
$$SW_{(X_0,C)}(a\mu (y)\omega )SW_{(X_0,C)}(a\mu (x)\omega )=SW_{(X_0,𝒥)}(a\widehat{S}W_{(C,𝒥)}(\omega )).$$
Proof. Clearly, the difference $`\overline{\rho }_C^{}(\mu (y))\overline{i}^{}(\mu (x))`$ is the first Chern class of the circle bundle $`\mathrm{Hom}_{S^1}(_x,_y)`$. Here, $`_z`$ denotes the moduli space based at $`z`$; see Section 9. To prove the proposition, we must verify that this bundle admits a section $`\sigma `$ in the complement of
$$\widehat{}(C,𝒥)\times _𝒥_{X_0}(𝒥)\overline{}_{X_0}(C)$$
(i.e. over $`_{X_0}(C)\overline{}_{X_0}(C)`$) and that, with respect to a trivialization of the circle bundle over a disk transverse to the submanifold, the restriction of the section to the boundary induces a map from the circle to the circle which has degree one.
The section $`\sigma `$ is induced by parallel transport, as follows. Let $`\gamma `$ be a half-infinite arc formed by joining $`[0,\mathrm{})\times y`$ to any arc which connects $`x`$ to $`0\times y`$. Over the point $`[A,\mathrm{\Phi }]_{X_0}(C)`$, parallel transport via $`A`$ along $`\gamma `$ induces a homomorphism in $`\mathrm{Hom}_{S^1}(_x,_y)`$.
We now verify that the trivialization induces a degree one map around circles transverse to the submanifold. For any point in the submanifold
$$[A_1,\mathrm{\Phi }_1]\times [A_2,\mathrm{\Phi }_2]\widehat{}(C,𝒥)\times _𝒥_{X_0}(𝒥),$$
fix fibers
$`[A_1,\mathrm{\Phi }_1,\lambda _1]_x|_{[A_1,\mathrm{\Phi }_1]}`$ and $`[A_2,\mathrm{\Phi }_2,\lambda _2]_y|_{[A_2,\mathrm{\Phi }_2]}.`$
These choices induce a trivialization of $`\mathrm{Hom}_{S^1}(_x,_y)`$ over a disk in $`\overline{}_{X_0}(C)`$ transverse to $`[A_1,\mathrm{\Phi }_1]\times [A_2,\mathrm{\Phi }_2]`$ (obtained by varying the gluing and translation parameters). Calculating the desired degree amounts to seeing how the holonomy along $`\gamma `$ varies as the gluing parameter is rotated. But holonomy along any path which crosses the gluing region once varies as a degree one function of the gluing parameter.
We can understand the action of $`𝔸(Y)`$ on $`H^{}(C)`$ explicitly, under the identification $`C\mathrm{Sym}^k(\mathrm{\Sigma })`$.
Before describing this, we begin with a few preliminaries about the homology of symmetric products of $`\mathrm{\Sigma }`$ (for an extensive discussion of this topic, see ). Recall that $`\mathrm{Sym}^k(\mathrm{\Sigma })`$ can be viewed as the quotient of the $`k`$-fold Cartesian product $`\mathrm{\Sigma }^{\times k}`$ by the action of the symmetric group on $`k`$ letters. We denote the quotient map by
$$q:\mathrm{\Sigma }^{\times k}\mathrm{Sym}^k(\mathrm{\Sigma }).$$
According to elementary properties of the transfer homomorphism,
$$q_{}:H_{}(\mathrm{\Sigma }^{\times k})H_{}(\mathrm{Sym}^k(\mathrm{\Sigma }))$$
is surjective. Dually, we have a map
$$q^{}:H^{}(\mathrm{Sym}^k(\mathrm{\Sigma }))H^{}(\mathrm{\Sigma }^{\times k})$$
which identifies $`H^{}(\mathrm{Sym}^k(\mathrm{\Sigma }))`$ with the elements of $`H^{}(\mathrm{\Sigma }^{\times k})H^{}(\mathrm{\Sigma })^k`$ which are invariant under the symmetric group action. In particular, by summing over the action, we obtain a map
$$\mathrm{Sym}^k:H^{}(\mathrm{\Sigma })H^{}(\mathrm{Sym}^k(\mathrm{\Sigma })).$$
Thus, if we fix any collection of points $`\{p_2,\mathrm{},p_k\}\mathrm{\Sigma }`$, given $`\omega H^{}(\mathrm{\Sigma })`$, $`\mathrm{Sym}^k(\mathrm{\Sigma })`$ is the class characterized by the property that
$$\mathrm{Sym}^k(\omega ),q_{}(Z\times p_2\times \mathrm{}\times p_k)=\omega ,Z,$$
for any cycle $`Z\mathrm{\Sigma }`$. Equivalently, given a cycle $`ZH_{}(\mathrm{\Sigma })`$, $`\mathrm{Sym}^k(\mathrm{PD}[Z])`$ is Poincaré dual to the cycle $`q(Z\times \mathrm{\Sigma }\times \mathrm{}\times \mathrm{\Sigma })`$. The above discussion works over rational coefficients (which suffices for our purposes), but in fact it works over $``$ as well, since $`H_{}(\mathrm{Sym}^k(\mathrm{\Sigma }))`$ has no torsion (see ).
###### Proposition 6.10.
Under the identification $`C\mathrm{Sym}^k(\mathrm{\Sigma })`$, the canonical map
$$𝔸(Y)H^{}(\mathrm{Sym}^k(\mathrm{\Sigma })),$$
induced from the inclusion of $`\mathrm{Sym}^k(\mathrm{\Sigma })=C^{}(Y)`$, takes $`\mu (y)`$ for $`yH_{}(Y)`$ to the cohomology class $`\mathrm{Sym}^k(\mathrm{PD}[\pi _{}(y)])H^2(\mathrm{Sym}^k(\mathrm{\Sigma }))`$, where $`\pi :Y\mathrm{\Sigma }`$ is the projection map.
Proof. We can reduce to a corresponding statement for configurations over $`\mathrm{\Sigma }`$, as follows. Let $`E`$ be a line bundle over $`\mathrm{\Sigma }`$, so that $`W\pi ^{}(E(K_\mathrm{\Sigma }))`$. Then, pull-back induces a map
$$\pi ^{}:(\mathrm{\Sigma },E)=𝒜(E)\times \mathrm{\Gamma }(E)/\mathrm{Map}(\mathrm{\Sigma },S^1)(Y,W),$$
to the configurations where the fiber-wise holonomy of the connection is constant, and the section is covariantly constant around each fiber. The identification between the critical manifolds and the symmetric powers $`C\mathrm{Sym}^k(\mathrm{\Sigma })`$ described in is obtained by proving that $`C`$ lies in the image of this pull-back map, and indeed that it lies in the pull-back of the vortex moduli space, which, according to (see also ), is in turn identified with the space of divisors, by looking at the zero-set of the section. The key points we need presently are that $`C`$ lies in $`\pi ^{}((\mathrm{\Sigma },E))`$, and that configurations are the pull-backs of configurations $`[A,\mathrm{\Phi }](\mathrm{\Sigma },E)`$, where $`\mathrm{\Phi }`$ is $`\overline{}_A`$-holomorphic section.
Over $`(\mathrm{\Sigma },E)`$, there is a universal line bundle $`(\mathrm{\Sigma })`$, defined in the usual manner. Note that
$$(Y)|_{\pi ^{}((\mathrm{\Sigma }))\times Y}\pi ^{}((\mathrm{\Sigma })),$$
so $`\mu (y)|_{\pi ^{}((\mathrm{\Sigma }))}`$ for $`yH_{}(Y)`$ agrees with $`\mu ([\pi _{}(y)])`$, where the former $`\mu `$-map is induced from $`(Y)`$, and the latter from $`(\mathrm{\Sigma })`$. We have thus reduced the proof of the proposition to a statement purely over $`\mathrm{\Sigma }`$; so for the duration of the proof, $``$ will refer to $`(\mathrm{\Sigma })`$, $``$ will refer to $`(\mathrm{\Sigma },E)`$, and all $`\mu `$-maps will be calculated over $`\mathrm{\Sigma }`$.
To facilitate the proof over $`\mathrm{\Sigma }`$, we pause for a discussion about the canonical section $`\sigma `$ of the universal line bundle $``$, which takes the configuration $`[A,\mathrm{\Phi }]\times \{x\}\times \mathrm{\Sigma }`$ to the based configuration $`[A,\mathrm{\Phi },\mathrm{\Phi }(x)]`$. This section has the property that, under the canonical identification of $`|_{[A,\mathrm{\Phi }]\times \mathrm{\Sigma }}E`$ (where $`E`$ is the bundle over $`\mathrm{\Sigma }`$ with Chern number $`k`$), the restriction $`\sigma |_{[A,\mathrm{\Phi }]\times \mathrm{\Sigma }}`$ is identified with the section $`\mathrm{\Phi }`$ of $`E`$. In particular, if $`\mathrm{\Phi }0`$ is a holomorphic section, then $`\sigma |_{[A,\mathrm{\Phi }]\times \mathrm{\Sigma }}`$ has at most $`k`$ zeros; moreover, if it has $`k`$ zeros, then each is transverse.
Now, if $`y\mathrm{\Sigma }`$ is a point (i.e. a generator of $`H_0(\mathrm{\Sigma };)`$), then by definition, $`\mu (y)`$ is the element of $`H^2(C)`$ whose pairing against any homology class $`[S]H_2(C)`$ is given by
$$\mu (y),[S]=c_1(_y),[S].$$
where, as usual, $`_y`$ denotes the restriction of $``$ to $`(\mathrm{\Sigma },E)\times \{y\}`$. Choose points $`\{p_2,\mathrm{},p_k\}\mathrm{\Sigma }`$ which are distinct from $`y`$. Recall that $`H_2(\mathrm{Sym}^k(\mathrm{\Sigma }))`$ is generated by the surface $`q(\mathrm{\Sigma }\times p_2\times \mathrm{}\times p_k)`$ (where $`p_2,\mathrm{},p_k`$ are points on $`\mathrm{\Sigma }`$), and the tori of the form $`q(C_1\times C_2\times p_3\times \mathrm{}\times p_k)`$, where $`C_1,C_2\mathrm{\Sigma }`$ are closed curves in $`\mathrm{\Sigma }`$, which we can choose to miss $`y`$. The canonical section $`\sigma `$ restricted to a torus of the form $`q(C_1\times C_2\times p_3\times \mathrm{}\times p_k)\times \{y\}`$ clearly vanishes nowhere (as all $`k`$ of the zeros have been constrained to lie in the set $`C_1C_2\{p_3,\mathrm{},p_k\}`$ which does not include the point $`y`$); thus,
$$c_1(_y),[q(C_1\times C_2\times \mathrm{}\times p_k)]=0.$$
Over $`q(\mathrm{\Sigma }\times p_2,\mathrm{}\times p_k)`$ the canonical section vanishes at the single point $`q(y\times p_2\times \mathrm{}\times p_k)`$. We verify transversality of this zero, as follows. View $`\sigma `$ as a section over $`\mathrm{\Sigma }\times \mathrm{\Sigma }=q(\mathrm{\Sigma }\times p_2\times \mathrm{}\times p_k)\times \mathrm{\Sigma }`$; we know that $`\sigma (y,y)0`$, and that $`D\sigma _{(y,y)}`$ induces an isomorphism from $`0T_y\mathrm{\Sigma }`$ to $`E_y`$ (i.e. that the zero of $`\sigma |_{\{y\}\times \mathrm{\Sigma }}`$ at $`y`$ is transverse). Differentiating the equation that $`\sigma (y,y)0`$, we see that
$$D\sigma _{(y,y)}(0,v)=D\sigma _{(y,y)}(v,0).$$
Thus, the section $`\sigma |_{\mathrm{\Sigma }\times \{y\}}`$ of $`_y|_{[\mathrm{\Sigma }]}`$ has a single, transverse zero, which shows that
$$c_1(),[q(\mathrm{\Sigma }\times p_2\times \mathrm{}\times p_k)]=\pm 1.$$
Moreover, the sign is positive since the section is holomorphic.
Hence, we have proved the result when $`yH_0(\mathrm{\Sigma })`$. Proving the result for classes coming from $`H_1(\mathrm{\Sigma })`$ amounts to proving that, if $`C_1`$ and $`C_2`$ closed curves in $`\mathrm{\Sigma }`$ which meet transversally, then
$$c_1(),q(C_1\times p_2\times \mathrm{}\times p_n)\times C_2=\mathrm{\#}C_1C_2=\mathrm{\#}C_2C_1.$$
Note first that the zeros of the canonical section $`\sigma `$, restricted to $`q(C_1\times p_2\times \mathrm{}\times p_k)\times C_2`$ are the points $`C_1C_2`$ (a zero of $`\sigma `$ corresponds to a point where the section $`\mathrm{\Phi }`$ vanishes at some point of $`C_2`$, but the zeros of $`\mathrm{\Phi }`$ lie in $`C_1\{p_2,\mathrm{},p_k\}`$, and $`\{p_2,\mathrm{},p_k\}C_2`$ is empty). We must now consider the local contribution of each zero (and check transversality).
Consider the map $`C_1\times C_2:S^1\times S^1\mathrm{Sym}^k(\mathrm{\Sigma })\times \mathrm{\Sigma }`$ defined by $`C_1\times C_2(s,t)=q(C_1(s)\times p_2\times \mathrm{}\times p_k)\times C_2(t)`$. Suppose for notational simplicity that $`C_1(0)=C_2(0)=y`$. We can view $`\sigma `$ as a section of $``$ pulled back to this torus. Now, evaluated on a typical tangent vector to the torus $`a\frac{}{s}+b\frac{}{t}`$, the derivative of $`\sigma `$ at the intersection point is given by
(20) $`D_{(y,y)}\sigma (C_1\times C_2)(a{\displaystyle \frac{}{s}}+b{\displaystyle \frac{}{t}})`$ $`=`$ $`aD_{(y,y)}\sigma ({\displaystyle \frac{dC_1}{ds}}(0),0)+bD_{(y,y)}\sigma (0,{\displaystyle \frac{dC_2}{dt}}(0))`$
$`=`$ $`aD_{(y,y)}\sigma (0,{\displaystyle \frac{dC_1}{ds}}(0))+bD_{(y,y)}\sigma (0,{\displaystyle \frac{dC_2}{dt}}(0)).`$
(We have used the chain rule and the derivative of the relation that $`\sigma (C_1(s),C_1(s))0`$.) Transversality of the intersection of $`C_1`$ and $`C_2`$ at $`0`$ ensures that the image of this differential is
$$D_{(y,y)}\sigma (0T\mathrm{\Sigma }_y);$$
so transversality of the section corresponding to $`q(y\times p_2\times \mathrm{}\times p_k)`$ at its zero $`y`$ ensures that the image of the differential surjective onto the fiber of $`E`$ over $`y`$; i.e. the canonical section is transverse. The sign is correct, as one can see by inspecting Equation (20).
###### Remark 6.11.
With the help of the above results, we can describe explicitly the invariant of the tube:
$$\widehat{S}W_{(C,𝒥)}:H^{}(C)H^{}(𝒥)𝔸(X_0),$$
which we do now for completeness. Let $`\mathrm{\Lambda }=\mathrm{\Lambda }^{}H_1(\mathrm{\Sigma })𝔸(Y)`$. According to the proof of Lemma 6.7, $`\widehat{S}W_{(C,𝒥)}`$ is a homomorphism of $`\mathrm{\Lambda }`$-modules; so, since $`𝔸(Y)=\mathrm{\Lambda }[U]`$ surjects onto $`H^{}(C)`$, the invariant is determined by $`\widehat{S}W_{(C,𝒥)}(U^i)`$, as $`i`$ ranges over the non-negative integers. Since the Poincaré dual of $`\mathrm{Sym}^k(\mathrm{\Sigma })T^{2g}`$ (which is the image of $`\widehat{}(C,𝒥)`$ under $`\rho _𝒥`$, according to Theorem 4.1), is
$$\frac{(_{i=1}^g\mu (A_i)\mu (B_i))^k}{k!},$$
it follows that
$$\widehat{S}W_{(C,𝒥)}(U^{\mathrm{}})=\frac{(_{i=1}^gA_iB_i)^{k+\mathrm{}}}{(k+\mathrm{})!}.$$
We will not use this formula, however. The results we prove in this paper require only the general properties of $`\widehat{S}W_{(C,𝒥)}`$ which follow from its definition, together with Proposition 6.9.
According to Lemma 6.7, if $`\gamma Y`$ is a curve which is null-homologous in $`X_0`$, then it annihilates the relative invariants, in the sense that
$$SW_{(X_0,C)}(a\mu (\gamma )\omega )=0.$$
If sufficiently many curves in $`Y`$ become null-homologous in $`X_0`$, then any class of sufficiently high degree in $`𝔸(Y)`$ annihilates the invariant, as follows.
###### Proposition 6.12.
Fix natural numbers $`k,\mathrm{}`$ with $`\mathrm{}k`$. Let $`I`$ denote the ideal generated by $`\mu (A_1),\mu (A_2),\mathrm{},\mu (A_{\mathrm{}})`$ in $`H^{}(\mathrm{Sym}^k(\mathrm{\Sigma }))`$. Then every element of $`H^{}(\mathrm{Sym}^k(\mathrm{\Sigma }))`$ of degree greater than $`k`$ lies in $`I`$.
Proof. The vector space $`H^{}(\mathrm{Sym}^k(\mathrm{\Sigma }))`$ is generated by homogeneous elements of the form
$$U^a\underset{q=1}{\overset{b}{}}(A_{i_q}B_{i_q})\underset{r=1}{\overset{c}{}}A_{i_{b+r}}\underset{s=1}{\overset{d}{}}B_{i_{b+c+s}},$$
where $`\{i_1,\mathrm{},i_{b+c+d}\}`$ is a subset of $`\{1,\mathrm{},g\}`$, and $`a`$, $`b`$, $`c`$, $`d`$ are integers with $`a+b+c+dk`$.
Clearly, it suffices to prove the proposition for homogeneous generators of degree $`k+1`$. Modulo $`I`$, such an element is equivalent to the element
$$\underset{p=1}{\overset{a}{}}(UA_pB_p)\underset{q=1}{\overset{b}{}}(A_{i_q}B_{i_q}A_{a+q}B_{a+q})\underset{r=1}{\overset{c}{}}A_{i_{b+r}}\underset{s=1}{\overset{d}{}}B_{i_{b+c+s}}.$$
Indeed, in light of the fact that $`a+bkd\mathrm{}d`$, we can arrange (after possibly simultaneously permuting the indices of the $`\{A_i\}_{i=1}^g`$ and $`\{B_i\}_{i=1}^g`$) that for each $`s=1,\mathrm{}d`$, $`a+b<i_{b+c+s}`$. Moreover, the original homogeneous element would automatically lie in $`I`$ unless we had that $`a+b<k\mathrm{}<i_{b+r}`$ for all $`r=1,\mathrm{},c`$. Put together, must consider elements of the above form which satisfy the constraint that $`a+b<i_j`$ for all $`j>b`$. If the degree of such an element is $`k+1`$, it must vanish in $`H^{}(\mathrm{Sym}^k(\mathrm{\Sigma }))`$.
This vanishing can be seen geometrically: $`U`$ is Poincaré dual to the subset (identified with $`\mathrm{Sym}^{k1}(\mathrm{\Sigma })`$) of $`\mathrm{Sym}^k(\mathrm{\Sigma })`$ where one point is constrained to lie in a specified point on $`\mathrm{\Sigma }`$: $`A_i`$ (resp. $`B_i`$) is Poincaré dual to the cycle where one point is constrained to lie on $`A_i`$ (resp. $`B_i`$). Thus, (if one chooses the point representing $`U`$ to be $`A_iB_i`$), then $`UA_iB_i`$ is Poincaré dual to the locus where two distinct points are constrained; one is to lie on $`A_i`$, the other on $`B_i`$. Similarly, the manifold Poincaré dual to $`A_iB_iA_jB_j`$ gives a constraint on two distinct points in the symmetric power. Finally, the remaining $`A_{i_{b+r}}`$ and $`B_{i_{b+c+s}}`$ give additional, disjoint constraints (these are disjoint, if one chooses that representing curves to be disjoint from the $`A_i`$ and $`B_i`$ for $`i=1,\mathrm{},a+b`$, which can be arranged since $`a+b<i_{b+r}`$ for all $`r1`$). Thus, since the total degree of the expression considered is $`k+1`$, we have put constraints on $`k+1`$ distinct points, forcing the intersection to be empty.
###### Corollary 6.13.
Suppose that $`C=\mathrm{Sym}^k(\mathrm{\Sigma })`$, and let $`\mathrm{}k`$ be an integer so that there is a symplectic basis $`\{A_i,B_i\}_{i=1}^g`$ for $`H_1(\mathrm{\Sigma })`$ so that $`i_{}(A_i)=0`$ in $`H_1(X_0;)`$ for $`i=1,\mathrm{},\mathrm{}`$. Then, for each $`b𝔸(\mathrm{\Sigma })`$ of degree $`d(b)>k`$, and each $`a𝔸(X_0)`$, $`\omega H^{}(C)`$, we have
$$SW_{(X_0,C)}(ab\omega )0.$$
Proof. By Proposition 6.12, $`b`$ lies in the ideal generated by $`\mu (A_1),\mathrm{},\mu (A_{\mathrm{}})`$. Now the proposition follows from Lemma 6.7.
We now have the promised proof of Proposition 3.8.
Proof of Proposition 3.8. Recall that we have constructed $`SW_{(X_0,C)}`$ so that
$$SW_𝔰^{irr}(ai_{}(b))=SW_{(X_0,C)}(ai_{}(b)1).$$
By Lemma 6.7 and Proposition 6.9, we can write
$$SW_{(X_0,C)}(ai_{}(b)1)=SW_{(X_0,C)}(ab)+SW_{(X_0,𝒥)}(ac).$$
for some $`c𝔸(\mathrm{\Sigma })`$. Note that $`c`$ lies in the ideal generated by $`H_1(\mathrm{\Sigma })`$, as it can be expressed in terms of Seiberg-Witten invariants of the tube, which take values in $`H^{}(𝒥)\mathrm{\Lambda }^{}(H_1(\mathrm{\Sigma }))H_1(\mathrm{\Sigma })𝔸(\mathrm{\Sigma })`$. By Corollary 6.13, the first term vanishes (using the homological hypothesis of the inclusion of $`\mathrm{\Sigma }`$ in $`X`$). The remaining term is identified with an absolute invariant, according to Proposition 5.9. $`\mathrm{}`$
The proof of Proposition 3.10, follows from the same argument as Proposition 3.8; only in that case, one must use the following (much simpler) analogue of Corollary 6.13.
###### Lemma 6.14.
Suppose that $`C=\mathrm{Sym}^k(\mathrm{\Sigma })`$. Then, for each $`b𝔸(\mathrm{\Sigma })`$ of degree $`d(b)>2k`$, and each $`a𝔸(X_0)`$, $`\omega H^{}(C)`$, we have
$$SW_{(X_0,C)}(ab\omega )0.$$
Proof. This follows immediately from the fact that $`dimC=2k`$.
## 7. The moduli spaces over $`N`$
The purpose of this section is to give the results about the neighborhood of $`\mathrm{\Sigma }`$ which were used in Section 5. Most of these results are applications of and . We assume for the duration of this section that the $`\mathrm{Spin}_{}`$ structure over $`N`$ satisfies $`c_1(𝔰),[\mathrm{\Sigma }]n(mod2n)`$. We return to the excluded cases in Section 8.
Over $`N`$, endowed with a cylindrical-end metric and a certain torsion connection on $`TN`$ , the Seiberg-Witten equations admit a complex interpretation analogous to the complex interpretation of the equations over a Kähler manifold (see Section 5 of for an explicit description of this connection, and especially Proposition 5.6 where the complex interpretation is proved). The Seiberg-Witten equations over $`N`$ can be written as equations for a connection $`A`$ over $`E`$, $`\alpha \beta (\mathrm{\Omega }^{0,0}\mathrm{\Omega }^{0,1})(N,E)`$:
(21) $`2\mathrm{\Lambda }F_A\mathrm{\Lambda }F_{K_N}`$ $`=`$ $`{\displaystyle \frac{i}{2}}(|\alpha |^2|\beta |^2)`$
(22) $`\mathrm{Tr}F_A^{0,2}`$ $`=`$ $`\overline{\alpha }\beta `$
(23) $`\overline{}_A\alpha +\overline{}_A^{}\beta `$ $`=`$ $`0,`$
where $`\mathrm{\Lambda }`$ denotes projection onto the $`(1,1)`$ form of the metric. As noted in , for finite energy solutions, decay estimates justify the usual integration-by-parts which shows that one of $`\alpha `$ or $`\beta `$ must vanish identically; i.e. the solutions over $`N`$ correspond to vortices over $`N`$.
When $`\beta 0`$, then $`A`$ induces an integrable $`\overline{}`$-operator on $`E`$, $`\overline{}_A`$, with respect to which $`\alpha `$ is holomorphic. Moreover, by the usual exponential decay results, together with the understanding of the solutions over $`Y`$ (Theorem 4.1), $`(A,\alpha )`$ exponentially approaches the pull-back of a vortex solution over $`\mathrm{\Sigma }`$. According to , the underlying holomorphic data extends to the ruled surface $`R`$ obtained by attaching a copy of $`\mathrm{\Sigma }`$ (denoted $`\mathrm{\Sigma }_+`$) to $`N`$ “at infinity.” We state the results here for convenience.
###### Definition 7.1.
Let $`\mathrm{\Phi }\mathrm{\Gamma }(N,W^+)`$, $`\mathrm{\Psi }\mathrm{\Gamma }(Y,W)`$ be a pair of spinors, and $`\delta >0`$ be some real number. Then, $`\mathrm{\Psi }`$ is said to $`\delta `$-decay to $`\mathrm{\Psi }`$ if for each $`k0`$,
$$\underset{t\mathrm{}}{lim}\underset{\{t\}\times Y}{sup}e^{\delta t}|^{(k)}\mathrm{\Psi }^{(k)}\pi ^{}(\mathrm{\Psi })|=0,$$
where $`^{(k)}`$ denotes the $`k`$-fold covariant derivative. More generally, $`\mathrm{\Phi }`$ is said to decay to $`\mathrm{\Psi }`$ if there is some $`\delta >0`$ so that $`\mathrm{\Psi }`$ $`\delta `$-decays to $`\mathrm{\Psi }`$. A similar notion can be defined for objects other than spinors, such as connections, differential forms, etc.
###### Definition 7.2.
Given a line bundle $`E`$ over $`Z`$, a holomorphic pair $`(A,\alpha )`$ in $`E`$ is a pair consisting of a $`\overline{}`$-operator $`\overline{}_A`$ over $`E`$, and a section $`\alpha `$ of $`E`$, so that $`F_A^{0,2}=0`$, and $`\overline{}_A\alpha =0`$.
###### Theorem 7.3.
Let $`(A,\alpha )`$ be a holomorphic pair on $`N`$ which decays to a the pull-back of a holomorphic pair $`(A_0,\alpha _0)`$ over $`\mathrm{\Sigma }`$. Then, there is a naturally associated line bundle $`\widehat{E}`$ over $`R`$ and holomorphic pair $`(\widehat{A},\widehat{\alpha })`$ in $`\widehat{E}`$, so that $`(\overline{}_{\widehat{A}},\widehat{\alpha })|_{R\mathrm{\Sigma }_+}(\overline{}_A,\widehat{\alpha })`$ and $`(\overline{}_{\widehat{A}},\widehat{\alpha })|_{\mathrm{\Sigma }_+}(\overline{}_{A_0},\alpha _0)`$.
The above theorem is essentially a restatement of Theorems 7.7 of , where it is stated for the cylinder, thought of as $`R`$ minus two copies of $`\mathrm{\Sigma }`$, rather than the neighborhood of $`\mathrm{\Sigma }`$, thought of as $`R`$ minus one copy of $`\mathrm{\Sigma }`$ (though the proof is no different). Analogous results for the anti-self-dual equations were obtained by Guo .
In a similar vein we have the following result, which allows us to deal with solutions with reducible boundary values. We state the result slightly differently from the above, since we will apply it in other contexts later.
###### Theorem 7.4.
Let $`A`$ be a connection on a line bundle $`E`$ over $`N`$, with $`F_A^{0,2}=0`$ and $`E|_{(0,\mathrm{})\times Y}\pi ^{}(E_0)`$, so curvature form $`F_A`$ decays to the pull-back of a closed two-form $`F_0`$ over $`\mathrm{\Sigma }`$ with
$$\left(\frac{i}{2\pi }_\mathrm{\Sigma }F_0\right)c_1(E_0),[\mathrm{\Sigma }]n.$$
Then, there is an associated line bundle $`\widehat{E}`$ over $`R`$ and complex structure $`\overline{}_{\widehat{A}}`$ with
$$(\widehat{E},\overline{}_{\widehat{A}})|_{R\mathrm{\Sigma }_+}(E,\overline{}_A),$$
and $`c_1(\widehat{E}),[\mathrm{\Sigma }_+]`$ is the greatest integer congruent to $`c_1(E),[\mathrm{\Sigma }]`$ moduli $`n`$ smaller than $`\frac{i}{2\pi }F_0`$. Furthermore, there is a natural identification
$$\mathrm{Ker}\overline{)}D_AL^2H^0(R,\widehat{E})H^2(R,\widehat{E}),$$
and
$$\mathrm{Coker}\overline{)}D_AL^2H^1(R,\widehat{E}).$$
The above is proved in Proposition 9.2 (see Corollary 9.11 and Theorem 10.6) of .
These results allow us to rule out the existence of certain solutions. Recall first the following standard fact about the cohomology of $`R`$ (see for example ):
###### Proposition 7.5.
Let $`R`$ denote the ruled surface over $`\mathrm{\Sigma }`$, which is given as the projectivization of $`L`$, $`(L)`$ (here, $`L`$ is some line bundle over $`\mathrm{\Sigma }`$). Let $`\widehat{E}`$ be a line bundle over the ruled surface $`R`$ and let $`E_0`$ denote the restriction of of $`\widehat{E}`$ to $`(0)\mathrm{\Sigma }`$ and let $`\mathrm{}`$ be the evaluation of $`c_1(\widehat{E})`$ on a fiber in the ruling. Then, if $`\mathrm{}0`$,
$$H^0(R,\widehat{E})\underset{j=0}{\overset{\mathrm{}}{}}H^0(\mathrm{\Sigma },E_0L^j);H^1(R,\widehat{E})\underset{j=0}{\overset{\mathrm{}}{}}H^1(\mathrm{\Sigma },E_0L^j);H^2(R,\widehat{E})=0;$$
and if $`\mathrm{}<0`$,
$$H^0(R,\widehat{E})=0;H^1(R,\widehat{E})\underset{j=1}{\overset{\mathrm{}1}{}}H^0(\mathrm{\Sigma },E_0L^j);H^2(R,\widehat{E})\underset{j=1}{\overset{\mathrm{}1}{}}H^1(\mathrm{\Sigma },E_0L^j).$$
In particular, if $`\mathrm{}=1`$, then $`H^{}(R,\widehat{E})=0`$.
We can apply these results to the case where $`N`$ is a neighborhood of a surface of self-intersection number with $`\mathrm{\Sigma }\mathrm{\Sigma }>2g2`$. Recall that according to Theorem 4.1, for each $`\mathrm{Spin}_{}`$ structure $`𝔰`$ on $`N`$, there are at most two components to the moduli space of the boundary, the reducible component $`𝒥`$ and the irreducible component $`C`$.
###### Proposition 7.6.
If
$$\left|c_1(𝔰),[\mathrm{\Sigma }]\right|<n,$$
the moduli space $`_N(𝒥)`$ contains only reducibles. Moreover, the space of reducibles is smoothly identified with the Jacobian torus $`𝒥`$ (i.e. the kernel and the cokernel of the Dirac operator coupled to any reducible vanishes). Furthermore, $`_N(C)`$ is empty.
Proof. We prove that both moduli spaces $`_N^{}(𝒥)`$ and $`_N(C)`$ are empty. Suppose there were some finite energy solution to the Seiberg-Witten equations in a $`\mathrm{Spin}_{}`$ structure with $`|c_1(𝔰),[\mathrm{\Sigma }]|<n`$. We know that the spinor lies entirely in one of the two summands in the splitting of the spinor bundle $`W^+E(K_{N}^{}{}_{}{}^{1}E)`$ (i.e. it is an $`\alpha `$\- or a $`\beta `$-spinor, in the notation of Equations (21)-(23)). By conjugating if necessary (which switches the two summands and sends the $`\mathrm{Spin}_{}`$ structure $`𝔰`$ to another one $`J𝔰`$ with $`c_1(J𝔰)=c_1(𝔰)`$), we can assume without loss of generality that the solution is an $`\alpha `$-solution.
According to Theorem 7.3 (and Theorem 7.4, when the boundary value is reducible), we can extend the data $`(E,\overline{}_A,\alpha )`$ over the associated ruled surface $`R`$, obtained by attaching the curve $`\mathrm{\Sigma }_+`$ at infinity. The fact that $`\widehat{E}`$ is an extension of $`E`$ says that
$`c_1(\widehat{E}),[\mathrm{\Sigma }_{}]`$ $`=`$ $`c_1(E),[\mathrm{\Sigma }]`$
$`=`$ $`{\displaystyle \frac{1}{2}}c_1(𝔰)+c_1(K_N),[\mathrm{\Sigma }]`$
$`=`$ $`g1+{\displaystyle \frac{n+c_1(𝔰),[\mathrm{\Sigma }]}{2}}`$
where $`\mathrm{\Sigma }_{}`$ is the curve in $`\widehat{E}`$ with self-intersection number $`n`$ (which is identified with $`\mathrm{\Sigma }N`$). By our hypothesis, then,
$$g1<c_1(\widehat{E}),[\mathrm{\Sigma }_{}]<n+g1.$$
On the other hand, Equation (21) says that $`\frac{i}{2\pi }F_A`$ converges to the pullback of a form over $`\mathrm{\Sigma }`$ whose integral is $`g1`$, so Theorem 7.4 guarantees that the Chern number of restriction to the other section of the ruling satisfies the bound
$$n+g1<c_1(\widehat{E}),[\mathrm{\Sigma }_+]<g1.$$
Now, since the Poincaré dual of a fiber is $`(\mathrm{PD}[\mathrm{\Sigma }_+]\mathrm{PD}[\mathrm{\Sigma }_{}])/n`$, we see that the evaluation of $`c_1(\widehat{E})`$ on a fiber is given by
$$\mathrm{}=\frac{c_1(\widehat{E}),[\mathrm{\Sigma }_+]c_1(\widehat{E}),[\mathrm{\Sigma }_{}]}{n}=1.$$
According to Proposition 7.5, it follows that $`\widehat{\alpha }`$ (and hence also $`\alpha `$) must vanish identically, contradicting the irreducibility hypothesis on $`(A,\alpha )`$.
The fact that the reducibles are smoothly cut out in this range follows in an analogous manner, using Theorem 7.4 and Proposition 7.5.
###### Remark 7.7.
Most of this result can be found in Proposition 2.5 of .
The above vanishing result is special to the particular $`\mathrm{Spin}_{}`$ structures considered, as it used the fact that the Dolbeault cohomology of certain line bundles over the ruled surface vanish. In general, the moduli spaces over $`N`$ typically do contain irreducibles. To study the deformation theory around these irreducibles, we use an infinitesimal version of Theorem 7.3; but first, we pause for a brief discussion of deformation theory for the Seiberg-Witten equations in general.
In general, on a four-manifold $`X_0`$ with a cylindrical end, the deformation complex around a solution $`(A,\mathrm{\Phi })`$ whose boundary value is smooth and irreducible, is given by
$$\mathrm{\Omega }^0(X_0,i)\mathrm{\Omega }^1(X_0,i)\mathrm{\Gamma }(X_0,W^+)\mathrm{\Omega }^+(X_0,i)\mathrm{\Gamma }(X_0,W^{}).$$
Here, terms in $`\mathrm{\Omega }^0(X_0,i)`$ are required to lie in $`L_{\delta ,k}^2`$, the $`\delta `$-decaying Sobolev space with $`k`$ derivatives (here we can choose any $`k3`$); i.e. functions for which
$$(f_{\delta ,k})^2=_{X_0}(|f|^2+|f|^2+\mathrm{}+|^{(k)}f|^2)e^{\delta \tau }<\mathrm{},$$
where $`\tau `$ is a smooth function on $`X_0`$ which agrees with the $`t`$ coordinate over the cylindrical end. Terms in $`\mathrm{\Omega }^1(X_0,i)\mathrm{\Gamma }(X_0,W^+)`$ are required to lie in $`L_{\delta ,k1}^2`$ extended by the tangent space to the moduli space at infinity at $`\rho [A,\mathrm{\Phi }]`$. Finally, terms in $`\mathrm{\Omega }^+(X_0,i)\mathrm{\Gamma }(X_0,W^{})`$ are required to lie in $`L_{k2}^2`$. The first map in the deformation complex is the linearization of the gauge group action on $`[A,\mathrm{\Phi }]`$ around the identity, while the second is the linearization of the Seiberg Witten equations around $`[A,\mathrm{\Phi }]`$. When the boundary value of $`[A,\mathrm{\Phi }]`$ is a smooth reducible, then the above specifies the deformation theory for the moduli space based at infinity. In either case, the moduli space of solutions about $`[A,\mathrm{\Phi }]`$ is transversally cut out by the Seiberg-Witten equations on $`X_0`$ if the $`H^2`$ of the above complex vanishes. (This discussion is modeled on the theory developed in .)
The space of divisors in a compact, complex surface $`X`$ admits a deformation theory, defined as follows. Consider the pair $`(\overline{}_{\widehat{A}},\widehat{\alpha })`$ where $`\overline{}_{\widehat{A}}`$ is an integrable $`\overline{}`$-operator, and $`\widehat{\alpha }`$ is $`\overline{}_{\widehat{A}}`$-holomorphic; i.e.
$`F_{\widehat{A}}^{0,2}`$ $`=`$ $`0`$
$`\overline{}_{\widehat{A}}\alpha `$ $`=`$ $`0.`$
This has a deformation complex
$$\begin{array}{ccccccc}\mathrm{\Omega }^{0,0}& & \mathrm{\Omega }^{0,1}\mathrm{\Omega }^{0,0}(E)& & \mathrm{\Omega }^{0,2}\mathrm{\Omega }^{0,1}(E)& & \mathrm{\Omega }^{0,2},\end{array}$$
whose cohomology groups are identified with the cohomology groups of the quotient sheaf $`/\widehat{\alpha }`$, obtained from the short exact sequence of sheaves:
$$\begin{array}{ccccccccc}0& & 𝒪_X& \stackrel{\widehat{\alpha }}{}& & & /\widehat{\alpha }& & 0.\end{array}$$
###### Theorem 7.8.
Let $`(A,\alpha )`$ correspond to a solution to the Seiberg-Witten equations over $`N`$, with irreducible boundary values. Then, the cohomology groups of the deformation complex of the Seiberg-Witten deformation complex are naturally isomorphic to the cohomology groups deformation complex of the divisor $`[\overline{}_{\widehat{A}},\widehat{\alpha }]`$ in the line bundle $`\widehat{E}`$ over $`R`$, provided by Theorem 7.3. When $`(A,\alpha )`$ has a reducible boundary value, then $`H^2`$ of the Seiberg-Witten deformation complex is identified with $`H^1(R,/\widehat{\alpha })`$, while the tangent space of the based moduli space is identified with $`H^0(R,/\widehat{\alpha })`$.
Proof. This follows exactly as in Theorem 9.14 (for irreducible boundary values) and Theorem 10.12 (for reducible boundary values) of . The key observation at this point is to note that
$$\mathrm{\Lambda }\overline{}+|\alpha |^2:L_{\delta ,k}^2L_{\delta ,k2}^2$$
is an isomorphism, which allows one to “unroll” parts of the Seiberg-Witten deformation complex to identify it with the deformation theory of divisors in $`N`$. As in (see also ), we can identify $`\mathrm{\Lambda }\overline{}`$ with the operator over the cylindrical end with
$$e^{2\lambda t}\frac{}{t}e^{2\lambda t}\frac{}{t}+\mathrm{\Delta }_Y,$$
where $`\lambda =\frac{\pi n}{\mathrm{Vol}(\mathrm{\Sigma })}`$. According to the theory of , the operator
$$\mathrm{\Lambda }\overline{}:L_{k,\delta }^2L_{k2,\delta }^2$$
is Fredholm for all weights $`0<\delta <4\lambda `$. In particular, it has the same index for all small $`0<\delta `$ as it has on the weight $`\delta =2\lambda `$, where it can be connected via Fredholm operators to the manifestly self-adjoint operator
$$d^_\lambda d:L_{k,\lambda }^2L_{k2,\lambda }^2,$$
where $`d^_\lambda `$ denotes the formal $`\lambda `$-weighted adjoint of $`d`$. It follows from the homotopy invariance of the index that $`\mathrm{\Lambda }\overline{}+|\alpha |^2`$ has index zero on $`L_{k,\delta }^2`$. From the maximum principle, it has no kernel, so it induces an isomorphism as claimed, identifying the deformation theory of the Seiberg-Witten equations with the deformation theory of divisors in $`N`$. Passing to the ruled surface then follows from Corollary 9.4 of .
###### Proposition 7.9.
Let $`N`$ be a disk bundle over a surface $`\mathrm{\Sigma }`$ with
$$\mathrm{\Sigma }\mathrm{\Sigma }=n<22g,$$
endowed with a $`\mathrm{Spin}_{}`$ structure $`𝔰`$ with
$$n<|c_1(𝔰),[\mathrm{\Sigma }]|n+2g2.$$
Let
$$e=\frac{n+2g2|c_1(𝔰),[\mathrm{\Sigma }]|}{2}.$$
Then, the expected dimensions of the moduli spaces over $`N`$ and $`X_0`$ are given by:
(24) $`\text{e-dim}_N(𝒥)`$ $`=`$ $`2e+1`$
(25) $`\text{e-dim}_N(C)`$ $`=`$ $`2e.`$
Moreover, $`_N^{}(𝒥)`$, $`_N(C)`$, are transversally cut out by the Seiberg-Witten equations (in particular, they are manifolds of the expected dimension). Furthermore, the boundary map
$$\rho :_N(C)C$$
is an orientation-preserving diffeomorphism onto its image.
Proof. The deformation theory around a solution $`[A,\alpha ]_N(C)`$ is identified the deformation theory around a corresponding divisor in the line bundle $`\widehat{E}`$ with $`c_1(\widehat{E}),[\mathrm{\Sigma }_\pm ]=e`$; i.e. with divisors in a line bundle which (topologically) pulls back from $`\mathrm{\Sigma }`$. According to Proposition 7.5, all such divisors actually pull back from the base $`\mathrm{\Sigma }`$; and indeed, the deformation theory corresponds to deformation theory of degree $`e`$ divisors in the base $`\mathrm{\Sigma }`$, which is unobstructed. Thus, $`_N(C)`$ is a manifold of real dimension $`2e`$, transversally cut out by the Seiberg-Witten equations.
The above transversality applies to $`_N(𝒥)`$ as well, except that the expected dimension is greater by one, as we saw in Theorem 7.8.
This identification of deformation theories of $`_N(C)`$ proves that $`\rho `$ is an orientation-preserving local diffeomorphism onto its image in $`\mathrm{Sym}^e(\mathrm{\Sigma })C`$. In fact, it is injective, as follows. As we saw, any two solutions with the same boundary values actually vanish over the same disks (with the same multiplicities). By the usual analysis of the vortex equations, any two such solutions must differ by a complex gauge transformation; i.e. a function $`u`$ which satisfies
$$\mathrm{\Lambda }\overline{}u+|\alpha |^2(e^{2u}1)=0,$$
where $`u`$ is a function which decays on the cylinder. By the maximum principle, such a function must vanish identically.
Having analyzed the moduli spaces over neighborhoods of $`\mathrm{\Sigma }`$, we close with a some general results concerning the rest of the moduli spaces of the complement of $`\mathrm{\Sigma }`$.
###### Proposition 7.10.
Let $`X_0`$ be as in Proposition 5.1. Then, letting $`\text{e-dim}(X)=d`$, we have
(26) $`\text{e-dim}_{X_0}(𝒥)`$ $`=`$ $`d+2g2e2`$
(27) $`\text{e-dim}_{X_0}(C)`$ $`=`$ $`d.`$
Moreover, $`_{X_0}(𝒥)`$, and $`_{X_0}(C)`$ are transversally cut out by the Seiberg-Witten equations (in particular, they are manifolds of the expected dimension).
Proof. By a standard excision argument, we have
$$\text{e-dim}_{X_0}(𝒥)+\text{e-dim}_N(𝒥)2g+1=\text{e-dim}_X(𝔰)=d,$$
which calculates $`\text{e-dim}_{X_0}(𝒥)`$, given Proposition 7.9. Similarly, we have
$$\text{e-dim}_{X_0}(C)+\text{e-dim}_N(C)2e=d,$$
which gives us $`\text{e-dim}_{X_0}(C)`$.
Smoothness of $`_{X_0}(𝒥)`$ and $`_{X_0}(C)`$ follows from adapting methods of .
## 8. Perturbations when $`e=g1`$
In our earlier discussion, we had to exclude one $`\mathrm{Spin}_{}`$ structure over $`Y`$. In this section, we introduce a perturbation of the equations which allows us to handle this case. We begin by adapting results of Section 4 to this perturbed equation, and then, we give a discussion which is parallel to that of Section 7. The perturbations used here are analogues of Taubes’ perturbations in the symplectic category , ; see also for a related discussion.
Recall that the Seiberg-Witten equations over $`Y`$ are obtained as the critical points of the Chern-Simons-Dirac functional $`\mathrm{CSD}`$ defined over the configuration space
$$(Y,W)=𝒜(W)\mathrm{\Gamma }(Y,W)/\mathrm{Map}(Y,S^1),$$
where $`𝒜(Y,W)`$ denotes the space of connection in the spinor bundle $`W`$ which are compatible with some fixed connection $``$ on $`TY`$. The functional is defined by
$$\mathrm{CSD}(B,\mathrm{\Psi })=_Y(BB_0)\mathrm{Tr}(F_B+F_{B_0})_Y\mathrm{\Psi },\overline{)}D_B\mathrm{\Psi },$$
where $`B_0𝒜(Y,W)`$ is some reference connection, $`BB_0\mathrm{\Omega }^1(Y;i)`$ denotes the difference $`1`$-form, and $`\mathrm{Tr}`$ denotes the trace of the corresponding connection on $`W`$. Its Euler-Lagrange equations (the three-dimensional Seiberg-Witten equations) are
(28) $`\mathrm{Tr}(F_B)i\tau (\mathrm{\Psi })=0`$
(29) $`\overline{)}D_B\mathrm{\Psi }=0,`$
where
$$\tau :\mathrm{\Gamma }(Y,W)\mathrm{\Omega }^1(Y;)$$
is adjoint to Clifford multiplication, in the sense that for all $`\gamma \mathrm{\Omega }^1(Y;)`$, $`\mathrm{\Psi }\mathrm{\Gamma }(Y,W)`$, we have
(30)
$$\frac{1}{2}i\gamma \mathrm{\Psi },\mathrm{\Psi }_W=\gamma ,\tau (\mathrm{\Psi })_{\mathrm{\Lambda }^1}.$$
Moreover, its upward gradient flow equations are the usual Seiberg-Witten equations on the four-manifold $`\times Y`$.
When $`Y`$ is a circle-bundle over a Riemann surface with Euler number $`n`$ satisfying
$$n>2g2,$$
recall that these equations are inconvenient in the $`\mathrm{Spin}_{}`$ structure when $`e=g1`$ (in the notation of Section 4). We will find it useful to consider a perturbed functional
$$\mathrm{CSD}_u:(Y,𝔱),$$
where $`u`$, given by
$$\mathrm{CSD}_u(B,\mathrm{\Psi })=\mathrm{CSD}(B,\mathrm{\Psi })+u_Yi\eta (\mathrm{Tr}F_B\mathrm{Tr}F_{B_0}),$$
where $`\eta `$ is the connection form for $`Y`$ over $`\mathrm{\Sigma }`$, and the reference connection $`B_0`$ satisfies $`\mathrm{Tr}(F_{B_0})0`$ (i.e. $`B_0𝒥`$). These give rise to perturbed Seiberg-Witten equations of the form
(31) $`\mathrm{Tr}(F_B)i\tau (\mathrm{\Psi })+iu(d\eta )=0`$
(32) $`\overline{)}D_B\mathrm{\Psi }=0,`$
whose moduli space of solutions is denoted $`𝒩_u(Y)`$. The gradient flow equations of the perturbed functional are solutions to the Seiberg-Witten equations on $`\times Y`$, perturbed by the self-dual component of $`iu(d\eta )`$, which can be collected into moduli spaces, denoted $`_u(C_1,C_2)`$, or their unparameterized versions
$$\widehat{}_u(C_1,C_2)=_u(C_1,C_2)/.$$
We have the following analogue of Theorem 4.1.
###### Theorem 8.1.
Let $`Y`$ be a circle-bundle over a Riemann surface with genus $`g>0`$ and Euler number $`n<22g`$. Let $`𝔱`$ be the $`\mathrm{Spin}_{}`$ structure corresponding to $`g1/nH^2(Y;)`$. For all $`u`$ with $`0<u<2`$, the moduli space contains two components, a reducible one $`𝒥`$, identified with the Jacobian torus $`H^1(\mathrm{\Sigma };/)`$, and a smooth irreducible component $`C`$ diffeomorphic to $`\mathrm{Sym}^{g1}(\mathrm{\Sigma })`$. Both of these components are non-degenerate in the sense of Morse-Bott. There is an inequality $`\mathrm{CSD}_u(𝒥)>\mathrm{CSD}_u(C)`$, so the space $`_u(𝒥,C)`$ is empty. The space $`\widehat{}_u(C,𝒥)`$ is smooth of expected dimension $`2g2`$; indeed it is diffeomorphic to $`\mathrm{Sym}^{g1}(\mathrm{\Sigma })`$.
Proof. Most of this is a straightforward adaptation of .
We begin with the identification of the moduli spaces over $`Y`$. As in , the equations over $`Y`$ reduce to vortex equations over $`\mathrm{\Sigma }`$. More specifically, the components of the moduli spaces $`𝒩_Y(𝔱)`$ correspond to line bundles $`E_0`$ over $`\mathrm{\Sigma }`$ with the property that
$$\pi ^{}(E_0K^1E_0)W,$$
the spinor bundle of $`𝔱`$ (here $`K`$ denotes the canonical line bundle over $`\mathrm{\Sigma }`$). The vortex equations are are equations for $`B𝒜(E_0)`$, $`\alpha \beta \mathrm{\Gamma }(\mathrm{\Sigma },E_0K^1E_0)`$, which, in the case at hand, take the form
(33) $`2F_BF_K+iu(d\eta )`$ $`=`$ $`i(|\alpha |^2|\beta |^2)(1)`$
(34) $`\overline{}_B\alpha +\overline{}_B^{}\beta `$ $`=`$ $`0`$
(35) $`\alpha \beta `$ $`=`$ $`0.`$
Thus, one of $`\alpha `$ or $`\beta `$ must vanish. In fact, in our case,
$$\mathrm{deg}E_0g1(modn).$$
In fact, if
$$\mathrm{deg}E_0g1,$$
then the solution space to these equations ($`0<u<2`$) is empty. More specifically, letting $`\mathrm{deg}E_0=g1+n\mathrm{}`$, we see that when $`\beta 0`$, then by integrating Equation (33) over $`\mathrm{\Sigma }`$ against $`i/2\pi `$, we get
$$2(g1+n\mathrm{})(2g2)+u\mathrm{deg}Y=2n\mathrm{}un0,$$
which forces $`\mathrm{}1`$ (since $`u>0`$). Since in this case $`\mathrm{deg}(E_0)>g1`$, $`H^1(\mathrm{\Sigma },E_0)=0`$, so $`\beta `$ must vanish. If, on the other other hand, it is $`\alpha 0`$, then we obtain in the same manner that
$$2n\mathrm{}un0,$$
which forces $`\mathrm{}0`$ (since $`u<2`$). Since $`\alpha `$ represents a class in $`H^0(\mathrm{\Sigma },E_0)`$, it follows that $`\mathrm{}=0`$.
So, all irreducibles correspond to $`\alpha `$-vortices in the line bundle $`E_0`$ with $`\mathrm{deg}E_0=g1`$. The identification of this space of vortices with the symmetric product follows from (see also ).
Non-degeneracy of the irreducible manifold $`C`$ follows exactly as in . To see non-degeneracy of $`𝒥`$, we appeal to results of Section 5.8 of . Consider the Dirac operator on the $`\mathrm{Spin}_{}`$ structure with spinors $`W=E(K^1)`$ with connection induced from a connection $`B𝒜(E)`$ whose curvature pulls up from $`\mathrm{\Sigma }`$. It is shown in Proposition 5.8.4 of that this Dirac operator admits no harmonic spinors unless the holonomy around a fiber circle in $`Y`$ is trivial. In fact this holonomy is trivial when the following integral is congruent to $`g1`$ modulo $`n`$:
$$\frac{i}{4\pi ^2}_YF_B\eta =g1\frac{u\mathrm{deg}(Y)}{2},$$
(we have used here Equation (31)). Since $`0<u<2`$, this holonomy is non-trivial, so the reducibles admit no harmonic spinors, i.e. $`𝒥`$ is smoothly cut out by the equations.
We now perform the Chern-Simons calculations (see the proof of Proposition 5.23 of ). Suppose $`[(B_1,\mathrm{\Psi }_1)]C`$, and $`[(B_0,0)]𝒥`$. Then, we have
$`2\mathrm{deg}B_0\mathrm{deg}K+u\mathrm{deg}(Y)`$ $`=`$ $`0;`$
$`2\mathrm{deg}B_1\mathrm{deg}K`$ $`=`$ $`0,`$
where by $`\mathrm{deg}B`$, we mean the integral $`\frac{i}{4\pi ^2}_YF_B\eta `$, which when $`B`$ is induced from a line bundle over $`\mathrm{\Sigma }`$, agrees with the degree of that line bundle. So,
$`\mathrm{CSD}_u(B_1)`$ $`=`$ $`{\displaystyle _Y}(B_1B_0)(2F_{B_1}+2F_{B_0}2F_K)+u{\displaystyle i\eta }(2F_{B_1}2F_{B_0})`$
$`=`$ $`{\displaystyle \frac{8\pi ^2}{\mathrm{deg}Y}}(\mathrm{deg}B_1\mathrm{deg}B_0)(\mathrm{deg}B_1+\mathrm{deg}B_0\mathrm{deg}K)`$
$`+u{\displaystyle i\eta }(2F_{B_1}2F_{B_0})`$
$`=`$ $`2\pi ^2u^2\mathrm{deg}Y,`$
which is negative; while $`\mathrm{CSD}_u(B_0)=0`$.
The smoothness of the space of flows, and its identification with the symmetric product, follows exactly as in the unperturbed case (see Section 4).
We now turn to the neighborhood of $`\mathrm{\Sigma }`$. We use a perturbation over $`N`$ which is compatible with the above perturbation over $`Y`$. Specifically, let
$$f:N$$
be a smooth function which is identically zero on the complement of the cylinder $`[0,\mathrm{})\times YN`$, and identically one on the subcylinder $`[1,\mathrm{})\times Y`$. We consider the Seiberg-Witten equations perturbed by the self-dual part of $`iuf(d\eta )`$. Note that this perturbing two-form is $`iu\lambda f`$ times the $`(1,1)`$ form of the standard cylindrical-end metric on $`N`$ (see ), where
$$\lambda =\frac{2\pi \mathrm{deg}Y}{\mathrm{Vol}(\mathrm{\Sigma })}$$
(here, $`\mathrm{Vol}(\mathrm{\Sigma })`$ denotes the volume of $`\mathrm{\Sigma }`$). Similarly, we can extend the perturbation over $`Y`$ to a self-dual two-form perturbation of the equations over $`X_0`$ (and, consequently, $`X(T)`$ to all $`T>2`$). Denote the corresponding moduli spaces by $`_{N,u}(𝒥)`$, $`_{N,u}(C)`$, $`_{X_0,u}(𝒥)`$, $`_{X_0,u}(C)`$, and $`_{X(T),u}`$. Strictly speaking, we still have to show that these perturbed moduli spaces $`_{X(T),u}(𝔰)`$ can be used to calculate the Seiberg-Witten invariant in either chamber. This is clear because we can always choose a compactly-supported perturbing two-form $`\eta _0`$ whose integral against $`\omega _g`$ dominates the integral of $`\omega _g`$ against $`iuf(d\eta )^+`$. The key point is that the latter integral is finite, since $`\omega _g`$ decays exponentially (see ).
We now have the following analogue of Proposition 5.5
###### Proposition 8.2.
Suppose $`c_1(𝔰),[\mathrm{\Sigma }]=n`$, and let $`u`$ be a real number with $`0<u<2`$. Then the perturbed moduli space $`_{N,u}(𝒥)`$ contains only reducibles. Moreover, the space of reducibles is smoothly identified with the Jacobian torus $`𝒥`$ (i.e. the kernel and the cokernel of the Dirac operator coupled to any reducible vanishes). Furthermore, $`_{N,u}(C)`$ is empty.
Proof. We begin by proving $`_{N,u}(C)`$ is empty. Note that $`C`$ consists entirely of $`\alpha `$-solutions, hence so must any section in $`_{N,u}(C)`$. Thus, a solution $`(A,\alpha )_{N,u}(C)`$ induces a non-zero element in $`H^0(\widehat{E})`$ with
$`c_1(\widehat{E}),[\mathrm{\Sigma }_{}]=n+g1`$ and $`c_1(\widehat{E}),[\mathrm{\Sigma }_+]=g1.`$
But $`H^{}(R,\widehat{E})0`$, according to Proposition 7.5. The same argument, now appealing to Theorem 7.4, shows that $`_{N,u}^{}(𝒥)`$ is empty, and that $`𝒥`$ is smooth.
###### Proposition 8.3.
Suppose that
$$c_1(𝔰),[\mathrm{\Sigma }]=n,$$
and let $`u`$ be any real number with $`0<u<2`$. Then according to Theorem 8.1, $`𝒩_u(𝔰|_Y)`$ has two components, $`𝒥`$ and $`C`$, where $`C`$ is diffeomorphic to $`\mathrm{Sym}^{g1}(\mathrm{\Sigma })`$. Furthermore, the expected dimensions of the moduli spaces over $`N`$ and $`X_0`$ are given by:
(36) $`\text{e-dim}_{N,u}(𝒥)`$ $`=`$ $`2g1`$
(37) $`\text{e-dim}_{N,u}(C)`$ $`=`$ $`2g2`$
(38) $`\text{e-dim}_{X_0,u}(𝒥)`$ $`=`$ $`2d`$
(39) $`\text{e-dim}_{X_0,u}(C)`$ $`=`$ $`2d.`$
Moreover, $`_{N,u}^{}(𝒥)`$, $`_{N,u}(C)`$, $`_{X_0,u}(𝒥)`$, and $`_{X_0,u}(C)`$ are transversally cut out by the Seiberg-Witten equations (in particular, they are manifolds of the expected dimension). Furthermore, the boundary map
$$\rho :_{N,u}(C)C$$
is an orientation-preserving diffeomorphism onto its image.
Proof. The proofs of Propositions 7.9 and 7.10 apply directly in this perturbed context.
## 9. Cohomology
The Seiberg-Witten invariant is obtained from pairings of certain canonical cohomology classes on the Seiberg-Witten moduli space. These cohomology classes are inherited from the configuration spaces in which the moduli spaces live. In this section, we recall the definitions of these classes and discuss natural geometric representatives for them. (See Chapter 5 of for the corresponding discussion of cohomology relevant to Donaldson invariants.)
Let $`X`$ be a Riemannian four-manifold with a $`\mathrm{Spin}_{}`$ structure $`𝔰`$ specified by the pair of Hermitian $`^2`$ bundles $`W^+`$ and $`W^{}`$, and the Clifford action
$$\rho :TXW^+W^{}.$$
The Seiberg-Witten pre-configuration space is the space
$$𝒞(W^+)=𝒜(W^+)\times \mathrm{\Gamma }(X;W^+)\mathrm{\Omega }^1(W;)\times \mathrm{\Gamma }(X;W^+),$$
where $`𝒜(W^+)`$ denotes the space of connections compatible with some fixed connection $``$ on $`TX`$, and the isomorphism above is induced by comparing any connection $`A`$ against some fixed connection $`A_0`$. The irreducible pre-configuration space $`𝒞^{}(W^+)`$ is the subset of $`𝒞(W^+)`$ consisting of pairs $`(A,\mathrm{\Phi })`$, where $`\mathrm{\Phi }0`$. Now, $`𝒞^{}(W^+)`$ is weakly contractable, and the space $`\mathrm{Map}(X;S^1)`$ acts freely on it, so the irreducible configuration space, which is
$$^{}(W^+)=𝒞^{}(W^+)/\mathrm{Map}(X;S^1)$$
is weakly homotopy equivalent to the classifying space of $`\mathrm{Map}(X;S^1)`$. Now,
$$\mathrm{Map}(X;S^1)\mathrm{Map}(X;S^1)_e\times \pi _0(\mathrm{Map}(X;S^1))S^1\times H^1(X;);$$
so
$$(W^+)^{\mathrm{}}\times \frac{H^1(X;)}{H^1(X;)},$$
and
$$H^{}((W^+);)[U]\mathrm{\Lambda }^{}H^1(X;),$$
where $`U`$ is a generator with grading two. More invariantly, we define
$$𝔸(X)=[H_0(X;)]\mathrm{\Lambda }^{}H_1(X;),$$
graded by declaring $`H_0(X;)`$ to have grading two and $`H_1(X;)`$ to have grading one. Then, we have seen that
$$H^{}((W^+);)𝔸(X).$$
We describe two functorial mechanisms for constructing generators in $`H^{}(^{}(W^+);)`$. Over the space $`X\times ^{}(W^+)`$, there is a universal line bundle $`=X\times S^1\times 𝒞^{}(W^+)/\mathrm{Map}(X,S^1)`$, where the action is defined by
$$u(x,\zeta ,A,\mathrm{\Phi })=(u,u(x)\zeta ,A+u^1du,u\mathrm{\Phi }).$$
Using this class we can define a “$`\mu `$-map”
$$\mu :(H_0H_1)(X;)H^{}(^{}(W^+)),$$
which sends a homology class of degree $`i`$ to a cohomology class of degree $`2i`$, by
$$\mu (x)=c_1()/x;$$
i.e. $`\mu (x)`$ is the cohomology class on $`^{}(W^+)`$ with the property that for any homology class $`cH_{}(^{}(W^+))`$,
$$\mu (x),c=c_1(),x\times c.$$
We describe another convenient mechanism for constructing one-dimensional cohomology in $`𝒞(W^+)`$ as follows. A closed curve $`x:S^1X`$ induces a map
$$\mathrm{Hol}_x:(W^+)S^1$$
which is defined to be the holonomy of the connection $`A`$ around the curve $`x`$. The pull-back of the volume form $`d\theta `$ of $`S^1`$ by this map gives rise to a one-dimensional cohomology class $`\mathrm{Hol}_x^{}(d\theta )`$ associated to $`x`$, which we call the holonomy class around $`x`$.
###### Proposition 9.1.
The cohomology groups of the configuration space $`(W^+)`$ are generated by the image of the $`\mu `$-map. Moreover, given $`xH_1(X;)`$, $`\mu (x)`$ is the holonomy class around $`x`$, $`\mathrm{Hol}_x^{}(d\theta )|_{^{}(W^+)}`$.
###### Remark 9.2.
Note that $`\mathrm{Hol}_x^{}(d\theta )`$ is naturally defined over the entire configuration space
$$(W^+)=𝒞(W^+)/\mathrm{Map}(X,S^1)\frac{H^1(X;)}{H^1(X;)}.$$
Proof. We begin by proving the second claim. Note that $``$ comes with a tautological connection along the $`X`$ factor, with the property that for any path $`\beta :S^1X`$ and connection $`A𝒞(W^+)`$,
(40)
$$\mathrm{Hol}_{\beta \times A}()=\mathrm{Hol}_\beta (A).$$
$`𝒞(W^+)`$ Now, fix a path in $`X`$,
$$\beta :S^1X.$$
We need to show that for all paths in the configuration space
$$\alpha :S^1(W^+)(𝔰),$$
we have that
$$c_1(),\alpha \times \beta =\mathrm{deg}(\mathrm{Hol}_\beta \alpha :S^1S^1).$$
This follows from the fact that for a line bundle $`L`$ over the the torus $`S^1\times S^1`$, the first Chern number is the degree of the map from $`S^1\times S^1`$ defined by
$$x\mathrm{Hol}_{x\times S^1}L$$
(a map which makes sense only after one puts a connection on $`L`$, but the degree is independent of this connection, so we left it out of the notation), together with the universal property of Equation (40). Thus, we have identified $`\mu `$ on any one-dimensional homology class.
The rest of the proposition is established, once we see that for a point $`xX`$, $`\mu [x]`$ generates $`H^2`$ of the configuration space. But this follows easily from the fact that $`\mathrm{Map}(X,S^1)_e`$ acts freely on the space of irreducible configurations.
With this concrete understanding of the $`\mu `$-classes, we turn to a discussion of submanifold representatives for them.
Given a point $`xX`$, let $`_x`$ denote the line bundle associated to the base fibration of $`X`$; i.e. it is the restriction of the universal line bundle $``$ to the slice $`^{}(W^+)\{x\}\times ^{}(W^+)X\times ^{}(W^+)`$. Given a point in the fiber $`\mathrm{\Psi }(x)W_x^+`$, we can construct a canonical section over $`(X,W^+)`$ by
$$[A,\mathrm{\Phi }][A,\mathrm{\Phi },\mathrm{\Phi }(x),\mathrm{\Psi }(x)].$$
The zero set of this section in $`(X,W^+)`$ is a codimension-two submanifold representing $`\mu [x]`$. The restriction of this section to a moduli space $`_X(𝔰)(X,𝔰)`$ is not, in general, transverse. However, by mollifying the construction appropriately, we can find a section which is generic over the moduli space, and hence obtain a divisor $`V(x)`$ representing $`\mu [x]`$, as follows.
###### Definition 9.3.
Fix a ball $`BX`$ around $`x`$ and a non-vanishing section $`\mathrm{\Psi }`$ of $`W^+|_B`$. Given a self-dual two-form $`\lambda `$ which is compactly supported over $`B`$, the $`\lambda `$-mollified section is the section of $`_x`$ defined by
$$[A,\mathrm{\Phi }][A,\mathrm{\Phi },_B\lambda \mathrm{\Phi },\mathrm{\Psi }].$$
###### Lemma 9.4.
There are $`L^2`$ sections $`\lambda `$ compactly supported in $`B`$ so that the $`\lambda `$-mollified section of $`_x`$, restricted to the moduli space $`_X`$, vanishes transversally.
Proof. Fix a compactly-supported cut-off function $`\beta `$ in $`B`$, and consider the section
$$[A,\mathrm{\Phi }]\times \lambda [A,\mathrm{\Phi },_B\lambda \mathrm{\Phi },\mathrm{\Psi }\beta ]$$
of $`\pi _2^{}(_x)`$, thought of as a line bundle over $`\mathrm{\Omega }^+(B)\times (X,𝔰)`$, giving $`\mathrm{\Omega }^+(B)`$ the $`L^2`$ topology. This transversality follows from the fact that, for any $`[A,\mathrm{\Phi }]_X`$, as we vary $`\lambda `$, the integral $`_B\lambda \mathrm{\Phi },\mathrm{\Psi }\beta `$ can take on any complex value. This, in turn, follows from the unique continuation theorem for elliptic differential operators, which guarantees that the section $`\mathrm{\Phi }`$ cannot vanish identically over $`B`$.
###### Remark 9.5.
In effect, the above lemma tells us how to construct a divisor representative $`V(x)`$ for $`\mu [x]`$ when $`[x]H_0(X)`$; this divisor is represented by the zero-set of the $`\lambda `$-mollified section of $`_x`$. Finding codimension-one representatives for $`\mu [\gamma ]`$, where $`\gamma H_1(X)`$ is even easier: one need only find a regular value $`\theta `$ for the map
$$\mathrm{Hol}_\gamma :_X(𝔰)S^1.$$
Then, $`\mathrm{Hol}_\gamma ^1(\theta )`$ is the submanifold $`V(\gamma )`$ representing $`\mu [\gamma ]`$ over $`_X(𝔰)`$.
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# Introduction
## Introduction
The interaction between surface acoustic waves (SAW’s) and mobile carriers in quantum wells is an important method to study the dynamic properties of two-dimensional (2D) systems. The SAW can trap carriers and induce acoustic charge transport as has been investigated in a number of systems in view of possible device applications . Also, the SAW-method was applied to study the quantum Hall effects , electron transport through a quantum-point contact , lateral nanostructures , and commensurability effects in a 2D system . However, all those experiments have been done in the regime of small signals and linear interaction. A recent Letter by Rotter et al. reports strongly nonlinear acousto-electric effects in a 2D electron gas (2DEG), which become possible in hybrid structures based on $`A_3B_5`$ semiconductors and $`LiNbO_3`$ \[9-11\]. In these experiments an intense SAW breaks a 2DEG into moving electron stripes and all characteristics of the acousto-electric interaction are strongly modified as compared to the linear case. In modern hybrid structures the SAW-induced potential amplitude can become comparable with the band-gap of a semiconductor. The previous Letter on nonlinear effects in the hybrid structures with a 2DEG includes a brief qualitative analysis. Here we present a detailed theoretical study of nonlinear acousto-electric effects in a 2D electron system and develop a coupled-amplitude method for intense SAW’s interacting with a 2DEG. Using our theoretical results we can explain main experimental observations. Including the effect of electron diffusion we find a good quantitative agreement between theory and experiment for the case of the SAW absorption coefficient in the nonlinear regime.
Nonlinear acoustic waves in bulk piezocrystals with free carriers were discussed in a number of theoretical papers \[12-16\]. The application of a dc voltage to the crystal can result in the current amplification of sound and in the formation of stationary nonlinear waves . Analytic results in the theory of nonlinear acoustic waves in bulk piezocrystals with free carriers were mostly obtained in the limit of small amplitudes or for the case of very intense acoustic waves . For nonlinear SAW’s in crystals with a 3D electron gas, a theory was developed in the limit of very high amplitudes, when the SAW bunches electrons near the crystal surface . Another theoretical aspect related to the generation of the second harmonic of a SAW was studied in Ref. by using the coupled-amplitude method and perturbation theory. A theory of acousto-electric interactions in a 2D electron system was developed mostly for the linear regime of interaction .
Here, we study both theoretically and experimentally the transition from the linear regime of the acousto-electric interaction to the limit of strongly nonlinear effects in a 2DEG. Our theoretical results are applied for a description of experimental data on hybrid structures . We pay attention to density-dependences of the absorption coefficient and the SAW-velocity shift because the electron density is a tunable parameter in experiments on 2D systems. Such dependences were not discussed in detail in the context of 3D systems \[12-16\]. It turns out that for low densities, the absorption coefficient is a decreasing function of the sound intensity caused by the trapping of electrons in the SAW piezoelectric potential. At sufficiently high electron density, however, the absorption coefficient is a non-monotonous function of the sound intensity. This behavior can be understood in terms of a dynamical screening effect. Also, our quantitative analysis shows that the absorption coefficient at room temperature is strongly reduced due to electron diffusion. A non-monotonous behavior with increasing sound intensity was also found for the intensities of higher harmonics in a short device.
The coupled-amplitude method with the introduction of fast and slow variables was used before for a description of bulk systems . However, the formulas from the bulk theory can not be directly applied to surface waves because of the complex character of the lattice vibrations in a SAW. To develop the coupled-amplitude method in a 2D system, we introduce a local velocity shift and a local absorption coefficient in an integral form through an electron current, a SAW potential, and an electro-mechanical coupling coefficient of a microstructure. The resulting formulas can then be applied to any type of SAW’s interacting with a 2DEG. Moreover, using our approach, we can find solutions for higher harmonics for arbitrary SAW intensity. This is in contrast with perturbation methods developed earlier for 3D systems and valid at small SAW intensities . At very high intensities, we find analytical asymptotic dependences for the high-harmonic intensities in a 2D system.
As has been shown before, the linear approximation holds when $`\delta nN_s`$, where $`N_s`$ is the equilibrium 2DEG density and $`\delta n`$ is a density perturbation due to a SAW. In the linear theory, the absorption coefficient $`\mathrm{\Gamma }^0`$ and the SAW-velocity shift due to the 2DEG $`\delta v_s^0`$ are given by the well-known relation
$`{\displaystyle \frac{\delta v_s^0}{v_s^0}}+i{\displaystyle \frac{\mathrm{\Gamma }^0}{2q}}={\displaystyle \frac{K_{eff}^2(q)}{2}}{\displaystyle \frac{i\sigma _0/\sigma _m}{1+i(\sigma _0/\sigma _m+D_eq/v_s^0)}},`$ (1)
where $`q`$ is the SAW wave vector, $`v_s^0`$ is the sound velocity in the absence of a 2DEG, $`\sigma _0`$ and $`D_e`$ are the 2DEG conductivity and the diffusion coefficient, respectively. $`K_{eff}^2`$ is the effective electro-mechanical coupling coefficient. $`\sigma _m(q)=v_s^0ϵ_{eff}(q)/(2\pi )`$, where $`ϵ_{eff}(q)`$ is the effective dielectric constant in a 2D system. In most piezoelectric crystals, the coupling $`K_{eff}^21`$. In $`GaAs`$, $`K_{eff}^2=0.00064`$, whereas in the hybrid structures as studied here, it is two orders of magnitude larger, in the range of $`0.010.05`$ , but still much less then unity. The goal of this paper is to describe the acousto-electric interaction for the case of large amplitude SAW’s, when $`\delta nN_s`$ and the perturbation theory is not longer valid. At the same time, the coupling $`K_{eff}^2`$ will be assumed to be much less than unity. Below, we will generalize the results following from Eq. 1 for the strongly nonlinear case, when $`\delta nN_s`$.
The paper is organized in the following way. In the first section, we will give the general equations for SAW’s on a piezoelectric crystal. The second section is devoted to a coupled-amplitude method developed for the case of intense SAW’s. The third and fourth sections are about phenomena related to large-amplitude SAW’s in a 2DEG. Then, we will discuss experimental data on the hybrid structures and apply our theoretical results for the interpretation of experiments.
## 1 Model and general equations
In usual GaAs-based microstructures it is very difficult to realize SAW’s with high-amplitude potentials because the electro-mechanical coupling in GaAs is relatively weak. A strong piezoelectric interaction can be achieved in hybrid structures (Fig. 1). Those consist of a semiconductor layer being bonded to a piezoelectric host crystal, in our case, $`LiNbO_3`$ . The semiconductor layer contains an $`InGaAsAlGaAs`$ quantum well (QW) with a high-quality 2DEG, to which Ohmic contacts are formed. The distance between the QW and the piezocrystal is only $`32nm`$, whereas the distance between the QW and the top transport gate is $`d=450nm`$. In our model $`x`$ and $`y`$ are the in-plane coordinates and $`z`$ is the normal one. The QW plane corresponds to $`z=d`$, and the SAW travels in the $`x`$-direction (see Fig. 1). By changing the transport-gate voltage $`V_t`$ one can tune the electron density $`N_s(V_t)`$ in a QW. In this structure, traveling SAW’s can induce very strong piezoelectric fields in the semiconductor layer due to the strong piezoelectricity of the host $`LiNbO_3`$ crystal. A SAW is induced and detected by the metallic interdigital transducers, $`IDT1`$ and $`IDT2`$, respectively, at room temperature . The acousto-electric current is measured between two Ohmic contacts labeled $`1`$ and $`2`$ in Fig. 1.
SAW’s in a piezoelectric crystal with a 2D plasma are described by the system of equations
$`\rho \ddot{u}_i=c_{iklm}_m_ku_l+p_{lik}_l_k\varphi `$ (2)
$`ϵ(x_3)_i_i\varphi 4\pi p_{ikl}_i_lu_k=4\pi en\delta (x_3d),`$ (3)
where $`e=|e|`$ is the electron charge, $`\rho `$ is the mass density, and $`ϵ`$ is the dielectric constant. Further, $`c_{iklm}`$ is the elastic tensor, $`p_{lik}`$ is the piezoelectric tensor, $`u_k(x,z,t)`$ is the lattice displacement, and $`\varphi (x,z,t)`$ and $`n(x,t)`$ are the electrostatic potential and the 2D electron density, respectively. In Eqs. 2,3, we have used the notations $`_lf=f/x_l`$ and $`\dot{f}=f/t`$, and the sum convention for repeated indexes. $`t`$ is the time, and $`x_3=z`$, $`x_2=y`$, and $`x_1=x`$ are the coordinates.
In our geometry, the SAW propagating in the x-direction is a purely Rayleigh wave, in which only two components of the displacement, $`u_x`$ and $`u_z`$, are nonzero. Hence, the electric field $`𝐄`$ is also polarized in the $`(xy)`$-plane. This case corresponds to the hybrid structures studied in experiments , where the $`128^0`$-rotated $`Y`$-cut of $`LiNbO_3`$ is used. The surface of the thin $`GaAs`$ film is $`(001)`$. The SAW propagates in the $`[110]`$-direction of $`GaAs`$ and $`X`$-direction of $`LiNbO_3`$.
The electron 2D plasma is described by the usual hydrodynamic equations
$`e\dot{n}(x,t)+{\displaystyle \frac{j(x,t)}{x}}=0,`$ (4)
$`j=\sigma {\displaystyle \frac{\varphi (x,d,t)}{x}}eD_e{\displaystyle \frac{n(x,t)}{x}},`$ (5)
where $`j`$ is the 2D electron current, $`\sigma =|e|\mu n_s(x,t)`$ is the 2DEG conductivity, $`\mu `$ is the mobility, and $`D_e`$ is the electron diffusion coefficient. Eq. 5 is valid in the long wavelength limit, when $`ql_e1`$, where $`l_e`$ is the electron mean free path.
The wave equations (2,3) should be solved together with the boundary conditions at the surface $`z=0`$ and at the interface $`z=d`$. At $`z=0`$, these conditions are the following: $`\varphi (z=0)=0`$ and $`\sigma _{zi}=c_{zilm}(_mu_l+_lu_m)/2+p_{lzk}_l\varphi =0`$. Here $`\sigma _{zi}`$ is the $`z`$-component of the stress. The top metal gate is thin and does not influence the boundary condition for the stress tensor. At the semiconductor-piezocrystal interface, $`\varphi `$, $`u_i`$, and $`\sigma _{zi}`$ should be continuous, and $`D_z(z=d+\delta )D_z(z=d\delta )=4\pi en`$, where $`\delta 0`$ and $`D_i=ϵ(z)_i\varphi +2\pi p_{ikl}(_ku_l+_lu_k)`$ is the electric displacement. For simplicity, we assumed above that the 2DEG is located directly at the semiconductor-piezocrystal interface.
## 2 Coupled-amplitude method
The system of nonlinear equations (2-5) can be simplified in the limit of weak electro-mechanical coupling, $`K_{eff}^2p^2/(cϵ)1`$. In this limit, we can introduce two coordinates, the ”slow” variable $`x`$ and the ”fast” variable $`x_1=xv_s^0t`$ . The solution is of the form (see Appendixes 1 and 2),
$`𝐮(x,z,t)=𝐮(x,x_1,z)=𝐮_0(x,z)+{\displaystyle \underset{n=1,2,\mathrm{}}{}}a_n(x)𝐔^0[z;q_n+\delta q_n(x)]e^{iq_nx_1}+c.c.`$ (6)
$`\varphi (x,z,t)=\varphi (x,x_1,z)=\varphi _0(x,z)+{\displaystyle \underset{n=1,2,..}{}}a_n(x)\mathrm{\Phi }^0[z;q_n+\delta q_n(x)]e^{iq_nx_1}+c.c.,`$ (7)
This solution is written as a sum of harmonics with wave vectors $`q_n=nq`$, where $`q>0`$ is the wave vector of the initially generated SAW near $`x=0`$. The vector $`𝐀^0[z;q]=(𝐔^0[z;q];\mathrm{\Phi }^0[z;q])`$ and the quantity $`\delta q_n`$ are determined by a linear system of equations as given in the Appendix 2. The envelope functions $`a_n(x)`$ are slowly changing on the scale of $`\lambda =2\pi /q`$. It is assumed that the SAW intensity related to the vector $`𝐀^0[z;q]`$ is unity and thus the total SAW intensity is $`I_{saw}=_{n=1,2,..}|a_n(x)|^2`$. The functions $`𝐮_0(x,z)`$ and $`\varphi _0(x,z)`$ describe the static spatial distributions, that can be induced by a SAW.
The functions $`j(x,x_1,t)`$ and $`n(x,x_1,t)`$ can be written in the standard form $`f(x,x_1,t)=f_0(x)+_{n=1,2,\mathrm{}}f_n(x)e^{iq_nx_1}+c.c.`$, where $`f_n(x)`$ is an envelope function. Also, the electric field $`𝐄`$ can be written in the way similar to Eq. 7.
The solution (6,7) is a sum of linear-like SAW’s, which slowly vary in the amplitude and in the $`z`$-profile. The $`z`$-distribution of lattice displacement $`𝐔^0`$ is a sum of exponential functions $`\mathrm{exp}[\gamma _j(q_n+\delta q_n(x))z]`$, where $`\gamma _j(q)`$ are the coefficients depending also on material constants . On short distances ($`\lambda `$) the envelope functions $`a_n(x)`$ can be regarded as constants and we can solve the equations (2-5) considering only the ”fast” variable $`x_1`$. $`n_n(x)`$ and $`\varphi _n(x)`$ should be considered here as the functions of the parameters $`a_1,a_2,\mathrm{}`$, $`n_0`$, and $`E_0`$. Then, having $`n_n`$ and $`\varphi _n`$ as functions of $`a_1,a_2,\mathrm{}`$, $`n_0`$, and $`E_0`$, we can find the behavior of $`a_n(x)`$ on long-range scale, $`x1/\mathrm{\Gamma }^0\lambda `$.
The electrostatic potential is written in a self-consistent way:
$`\varphi (x,x_1,z)=\varphi ^{ind}+\varphi ^{saw},`$ (8)
where $`\varphi ^{ind}`$ and $`\varphi ^{saw}`$ are the potentials induced by a 2DEG and by piezoelectric charges of a SAW. Using Eq.3 we write
$`ϵ(z)_i_i\varphi ^{ind}=4\pi en\delta (zd)`$ (9)
$`ϵ(z)_i_i\varphi ^{saw}=4\pi p_{ikl}_i_lu_k.`$ (10)
$`\varphi ^{ind}`$ and $`\varphi ^{saw}`$ can be expressed by the harmonic amplitudes $`\varphi _n^{ind}`$ and $`\varphi _n^{saw}`$. For example, $`\varphi ^{ind}=\varphi _0^{ind}(x,z)+_{n=1,2,\mathrm{}}\varphi _n^{ind}(x,z)e^{iq_nx_1}+c.c.`$. In the limit $`\mathrm{\Gamma }_{max}\lambda K_{eff}^21`$, we find from Poisson’s equation (see Appendix 1) and from the conservation of charge
$`\varphi _n^{ind}(x,d)={\displaystyle \frac{2\pi en_n(x)}{ϵ_{eff}(q_n)q_n}},E_{nx}^{ind}(x,d)=iq_n\varphi _n^{ind}(x,d),`$ (11)
$`j_n(x)=v_s^0en_n(x),`$
where $`n=1,2,3\mathrm{}`$ . Here $`ϵ_{eff}(q)=[ϵ_p+ϵ_scoth(|q|d)]/2`$ is the effective dielectric constant including the gate electrode effect, and $`ϵ_p`$ and $`ϵ_s`$ are the dielectric constants of a host piezocrystal and a semiconductor, respectively .
The $`n`$th-harmonic of the SAW-potential $`\varphi _n^{saw}`$ is given only by $`a_n(x)`$ and by material constants and can be easily found from the Poisson equation (10). At $`z=d`$ we have $`\varphi _n^{saw}(x,d)=C_na_n(x)`$, where the coefficient $`C_n`$ depends on the geometry and the material constants. For example, in a crystal of the type of $`GaAs`$, $`C_n=p_4g(q_n)`$, where $`g(q_n)`$ is a complicated function of $`q_n`$. Below, we will give the necessary relations for the hybrid structures.
To find the harmonic amplitudes $`n_n`$ and $`\varphi _n`$, we have to solve Eqs. 3, 4, and 5 regarding the ”slow” variable $`x`$ as the constant. The slowly-varying quantities $`n_n(x)`$ and $`j_n(x)`$ can be found as Fourier components of the solution $`n(x_1,x)`$ from a ”fast” equation in terms of $`x_1`$. In a self-consistent approach, the electron density $`n(x_1,x)`$ is determined by the SAW-induced potential at $`z=d`$, $`\varphi ^{saw}(x_1,x,d)=\varphi _0^{saw}(x,d)+_{n=1,2,\mathrm{}}C_na_n(x)e^{iq_nx_1}+c.c.`$, through a non-linear equation . By using Poisson’s equation, the results of Appendix 1, and Eqs. 4 and 5, we obtain
$`|e|n(x_1,x)\mu [{\displaystyle \frac{d}{dx_1}}\{{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑x_1^{}G(x_1x_1^{})n(x_1^{},x)\}+E^{saw}(x_1,x)]`$ (12)
$`eD_e{\displaystyle \frac{dn(x_1,x)}{dx_1}}ev_{s0}n(x_1,x)=b_0,`$
$`G(x_1x_1^{})=e{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑k{\displaystyle \frac{e^{ik(x_1x_1^{})}}{|k|ϵ_{eff}(|k|)}},`$ (13)
where $`b_0`$ is a constant, which occurs after one integration in the conservation-of-charge equation (4). Taking into account $`\mathrm{\Gamma }^0\lambda 1`$, we assume that the solution of Eq. 12 is periodic in $`x_1`$, $`n_s(x_1,x)=n_s(x_1+\lambda ,x)`$. Also, $`<n_s(x_1,x)>=_0^\lambda n_s(x_1,x)𝑑x_1/\lambda =n_0(x)`$, where $`n_0(x)`$ plays the role of a ”local” 2D density.
The functions $`a_n(x)`$ are connected by the system of nonlinear equations (see Appendix 2)
$`{\displaystyle \frac{da_n}{dx}}=i\delta q_n[a_1,a_2,\mathrm{};n_0,E_0]a_n.`$ (14)
where
$`\delta q_n[a_1,a_2,\mathrm{};n_0,E_0]={\displaystyle \frac{K_{eff}^2(q_n)}{2}}{\displaystyle \frac{2\pi }{ϵ_{eff}(q_n)}}{\displaystyle \frac{en_n(x)}{\varphi _n^{saw}(x)}}=i{\displaystyle \frac{K_{eff}^2(q_n)}{2q_n\sigma _m(q_n)}}{\displaystyle \frac{j_n(x)E_n^{saw}}{|\varphi _n^{saw}(x)|^2}}`$ (15)
and $`n=1,2,3\mathrm{}`$ .
The local velocity change of a SAW $`\delta v_n`$ and the local absorption coefficient $`\mathrm{\Gamma }_n`$ can be expressed by $`\delta q_n`$:
$`\delta q_n=q_n{\displaystyle \frac{\delta v_n}{v_s^0}}+i{\displaystyle \frac{\mathrm{\Gamma }_n}{2}}.`$ (16)
Using the relation
$`K_{eff}^2(q_n)={\displaystyle \frac{2|E_n^{saw}|^2\sigma _m(q_n)}{q_nI_n}},`$ (17)
(see Ref. ) and Eqs. 11 we write
$`{\displaystyle \frac{\delta v_n(x)}{v_s^0}}={\displaystyle \frac{<\stackrel{~}{\varphi }_n^{saw}(x_1,x)j(x_1,x)>}{2I_n(x)}}`$ (18)
$`\mathrm{\Gamma }_n(x)={\displaystyle \frac{<\stackrel{~}{E}_n^{saw}(x_1,x)j(x_1,x)>}{I_n(x)}}.`$ (19)
Here we use the notations $`<f(x_1,x)>=_0^\lambda f(x_1,x)𝑑x_1/\lambda `$, and $`\stackrel{~}{f}_n(x_1,x)=f_n(x)exp(iq_nx_1)+c.c.`$ .
In this section, we have assumed that $`K_{eff}^21`$ and neglected the terms $`d^2a_n/dx^2\delta q_nda_n/dxK_{eff}^4`$ and $`d\delta q_n/dxK_{eff}^4`$. The static electric fields $`\varphi _0(x,z)`$ in Eq. 7, which can be induced by a SAW, will not play an important role in this paper because we will consider the case with no voltage applied to the Ohmic contacts ($`V_1=V_2`$) and a relatively short device with $`L1/\mathrm{\Gamma }^0`$. Thus, $`\varphi _0(x,d)const`$ and the static electric fields $`E_{0x}(x,d)`$ can be neglected. Eqs. 14 and 15 can now be applied to various types of SAW’s by introducing specific electro-mechanical coupling coefficients.
## 3 Acousto-electric transport in a two-dimensional plasma
Here, we intend to consider the SAW absorption and the acoustic charge transport in the regime $`V_1=V_2`$. The behavior of a 2DEG in an intense SAW at small distances ($`\lambda `$) can be assumed to be periodic and is described by nonlinear equation (12). At long distances ($`\lambda `$) the SAW behavior is determined by complex amplitudes $`a_n(x)`$, that can be found by Eqs. 14 and 15. First, we define the boundary conditions at $`x=0`$: $`a_1(0)=\sqrt{I_1(0)}`$ and $`a_n(0)=0`$ for $`n=2,3,\mathrm{}`$. Here, $`I_1(0)`$ denotes the SAW intensity generated by IDT1. In this case, for a relatively short sample the SAW contains mostly the fundamental SAW harmonic $`n=1`$ with a small admixture of the higher harmonics $`n=2,3,\mathrm{}`$. So, it follows from Eqs. 14 and 15
$`a_1(x)=a_1(0)\left(1+ix\delta q_1[a_1(0),0,0,\mathrm{};n_0,0]\right).`$ (20)
We assume that $`[I_1(0)I_1(L)]I_1(0)`$, where $`L`$ is the length of a semiconductor film and $`I_1(L)`$ is the intensity detected at the IDT2. The latter is valid in a short sample where $`\mathrm{\Gamma }_1L1`$. The first-harmonic absorption coefficient per unit length, which is measured in the experiments, is $`[I_1(0)I_1(L)]/(I_1(0)L)=Im\left(\delta q_n[a_1(0),0,0,\mathrm{};n_0,0]\right)=\mathrm{\Gamma }_1(a_1,0,0,\mathrm{};n_0,0)\mathrm{\Gamma }_1`$ and the velocity shift is $`\delta v_1(a_1,0,0,\mathrm{};n_0,0)/v_s^0=Re\left(\delta q_1[a_1,0,0,\mathrm{};n_0,0]/q\right)\delta v_1/v_s^0`$.
The amplitudes of higher harmonics in a short sample turn out to be given by:
$`a_n(x)=ix\underset{a_n0}{lim}a_n\delta q_n[a_1(0),0,0,..,a_n,\mathrm{};n_0,0]=ix{\displaystyle \frac{\pi K_{eff}^2(q_n)}{ϵ_{eff}(q_n)}}{\displaystyle \frac{en_n(0)}{C_n}},`$ (21)
where $`n=2,3,4,\mathrm{}`$ . Here, we assume that $`[I_1(0)I_1(L)]I_n`$, which is valid for a short sample where $`\mathrm{\Gamma }_1L1`$, and we take into account that for the strongly nonlinear case $`\mathrm{\Gamma }_1\mathrm{\Gamma }_n`$. The intensity of the $`n`$th-harmonic is given by $`I_n(x)=|a_n(x)|^2`$ and
$`|C_n|^2={\displaystyle \frac{|\varphi _n^{saw}|^2}{|a_n|^2}}={\displaystyle \frac{K_{eff}^2(q_n)}{2q_n\sigma _m(q_n)}}`$ (22)
(see Eq. 17). By using Eq. 21 we find the intensities of high harmonics at $`x=L`$: $`I_n(L)=lim_{a_n0}|a_n\delta q_n[a_1,0,0,\mathrm{},a_n,\mathrm{};0,0]|^2L^2=\pi L^2(K_{eff}^2/ϵ_{eff})q_nv_s^0e^2|n_n(0)|^2`$.
To calculate $`\mathrm{\Gamma }_1`$ and $`\delta v_1`$, we numerically solve Eq. 12 for the parameters $`a_10`$ and $`a_n=0`$ when $`n=2,3,\mathrm{}`$. Then we can find the Fourier components of $`n(x_1,a_1)`$. We now calculate the electron density $`n(x_1,a_1)`$ in the long wave-length limit $`qd1`$. In this limit the function given by Eq. 13 is reduced to $`G(x_1x_1^{})=(4\pi ed/ϵ_s)\delta (x_1x_1^{})`$ and Eq. 12 now writes as
$`|e|n(x_1,a_1)\mu \{{\displaystyle \frac{4\pi ed}{ϵ_s}}{\displaystyle \frac{dn(x_1,a_1)}{dx_1}}+F_1^{saw}cos(qx_1+\chi _1)\}eD_e{\displaystyle \frac{dn(x_1,a_1)}{dx_1}}`$ (23)
$`ev_s^0n(x_1,a_1)b_0=0,`$
where the SAW-induced piezoelectric field is taken in the form $`E^{saw}=F_1^{saw}cos(qx_1+\chi _1)`$, and $`\chi _1`$ is the phase of the first SAW harmonic. For simplicity, we put in the following $`\chi _1=0`$. The constant $`b_0`$ in Eq. 23 is directly connected with the kinetic motion of a SAW and thus vanishes for $`v_s^00`$. Dividing Eq. 23 by $`n`$ and then integrating over $`x_1`$ we get
$`|e|\mu \{{\displaystyle \frac{4\pi ed}{ϵ_s}}n(x_1)+{\displaystyle \frac{F_1^{saw}}{q}}\mathrm{sin}(qx_1)\}eD_eln[n(x_1)]=`$ (24)
$`const_1+ev_s^0x_1+b_0{\displaystyle _0^{x_1}}{\displaystyle \frac{dx_1^{}}{n(x_1^{})}}.`$
The left-hand side of the above equation is periodic and thus the right-hand side should be periodic as well. It implies that $`b_0_0^\lambda 𝑑x_1^{}/n(x_1^{})=ev_s^0\lambda `$.
A numerical solution of Eq. 23 $`n(x_1)`$ at $`T=300K`$ and at various averaged densities $`N_s=<n(x_1)>`$ is shown in the inset of Fig. 3. Here, we use the following parameters: $`\lambda =33\mu m`$, $`\mu =5000cm^2/Vs`$, $`D=\mu (KT/e)`$, $`ϵ_s=12.5`$, and $`v_s^0=3.810^5cm/s`$. It is seen that that with decreasing $`N_s`$ the formerly homogenous 2DEG turns into moving electron stripes.
It follows from Eqs. 18,19:
$`{\displaystyle \frac{\delta v_1}{v_s^0}}={\displaystyle \frac{<\mathrm{\Phi }_1^{saw}sin(qx_1)j(x_1)>}{2I_1}},\mathrm{\Gamma }_1={\displaystyle \frac{<F_1^{saw}cos(qx_1)j(x_1)>}{I_1}}.`$ (25)
The calculated absorption coefficient $`\mathrm{\Gamma }_1`$ as a function of the electron density $`N_s`$ for various potential amplitudes $`\mathrm{\Phi }_1^{saw}=F_1^{saw}/q`$ is shown in Fig. 2. It was calculated from Eqs. 25 and the numerical solution for $`n(x_1)`$. The potential amplitude $`\mathrm{\Phi }_1^{saw}`$ can be easily connected with the SAW intensity $`I_1`$ by using Eqs. 17 and 22 with $`K_{eff}^2`$ found numerically in Ref. . As an example, in the hybrid structures at SAW frequencies $`f=114MHz`$ and $`340MHz`$, $`K_{eff}^2`$ is about $`0.015`$ and $`0.035`$, respectively.
We see in Fig. 2 that with increasing $`\mathrm{\Phi }_1^{saw}`$ the absorption coefficient in general decreases and its maximum is shifted to the higher values of $`N_s`$. This nonlinear behavior can be understood qualitatively as follows. For the densities $`N_sN_s^{max}`$, the electron plasma forms the moving charge stripes and the electron velocity $`j/(eN_s)`$ is very close to its maximum $`v_s^0`$. Here $`N_s^{max}`$ denotes the density corresponding to the maximum of the function $`\mathrm{\Gamma }_1(N_s)`$ at fixed $`\mathrm{\Phi }_1^{saw}`$. We will show below that in the case $`N_s<N_s^{max}`$ $`\mathrm{\Gamma }_1|e|N_s(v_s^0)^2/(I_1\mu )N_s/I_1`$ . This asymptotic behavior follows from the Weinreich relation . In the region $`N_sN_s^{max}`$, the electric current $`\sigma E_x`$ decreases with increasing $`N_s`$ because of the screening effect. Thus, the absorption coefficient decreases as well.
In Fig. 3 we plot $`\mathrm{\Gamma }_1`$ as a function of the SAW potential $`\mathrm{\Phi }_1^{saw}`$ for a fixed 2DEG density. For electron densities less than about $`10^{10}cm^1`$, the function $`\mathrm{\Gamma }_1(\mathrm{\Phi }_1^{saw})`$ is always decreasing. At higher densities $`\mathrm{\Gamma }_1(\mathrm{\Phi }_1^{saw})`$ has a maximum. We attribute this behavior to the screening effect in a 2DEG. At high density and small $`\mathrm{\Phi }_1^{saw}`$, the absorption is strongly suppressed because of screening. The intense SAW, however, modulates the 2DEG and, consequently, reduces screening. Thus, the absorption coefficient starts to increase with increasing $`\mathrm{\Phi }_1^{saw}`$, when $`\mathrm{\Phi }_1^{saw}`$ is not so large. At larger $`\mathrm{\Phi }_1^{saw}`$, the plasma becomes broken into stripes and $`\mathrm{\Gamma }_1`$ decreases with $`I_1`$ as $`1/I_1`$ .
The velocity shift $`\delta v_1`$ numerically calculated by Eqs. 25 is shown in Fig. 4. With increasing SAW intensity $`\delta v_1`$ increases and approaches to the velocity $`v_s^0`$ that corresponds to the case of a totally depleted 2DEG. By analyzing Eqs. 25 and the function $`j(x_1)`$, we can see that, in the limit $`I_1\mathrm{}`$, $`\delta v_11/\mathrm{\Phi }_1^{saw}1/\sqrt{I_1}`$, like in a 3D plasma .
In Fig. 5 we show results of numerical calculations for the intensities of the high harmonics $`I_n/I_{n,max}=|n_n|^2/N_s^2`$, where
$`I_{n,max}=I_n(\mathrm{\Phi }_1^{saw}\mathrm{})=\pi L^2{\displaystyle \frac{K_{eff}^2}{ϵ_{eff}}}q_nv_s^0e^2N_s^2.`$
This formula was obtained taking into account that $`|n_n|N_s`$ when $`\mathrm{\Phi }_1^{saw}\mathrm{}`$. The quantities $`I_n`$ were calculated from the Fourier components of the function $`n(x_1)`$. At small amplitudes $`\mathrm{\Phi }_1^{saw}`$, $`I_n(\mathrm{\Phi }_1^{saw})^n`$, being typical for weak nonlinearity. At larger $`\mathrm{\Phi }_1^{saw}`$, the behavior of $`I_n`$ is quite complex and strongly differs from the weak-nonlinearity behavior. At very large $`\mathrm{\Phi }_1^{saw}`$, the high-harmonics intensities tends to saturate: $`I_nI_{n,max}`$. At higher electron densities the saturation of $`I_n`$ occurs at larger SAW potentials (see Fig. 5).
## 4 Analytic results for acousto-electric effects in a 2D electron plasma
In this section, we will give some analytic expressions describing various acousto-electric effects in a system with a 2DEG. For simplicity, we will not take into account diffusion, hence assuming $`D_e=0`$ in the formula for the current (5). This assumption is well justified at large electron densities or at low temperatures. It is convenient to start with the Weinreich relation for a 2D system in the nonlinear regime . Assuming that $`n(x_1)`$ and $`j(x_1)`$ are periodic functions we can rewrite Eq. 12 in the form
$`j(x_1)v_s^0en(x_1)=b_0.`$ (26)
Obviously, it is valid in the limit $`K_{eff}^21`$. Using Eq. 26 we can write $`<j>=|e|\mu <nE_x>=\mu <jE_x>/v_s^0`$. From this equation we now get the Weinreich relation:
$`{\displaystyle \frac{<jE_x>}{<j>}}={\displaystyle \frac{v_s^0}{\mu }},`$ (27)
where $`<jE_x>`$ is the dissipation in a SAW. To get Eq. 27, we have neglected the averaged electric fields $`<E_x>=0`$ assuming the case $`V_1=V_2`$.
We now consider Eq. 23 with $`D_e=0`$. At fixed $`\mathrm{\Phi }_1^{saw}`$, there is a critical density, $`N_{crit}`$, for the formation of stripes in a 2DEG. When $`N_s<N_{crit}`$ the plasma is split into electron stripes. If $`N_s>N_{crit}`$, the plasma is continuous but can be strongly modulated in space. For the case $`N_s<N_{crit}`$, the constant $`b_0`$ in Eq. 23 becomes zero and thus Eq. 23 has formally two solutions:
$`f_1(x_1)=const_2{\displaystyle \frac{ϵ_s\mathrm{\Phi }_1^{saw}}{4\pi |e|d}}\mathrm{sin}(qx_1){\displaystyle \frac{v_s^0}{\mu }}{\displaystyle \frac{ϵ_s}{4\pi |e|d}}x;f_2(x_2)=0.`$ (28)
The solution has to be a periodic continuous combination of these two functions. It follows from Eq. 28 that in the limit $`\sigma _0\sigma _m`$ $`N_{crit}\mathrm{\Phi }_1^{saw}ϵ_s/(4\pi |e|d)`$, where $`\sigma _0=|e|\mu N_s`$.
When $`N_s<N_{crit}`$ the plasma turns into stripes, which means that electrons are totally trapped and the local electron velocity in a 2DEG reaches its maximum $`j/(eN_s)=v_s^0`$. From the Weinreich relation, we get $`\mathrm{\Gamma }_1=<jE_x>/I_1=|e|N_s(v_s^0)^2/(\mu I_1)1/I_1`$ .
Eqs. 23 and 24 can be used to find an asymptotic formula for $`\mathrm{\Gamma }_1`$ and $`<j>`$ in the large density limit when $`\sigma _0\sigma _m`$. In the limit $`\sigma _0\sigma _m`$, $`nN_s(ϵ_s\mathrm{\Phi }_1^{saw})/(4\pi |e|d)\mathrm{sin}(qx_1)`$ (see Eq. 28). By using the above and Eqs. 26 and 27, we have in the limit $`\sigma _0\sigma _m`$ and in the region $`\mathrm{\Phi }_1^{saw}<\mathrm{\Phi }_{crit}`$:
$`\mathrm{\Gamma }_1=\mathrm{\Gamma }_{max}{\displaystyle \frac{4\sigma _m}{\sigma _0}}({\displaystyle \frac{\mathrm{\Phi }_{crit}}{\mathrm{\Phi }_1^{saw}}})^2\left(1\sqrt{1[{\displaystyle \frac{\mathrm{\Phi }_1^{saw}}{\mathrm{\Phi }_{crit}}}]^2}\right),`$ (29)
where $`\mathrm{\Gamma }_{max}=qK_{eff}^2/4`$ and $`\mathrm{\Phi }_{crit}=N_s4\pi |e|d/ϵ_s`$. This equation is valid when $`(\mathrm{\Phi }_{crit}\mathrm{\Phi }_1^{saw})/\mathrm{\Phi }_{crit}\sigma _m/\sigma _0`$. Eq. 29 reproduces the numerical data for $`\mathrm{\Gamma }_1(\mathrm{\Phi }_1^{saw})`$ in Fig. 3 at large densities. The asymptotic formula (29) was given before in Ref. without noting the condition $`\sigma _0\sigma _m`$. In the linear regime of interaction $`\mathrm{\Gamma }_1`$ and $`\delta v_1`$ are given by the formula (1).
In the end of this section we consider an asymptotic behavior for the high-harmonic intensities $`I_n`$ in the limit $`I_1\mathrm{}`$. The electron density $`n_s(x_1)`$ at high $`\mathrm{\Phi }_1^{saw}`$ can be written in a parabolic approximation. Then, by calculating the Fourier components $`n_n`$, we find that $`I_{n,max}I_n(\mathrm{\Phi }_1^{saw})^{2/3}I_1^{1/3}`$, where $`n=2,3,\mathrm{}`$.
## 5 Comparison with experimental data
The experiments involving SAW’s were performed on the hybrid semiconductor-$`LiNbO_3`$ structures fabricated by the epitaxial lift-off (ELO) technique developed by Yablonovich et al. . The structures contain a $`12nm`$-thick high-quality $`In_{0.2}Ga_{0.8}As`$ quantum well (QW) embedded in modulation doped $`Al_{0.2}Ga_{0.8}As`$ barriers. In these structures, the thin semiconductor layered system including a QW was tightly bound to the lithium niobate host crystal by the van der Waals forces . The MBE grown quantum well structure is removed from its native $`GaAs`$ substrate by etching an $`AlAs`$ sacrificial layer below the active semiconductor system. The thin ELO film with a thickness of only $`500nm`$ is then transferred onto the host $`LiNbO_3`$ crystal. The parameters for this structure were described already in the beginning of Sec. 1. The geometry of the structure is shown in Fig. 1. For further details related to the fabrication procedure of such quasi-monolithic structures we refer the reader to Refs. \[8-10\]. The experiments were performed for two SAW frequencies $`f=340MHz`$ and $`f=114MHz`$ at room temperature.
The SAW in our experiments can be strong enough to break up an initially homogenous 2D plasma into moving stripes. The transition to the regime of moving electron stripes was directly observed in the experiments on acoustic charge transport (ACT) in samples with specially-designed injection and detection dates . In these experiments the velocity of the ACT-signal first increases with the SAW-intensity and finally saturates at the sound velocity. The latter manifests the formation of stripes. Strongly nonlinear effects are also observed in the attenuation data. The attenuation of a SAW with $`f=114MHz`$ for different intensities is plotted in the insert of Fig. 2 as a function of the transport-gate voltage, which determines the averaged electron density in a 2DEG. At small SAW-intensities, the electronic sound attenuation $`\mathrm{\Gamma }^0`$ as a function of the conductivity $`\sigma _0`$ is described by the well-known linear-theory equation (1) and exhibits a maximum. This linear regime is realized in our experiments at the smallest SAW intensities of about $`12dBm`$ (insert of Fig. 2). It is seen from the insert of Fig. 2 that at high SAW amplitudes the attenuation is strongly suppressed and its maximum is shifted to higher gate bias or conductivity, respectively. The experimental data for a SAW with the frequency $`f=340MHz`$ look qualitatively similar to those for $`f=114MHz`$ and were given earlier in Ref. .
The nonlinear regime of interaction is described by the theory given in Secs. 2, 3, and 4. To quantitatively compare theory and experiment, we now express the SAW potential amplitude $`\mathrm{\Phi }_1^{saw}`$ through the input radio frequency (RF) power $`P`$. The SAW intensity can be written as $`I_1=I_{saw}=2(P/w)10^{IL/10}`$, where the width of the transducer $`w=0.55mm`$. The insertion losses $`(IL)`$ in the transducers were measured to be $`15dB`$. Then, the SAW potential can be written using Eq. 22 as $`\mathrm{\Phi }_1^{saw}=K_{eff}\sqrt{2I_1/(q\sigma _m)}`$, where $`K_{eff}^2=0.015`$ for $`f=114MHz`$ . In Fig. 3 we also show the experimentally measured absorption coefficient at the gate voltage $`V_t=7.5V`$. This voltage corresponds to the maximal attenuation for the smallest RF power, $`12dBm`$ (see insert of Fig. 2). From the linear theory we find that the absorption coefficient is maximal at $`N_s=0.910^{10}1/cm^2`$. One can see from Fig. 3 that the experimentally measured function $`\mathrm{\Gamma }_1(\mathrm{\Phi }_1^{saw})`$ for $`V_t=7.5V`$ is in a very good agreement with the calculated one for $`N_s=0.810^{10}cm^2`$. Here we did not use any fitting parameters. This quantitative agreement becomes possible if we account for the diffusion coefficient. The maximal absorption coefficient as calculated from the linear theory is $`\mathrm{\Gamma }_{max}^{diff}=7.6cm^1`$. Without diffusion this value is about $`\mathrm{\Gamma }_{max}=14.3cm^1`$. Thus, the diffusion strongly suppresses the SAW absorption.
At fixed SAW power and a sufficiently small density $`N_s`$, the 2DEG is divided into stripes and $`\mathrm{\Gamma }_1`$ increases with increasing the gate voltage. In our theory, $`\mathrm{\Gamma }_1N_s`$ in the regime of stripes, which explains the increase of the attenuation at small gate voltages in the insert of Fig. 2. At a sufficiently large gate voltage, the absorption coefficient as a function of $`V_t`$ starts to decrease because of the screening effect in a high-density 2DEG modulated by a SAW. The interplay of these two effects leads to the shift of the maximum of the function $`\mathrm{\Gamma }_1(V_t)`$ shown in the insert of Fig. 2.
The change of the SAW-velocity due to the electron plasma is shown in the insert of Fig. 4. With increasing RF power the curves in Fig. 4 are again shifted towards larger electron conductivity which can be understood in terms of screening. With increasing SAW intensity the electron plasma is strongly modulated or even split into stripes and the screening of piezoelectric fields by electrons becomes not so effective. Thus, the shift of the SAW-velocity due to electrons decreases with increasing the SAW intensity. The experimentally observed shift is in qualitative agreement with our modeling shown in Fig. 4. However, our theory does not reproduce the character of $`\delta v_1(V_t)`$ in the region of small electron densities. Likely, the measurement of the SAW velocity is not so sensitive to a low-density electron system in comparison with the attenuation method.
In Fig. 6, we show the quantity $`\mathrm{\Gamma }_1P/I_{ae}`$ as a function of the RF power to verify the Weinreich relation. The acousto-electric current $`<j>=I_{ae}(P)`$ was measured in a ”short circuit geometry”, where the Ohmic contacts are directly (without resistor) connected to the current measurement instrument. For more details on acousto-electric current measurements we refer to Ref. . We see from Fig. 6 that the ratio $`\mathrm{\Gamma }_1P/I_{ae}`$ has a weak power-dependence. Thus, our experimental data are well described by the Weinreich relation. Slight deviations from the Weinreich relation seen in Fig. 6 can come from the density-dependence of the mobility $`\mu (N_s)`$, that is expected to be relatively weak at room temperature. The reason is that the main electronic scattering mechanism at high temperatures is due to acoustic phonons and is relatively insensitive to the 2D density.
## 6 Conclusions
The theoretical results obtained in Sec. 2, can also be applied to study dynamics of SAW’s at large distances in a long sample, where the contribution of high harmonics can be very important . This long-distance transformation to high harmonics was studied experimentally for SAW’s interacting with a 3D electron gas of a semiconductor on a piezocrystal . In the presence of a dc voltage applied to the crystal, it is possibly to expect the appearance of nonlinear waves with a stationary profile or with a stationary energy flow . In a wave with a stationary energy flow the wave shape is periodically changed in space . Another scenario can relate to chaotic dynamics in an acousto-electric system . Eq. 14 can be used to numerically model these phenomena in 2D electron systems at long distances.
To conclude, we have studied strongly nonlinear acousto-electric phenomena caused by the interaction between a SAW and a two-dimensional electron system. In the experimental measurements performed on hybrid semiconductor-piezocrystal structures the SAW attenuation, the SAW velocity change, and the acousto-electric current are strongly modified in the nonlinear regime due to the formation of moving electron stripes. By using a coupled-amplitude method we have modeled the decay and transformation of SAW’s in the nonlinear regime. Using our theoretical results, we could explain our experimental findings and distinguished between different regimes of the nonlinear acousto-electric interaction at large SAW intensities.
## Acknowledgements
We would like to thank D. Bernklau and H. Riechert for the fabrication of the excellent MBE material, W. Ruile for the strong support from the SAW device side, and A. V. Chaplik and M. K. Balakirev for helpful discussions. We gratefully acknowledge financial support by the Volkswagen-Stiftung and by the Russian Foundation for Basic Research (grants 99-02-17019 and 99-02-17127).
## Appendix 1
Here we intend to briefly discuss the electrostatics of the hybrid structure. The spacing between the 2DEG and the top metal gate in the fabricated structures is much larger than the distance from the 2DEG to the $`AlGaAsLiNbO_3`$ interface. Thus, to model the screening effects, we will assume that the 2DEG is located right on the $`AlGaAsLiNbO_3`$ interface. It is convenient to solve this problem by a Fourier transform in terms of $`x`$ and by remaining the vertical coordinate $`z`$. The relation between the Fourier components of the electrostatic potential $`W^{ind}[z;k]`$ induced by 2D electrons and the 2DEG density $`n_k`$ is found from the Poisson equation and from the corresponding boundary conditions,
$`W^{ind}[z;k]={\displaystyle \frac{2\pi en_k}{|k|ϵ_{eff}(k)}}G(z),`$ (30)
where $`ϵ_{eff}(k)=(ϵ_p+ϵ_scoth(|k|d))/2`$. The function $`G(z)=\frac{\mathrm{sinh}|k|z}{\mathrm{sinh}|k|d}`$, when $`0<z<d`$, and $`G(z)=e^{|k|(zd)}`$ for $`d<z`$ .
Taking the electron density in the form of $`n_s(x,x_1,t)=n_0(x)+_{n=1,2,\mathrm{}}n_n(x)e^{iq_nx_1}+c.c.`$, we can write for the induced electrostatic potential: $`\varphi ^{ind}(x,x_1,z,t)=\varphi _0^{ind}(x,z)+_{n=1,2,\mathrm{}}\varphi _n^{ind}(x,z)e^{iq_nx_1}+c.c.`$, where $`q_n=nq`$. Using Eq. 30, we find $`\varphi _n^{ind}(x,z)=2\pi en_n(x)G(z)/\left(q_nϵ_{eff}(q_n)\right)+\delta \varphi _n(x,z)`$, where $`n=1,2,\mathrm{}`$. The correction $`\delta \varphi _n(x,z)da_n(x)/dx\delta q_nK_{eff}^2`$ (see Eq. 15 for $`\delta q_n`$) and is small compared to $`W^{ind}[z;q_n]`$. Thus, regarding $`n_n(x)`$ as the constants, we obtain at $`z=d`$: $`\varphi _n^{ind}(x,d)2\pi en_n(x)/[q_nϵ_{eff}(q_n)]`$ and $`E_{nx}(x,d)iq_n\varphi _n^{ind}(x,d)`$. In other words, we consider the envelope functions $`n_n(x)`$ as constants and solve Poisson’s equation in terms of the ”fast” variable $`x_1`$. Again, it is valid in the limit $`\mathrm{\Gamma }^0/qK_{eff}^21`$.
As an example, we now calculate the correction $`\delta \varphi _n`$ in the linear regime of interaction, when $`n(x,x_1,t)=N_s+n_1(x)e^{iqx_1}+n_1^{}(x)e^{iqx_1}`$ with $`n_1(x)=\stackrel{~}{n}_1e^{\mathrm{\Gamma }^0x/2iq\delta v_s^0/v_s^0x}`$. From the Poisson equation we find $`\delta \varphi _1(x,z)(\mathrm{\Gamma }^0/2+iq\delta v_s/v_s^0)n_1(x)F(z)`$, where $`F(z)1`$ for $`z1/q`$.
## Appendix 2
In order to solve a system of nonlinear equations (2-5) we will use some of results from a linear-response theory . In a linear theory the total electrostatic potential and the 2D density can be written as $`W[x,z,t]=W[z,k]e^{ikxi\omega t}`$ and $`n(x,t)=n_ke^{ikxi\omega t}`$, respectively. Here $`\omega `$ is the SAW frequency. It is convenient to introduce a quantity $`\mathrm{\Pi }_k`$ by means of the relation $`en_k=\mathrm{\Pi }_kW[d;k]`$, where $`W[d,k]`$ is the Fourier component of the electrostatic potential at $`z=d`$. Eqs. 2,3 are now written as
$`\omega ^2\rho u_i+c_{iklm}_m_ku_l+p_{lik}_l_kW[x,z,t]=0,`$ (31)
$`4\pi p_{ikl}_i_lu_k+\left(ϵ_i_i4\pi \mathrm{\Pi }_k\delta (x_3d)\right)W[x,z,t]=0.`$ (32)
Above equations should be solved together with the necessary boundary conditions considered in Sec. 1. Then, we rewrite Eqs. 31,32 in the form
$`\widehat{L}_{lin}𝐀_{lin}=0,`$ (33)
Here $`\widehat{L}_{lin}`$ is a linear operator and $`𝐀_{lin}=(𝐮[x,z,t],W[x,z,t])=(𝐔^0[z;k],W^0[z;k])e^{ikxi\omega t}`$. It follows from the boundary conditions that $`k=q+\delta q_{lin}`$ , where $`q=\omega /v_s^0`$ and
$`\delta q_{lin}={\displaystyle \frac{qK_{eff}^2}{2}}{\displaystyle \frac{\mathrm{\Pi }_q/\mathrm{\Pi }_q^0}{1+\mathrm{\Pi }_q/\mathrm{\Pi }_q^0}},\mathrm{\Pi }_q^0={\displaystyle \frac{qϵ_{eff}(q)}{2\pi }}.`$ (34)
Now a solution of Eq. 33 can be written as $`𝐀_{lin}=𝐀^0[z;q+\delta q_{lin}]f_0(x)e^{iqx_1}`$, where $`f_0(x)=e^{i\delta q_{lin}x}`$ and $`x_1=xv_s^0t`$. $`𝐀^0[z;q+\delta q_{lin}]`$ is a vector, that can be found from the matrix given by the boundary conditions .
Eq. 33 contains first and second spatial derivatives and can be written as $`\widehat{L}_{lin}𝐀_{lin}=e^{iqx_1}[f_0(x)\widehat{L}_0(q)𝐀^0+f_0^{}(x)\widehat{L}_1(q)𝐀^0+f_0^{\prime \prime }(x)\widehat{L}_2(q)𝐀^0]=0`$, where $`f^{}=df/dx`$. Neglecting the second derivative $`f_0^{\prime \prime }(x)`$, that is $`\delta q_{lin}^2K_{eff}^4`$, we have
$`\widehat{L}_0(q)𝐀^0+i\delta q_{lin}\widehat{L}_1(q)𝐀^0=0.`$ (35)
Now we turn to the nonlinear theory, where the equation $`\widehat{L}𝐀=0`$ is also valid. The operator $`\widehat{L}`$ is determined by the equations similar to Eqs. 31-33 but with nonlinear quantity $`\mathrm{\Pi }_n(x)`$, that is defined by $`\mathrm{\Pi }_n(x)=en_n(x)/\varphi _n(x,d)`$. In the nonlinear case the vector $`𝐀=𝐀_0(x,z)+_{n=1,2,\mathrm{}}a_n(x)𝐀_n(x,z)e^{iq_nx_1}+c.c.`$. Each term in the equation $`\widehat{L}𝐀=0`$ should be zero, and so $`\widehat{L}a_n(x)𝐀_n(x,z)e^{iq_nx_1}e^{iq_nx_1}[a_n(x)\widehat{L}_0(q_n)𝐀_n+a_n^{}(x)\widehat{L}_1(q_n)𝐀_n]=0`$. To get the latter equation, we have neglected $`d\delta q_n/dxK_{eff}^4`$ and $`d^2a_n/dx^2K_{eff}^4`$. We can solve the equation $`\widehat{L}𝐀=0`$ if we choose $`𝐀_n(x,z)=𝐀^0[z,q_n+\delta q_n(x)]`$, where the vector $`𝐀^0[z;q]`$ is defined above in the linear theory. Using Eq. 35 we get
$`{\displaystyle \frac{da_n(x)}{dx}}=i\delta q_n(x)a_n(x),`$ (36)
where $`\delta q_n(x)`$ is given by the equation for $`\delta q_{lin}`$ (Eq. 34) with corrections $`\mathrm{\Pi }_q\mathrm{\Pi }_n(x)`$ and $`\mathrm{\Pi }_q^0\mathrm{\Pi }_{q_n}^0=\mathrm{\Pi }_n^0`$. Eq. 36 is used in Secs. 1,2 to describe the acousto-electric phenomena in a 2DEG. Using Eqs. 8 and 11 and the results of Appendix 1, the denominator in $`\delta q_n(x)`$ (see Eq. 34) is rewritten as $`1+\mathrm{\Pi }_n(x)/\mathrm{\Pi }_n^0=\mathrm{\Pi }_n(x)\varphi _n^{saw}(x)/en_n(x)`$. Then, by using Eq. 34 and the conservation-of-charge equation $`j_n(x)=v_s^0en_n(x)`$, we obtain
$`\delta q_n(x)={\displaystyle \frac{K_{eff}^2(q_n)}{2}}{\displaystyle \frac{2\pi }{ϵ_{eff}(q_n)}}{\displaystyle \frac{en_n(x)}{\varphi _n^{saw}(x)}}=i{\displaystyle \frac{K_{eff}^2(q_n)}{2\mathrm{\Pi }_n^0v_s^0}}{\displaystyle \frac{j_n(x)E_n^{saw}}{|\varphi _n^{saw}(x)|^2}}.`$ (37)
We use this equation in Sec. 2 (see Eq. 15). To obtain above results, we have neglected the terms like $`a_n^{\prime \prime }(x)`$ and $`\delta q_n^{}(x)`$ assuming that $`K_{eff}^2`$ is a small parameter.
## Figure captions
Fig. 1. The cross section of a hybrid semiconductor-piezocrystal structure. An epitaxial lift-off film has a thickness $`0.5\mu m`$. The Ohmic contacts are formed to a 2D electron gas. The transport gate with applied voltage $`V_t`$ is used to control the conductivity of the electron plasma. A high-frequency (RF) signal is applied to the metal interdigital transducer $`IDT1`$ in order to generate surface acoustic waves. A surface acoustic wave propagates though a sample and is detected by the transducer $`IDT2`$.
Fig. 2. The calculated absorption coefficient of the first harmonic $`\mathrm{\Gamma }_1`$ as a function of the carrier density $`N_s`$ for various potential amplitudes $`\mathrm{\Phi }_1^{saw}`$. $`\lambda =33\mu m`$, $`\mu =5000cm^2V/s`$, and $`T=300K`$. Insert: The measured attenuation of a SAW as a function of the gate voltage for different high-frequency (RF) powers applied to the IDT1; $`f=114MHz`$.
Fig. 3. The calculated absorption coefficient of the first harmonic $`\mathrm{\Gamma }_1`$ as a function of the potential amplitude $`\mathrm{\Phi }_1^{saw}`$ induced by a SAW for various fixed densities $`N_s`$. The parameters are similar to those in Fig. 2. The dots show the experimentally measured absorption coefficient $`\mathrm{\Gamma }_1`$ at the gate voltage $`7.5V`$. This voltage corresponds to the maximal attenuation at the smallest RF power. In the inset we plot the calculated local carrier concentration $`n`$ as a function of the in-plane coordinate $`x_1`$ for different total carrier concentration $`N_s`$. The numbers attached to the plots correspond to $`N_s`$ in units of $`10^{10}cm^2`$.
Fig. 4. The calculated SAW-velocity change $`\delta v_1`$ as a function of the electron density for various potential amplitudes $`\mathrm{\Phi }_1^{saw}`$. The parameters are similar to those in Fig. 2. Insert: The measured velocity change of a SAW as a function of the gate voltage for different high-frequency (RF) powers applied to the IDT1; $`f=114MHz`$.
Fig. 5. The calculated intensities of higher harmonics with $`n=2,3`$ and $`4`$ as functions of the potential amplitude $`\mathrm{\Phi }_1^{saw}`$ for two electron densities $`N_s=10^{10}cm^2`$ (upper part) and $`N_s=10^{11}cm^2`$ (lower part). The parameters are similar to those in Fig. 2.
Fig. 6. The measured ratio $`\mathrm{\Gamma }_1P/I_{ae}`$ as a function of the high-frequency (RF) power $`P`$ for $`f=340MHz`$; $`T=300K`$. The acousto-electric current $`I_{ae}`$ was measured in its maximum.
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# Coarse Obstructions to Positive Scalar Curvature Metrics in Noncompact Arithmetic Manifolds
## Abstract.
Block and Weinberger show that an arithmetic manifold can be endowed with a positive scalar curvature metric if and only if its $``$-rank exceeds $`2`$. We show in this article that these metrics are never in the same coarse class as the natural metric inherited from the base Lie group. Furthering the coarse $`C^{}`$-algebraic methods of Roe, we find a nonzero Dirac obstruction in the $`K`$-theory of a particular operator algebra which encodes information about the quasi-isometry type of the manifold as well as its local geometry.
I. Introduction
In the course of showing that no manifold of non-positive sectional curvature can be endowed with a metric of positive scalar curvature, Gromov and Lawson were led to consider what we would now call restrictions on the coarse equivalence type of complete noncompact manifolds of such positively curved metrics. In particular, they showed that such metrics cannot exist in manifolds for which there exists a degree one proper Lipschitz map from the universal cover to $`^n`$, now understood to be essentially a coarse condition. Block and Weinberger investigate the situation in which no coarse conditions are imposed upon the complete metric, focusing on quotients $`\mathrm{\Gamma }\backslash G/K`$ of symmetric spaces associated to a lattice $`\mathrm{\Gamma }`$ in an irreducible semisimple Lie group $`G`$. They show that the space $`M=\mathrm{\Gamma }\backslash G/K`$ can be given a complete metric of uniformly positive scalar curvature $`\kappa \epsilon >0`$ if and only if $`\mathrm{\Gamma }`$ is an arithmetic group of $``$-rank exceeding $`2`$.
Note that the theorem of Gromov and Lawson mentioned above establishes this theorem in the case of $`\text{rank}_{}\mathrm{\Gamma }=0`$. In the higher rank cases, for which the resulting quotient space is noncompact, the metrics constructed by Block and Weinberger are however wildly different in the large when compared to the natural one on $`M`$ inherited from the base Lie group $`G`$. In fact, their examples are all coarse quasi-isometric to rays. Their theory evokes a natural question: Can the metric be chosen so that it is simultaneously uniformly positively curved and coarsely equivalent to the natural metric induced by $`G`$?
One of the important developments in analyzing positive scalar curvature in the context of noncompact manifolds, especially when restricted to the coarse quasi-isometry type, is introduced by Roe , , who considers a higher index, analogous to the Novikov higher signature, that lives naturally in the $`K`$-theory of the $`C^{}`$-algebra $`C^{}(M)`$ of operators on $`M`$ with finite propagation speed. He describes a map from the $`K`$-theory group $`K_{}(C^{}(M))`$ to the $`K`$-homology $`K_{}(\nu M)`$ of the Higson corona space which admits a dual transgression map $`H^{}(\nu M)HX^{}(M)`$. If the Dirac operator on $`M`$ is invertible, then the image of its index in $`K_{}(\nu M)`$ vanishes, leading to vanishing theorems for the index paired with coarse classes from the transgression of $`\nu M`$. Roe’s construction is used to show that a metric on a noncompact manifold cannot be uniformly positively curved if the Higson corona of the manifold contains an essential $`(n1)`$-sphere. Such spaces are called ultraspherical manifolds.
The usual Roe algebra, however, is unsuited to provide information about the existence of positive scalar curvature metrics that exist on arithmetic manifolds, because their coronae are too anemic. For example, the space at infinity of a product of punctured two-dimensional tori is a simplex and therefore contractible. As a coarse object, the $`K`$-theory of the Roe algebra associated to this multi-product space can be identified with $`K_{}(C^{}(_0^n))`$. Yet Higson, Roe and Yu have shown that the Euclidean cone $`cP`$ on a single simplex $`P`$ must satisfy $`K_{}(C^{}(cP))=0`$. Since the Euclidean hyperoctant $`_0^n`$ is simply the cone on an $`(n1)`$-simplex, we find that $`K_{}(C^{}(_0^n))`$ is the trivial group and hence no obstructions are detectable. Even by considering the fundamental group of the manifold by tensoring the Roe algebra with $`C^{}\pi _1(M)`$ is this detection process unfruitful, since the $`K`$-theory group $`K_{}(C^{}(M)C^{}\pi _1(M))`$ vanishes as well. What seems to be critical is how different elements of the fundamental group at infinity can be localized to different parts of the space at infinity.
In this article, we shall provide coarse indicial obstructions in the following noncompact manifolds: a finite product of punctured two-dimensional tori, a finite product of hyperbolic manifolds, the double quotient space $`\mathrm{SL}_n()\backslash \mathrm{SL}_n()/\mathrm{SO}_n()`$ of unit volume tori, and more generally the double quotient space $`\mathrm{\Gamma }\backslash G/K`$, where $`G`$ is an irreducible semisimple Lie group, $`K`$ its maximal compact subgroup and $`\mathrm{\Gamma }`$ an arithmetic subgroup of $`G`$. Note that the first two do not correspond to irreducible quotients, but an analysis of these spaces gives us the proper insight to attack the more general cases. A further research project will analyze this problem without the irreducibility assumption. The key feature in these particular manifolds $`M`$ is that they contain hypersurfaces $`V`$ that are coarsely equivalent to a product $`E\times U`$ of Euclidean space $`E`$ with some iterated circle bundle $`U`$ (i.e. a torus, Heisenberg group, or more generally a group of unipotent matrices). Moreover such a hypersurface decomposes the manifold $`M`$ into a coarsely excisive pair $`(A,B)`$ for which $`AB=M`$ and $`AB=V`$. A generalized form of the Mayer-Vietoris sequence constructed by Higson, Roe and Yu provides the following:
$$\mathrm{}K_{}(C_G^{}(A))K_{}(C_G^{}(B))K_{}(C_G^{}(M))K_1(C_G^{}(V))\mathrm{}$$
The boundary map $`:K_{}(C_G^{}(M))K_1(C_G^{}(E\times U))`$ sends $`\text{ Ind }_M(D)`$, the index of the spinor Dirac bundle on the universal cover lifted from that on $`M`$, to $`\text{ Ind }_{E\times U}(D)`$. To see that these indices are indeed nonzero, we note that there is a boundary map $`K_1(C_G^{}(E\times U))K_{\mathrm{dim}E}(C_G^{}(\times U))`$, which sends index to index. We show that the index of the Dirac operator in the latter group, however, is nonzero by noting that the Gromov-Lawson-Rosenberg conjecture is true for nilpotent groups and hence provides an appropriate nonzero obstruction.
I would like to thank Alex Eskin, Benson Farb, Nigel Higson, Thomas Nevins, Mel Rothenberg, John Roe, Stephan Stolz and Guoliang Yu for very useful conversations. In particular, I would like to acknowlege the role of my advisor Shmuel Weinberger in pointing out the strength of certain tools in the realization of these theorems.
II. The Generalized Roe Algebra
The coarse category is defined to contain metric spaces as its objects and maps $`f:(X,d_X)(Y,d_Y)`$ between metric spaces as its morphisms satisfying the following expansion and properness conditions: (a) for each $`R>0`$ there is a corresponding $`S>0`$ such that, if $`d_X(x_1,x_2)R`$ in $`X`$, then $`d_Y(f(x_1),f(x_2))S`$, (b) the inverse image $`f^1(B)`$ under $`f`$ of each bounded set $`BY`$ is also bounded in $`X`$. Such a function will be designated a coarse map, and two coarse maps $`f,g:XY`$ are said to be coarsely equivalent if their mutual distance of separation $`d_Y(f(x),g(x))`$ is uniformly bounded in $`x`$. Naturally two metric spaces are coarsely equivalent if there exist maps from one to the other whose compositions are coarsely equivalent to the appropriate identity maps. Two metrics $`g_1`$ and $`g_2`$ on the same space $`M`$ are said to be coarsely equivalent if $`(M,g_1)`$ and $`(M,g_2)`$ are coarsely equivalent metric spaces.
Following Roe , we recall a Hilbert space $`H`$ is an $`M`$-module for a manifold $`M`$ if there is a representation of $`C_0(M)`$ on $`H`$, that is, a $`C^{}`$-homomorphism $`C_0(M)B(H)`$. We will say that an operator $`T:HH`$ is locally compact if, for all $`\phi C_0(M)`$, the operators $`T\phi `$ and $`\phi T`$ are compact on $`H`$. We define the support of $`\phi `$ in an $`M`$-module $`H`$ to be the smallest closed set $`KM`$ such that, if $`fC_0(M)`$ and $`f\phi 0`$, then $`f|_K`$ is not identically zero. Consider the $`\stackrel{~}{M}`$-module $`H=L^2(\stackrel{~}{M})`$, where $`\stackrel{~}{M}`$ is the universal cover of $`M`$ endowed with the appropriate metric lifted from the base space. Let $`\pi :\stackrel{~}{M}M`$ be the usual projection map and for any $`\phi ,\psi C_0(\stackrel{~}{M})`$ consider the collection $`\mathrm{\Gamma }(\phi ,\psi )`$ of paths $`\gamma :[0,1]\stackrel{~}{M}`$ in $`\stackrel{~}{M}`$ originating in $`\text{Supp}(\phi )`$ and ending in $`\text{Supp}(\psi )`$. Denote by $`L[\gamma ]`$, for $`\gamma \mathrm{\Gamma }(\phi ,\psi )`$, the maximum distance of any two points on the projection of the curve $`\gamma `$ in $`M`$ by $`\pi `$, i.e. $`L[\gamma ]=sup_{x,y[0,1]}d(\pi \gamma (x),\pi \gamma (y))`$.
Definition: Let $`M`$ be a manifold with universal cover $`\stackrel{~}{M}`$. We say that an operator $`T`$ on $`L^2(\stackrel{~}{M})`$ has generalized finite propagation if there is a constant $`R>0`$ such that $`\phi T\psi `$ is identically zero in $`B(H)`$ whenever $`\phi ,\psi C_0(\stackrel{~}{M})`$ satisfies
$$\underset{\gamma \mathrm{\Gamma }(\phi ,\psi )}{inf}L[\gamma ]>R.$$
The infimum of all such $`R`$ will be the generalized propagation speed of the operator $`T`$. If $`G=\pi _1(M)`$ is the fundamental group of $`M`$, we denote by $`D_G^{}(M)`$ to be the norm closure of the $`C^{}`$-algebra of all locally compact, $`G`$-equivariant, generalized finite propagation operators on $`H`$.
Let $`M`$ be a manifold and $`\stackrel{~}{M}`$ its universal cover. Let $`T:HH`$ be an operator on $`H=L^2(\stackrel{~}{M})`$. Consider the subset $`Q\stackrel{~}{M}\times \stackrel{~}{M}`$ of pairs $`(m,m^{})`$ for which there exist functions $`\phi ,\psi C_0(\stackrel{~}{M})`$ such that $`\phi (m)0`$, $`\psi (m^{})0`$ and $`\phi T\psi `$ does not identically vanish. We will say that the support of $`T`$ is the complement in $`\stackrel{~}{M}\times \stackrel{~}{M}`$ of $`Q`$. For such two points $`m,m^{}\stackrel{~}{M}`$, let $`\gamma _{mm^{}}:[\mathrm{\hspace{0.17em}0},1]\stackrel{~}{M}`$ be the path of least length joining $`m`$ and $`m^{}`$ in $`\stackrel{~}{M}`$. We consider the projection of this path into $`M`$ by $`\pi `$ and take the greatest distance between two points on this projected path. Then it is easy to see that an operator $`T`$ has generalized finite propagation, as previously defined, if
$$\underset{m,m^{}}{sup}\underset{x,y[0,1]}{sup}d(\pi \gamma _{mm^{}}(x),\pi \gamma _{mm^{}}(y))<\mathrm{}.$$
Definition: Consider the norm closure $`I`$ of the ideal in $`D_G^{}(M)`$ generated by operators $`T`$ whose matrix representation, parametrized by $`\stackrel{~}{M}\times \stackrel{~}{M}`$, satisfies the condition that $`(\pi \times \pi )(\text{Supp}T)`$ is bounded in $`M\times M`$. Then the generalized Roe algebra, denoted by $`C_G^{}(M)`$, is obtained as the quotient $`D_G^{}(M)/I`$. Two operators in $`D_G^{}(M)`$ belong to the same class in $`C_G^{}(M)`$ if their nonzero entries differ on at most a bounded set when viewed from the perspective of the base space.
Examples:
(1) Let $`T:L^2()L^2()`$ be operator on $`L^2`$-functions on the real line given by $`(Tg)(x)=g(x+1)`$ for all $`gL^2()`$ and $`x`$. Then for any $`\phi ,\psi C_0()`$, $`(\phi T\psi )g(x)=\phi (x+1)g(x+1)\psi (x)`$. If $`\phi `$ is supported at $`m=1`$ and $`\psi `$ is supported at $`m^{}=0`$, then $`(\phi T\psi )g`$ is nonzero for any $`g`$ supported at $`x=1`$. Hence $`(0,1)\text{Supp}T`$. It is easy to see that $`(m,m^{})\text{Supp}T`$ if and only if $`m^{}m=1`$. The propagation speed of $`T`$ is $`1`$. If we write $`T`$ as a matrix parametrized by $`\times `$, all the nonzero entries will lie at distance one from the diagonal.
(2) Let $`M`$ be the cylinder $`S^1\times `$ with its universal cover $`\stackrel{~}{M}=^2`$. An operator in the algebra $`D_G^{}(M)`$ will be some $`T:HH`$ on $`L^2(^2)`$, which is of finite propagation speed (in the usual sense) in the direction projecting down to the noncompact direction in $`M`$, but has no such condition in the orthogonal direction corresponding to the compact direction of $`M`$. In this direction, however, the operator is controlled by the condition that it be $``$-equivariant. It is apparent that the operator, when restricted to individual fibers, has finite propagation speed, there is no requirement that the speed to be uniformly bounded across all fibers.
(3) Let $`M=\stackrel{}{^n}`$, $`n3`$, the once-punctured real projective space, expressible as the quotient $`(S^{n1}\times )/_2`$. Certainly $`M`$ is coarsely equivalent to the ray $`[\mathrm{\hspace{0.17em}0},\mathrm{})`$ and is covered by the space $`\stackrel{~}{M}=S^{n1}\times `$, where the points $`(s,r)`$ and $`(s,r)`$ are identified by the projection map to $`M`$. Let $`T:L^2(\stackrel{~}{M})L^2(\stackrel{~}{M})`$ be given by the reflection $`(Tf)(s,r)=f(s,r)`$. Consider $`\phi _i,\psi _iC_0(\stackrel{~}{M})`$ compactly supported on $`S^{n1}\times [i1,i]`$ and $`S^{n1}\times [i,i+1]`$, respectively. Notice that $`\phi T\psi `$ will never be identically zero, and yet the length $`L_i[\gamma ]`$ associated to $`\phi _i`$ and $`\psi _i`$ will always be at least $`i`$. Hence the operator $`T`$ is not of generalized finite propagation speed and therefore not an element of the generalized Roe algebra $`C_G^{}(M)`$.
Lemma 1: Let $`D`$ a generalized elliptic operator in $`L^2(M,S)`$. Suppose that $`\stackrel{~}{D}`$ is the lifted operator on $`\stackrel{~}{M}`$. If $`\mathrm{\Phi }:`$ is compactly supported, then $`\mathrm{\Phi }(\stackrel{~}{D})`$ lies in the generalized Roe algebra $`C_G^{}(M)`$.
Proof: (cf. , ) Suppose that $`\mathrm{\Phi }`$ has compactly supported Fourier transform and denote by $`\widehat{\mathrm{\Phi }}`$ the Fourier transform of $`\mathrm{\Phi }`$. We may write
$$\mathrm{\Phi }(\stackrel{~}{D})=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}\widehat{\mathrm{\Phi }}(t)e^{it\stackrel{~}{D}}𝑑t.$$
It is known that $`e^{it\stackrel{~}{D}}`$ has finite propagation speed, and since $`\widehat{\mathrm{\Phi }}`$ is compactly supported, the integral is defined and has a generalized propagation bound. Moreover, by construction $`\stackrel{~}{D}`$ is $`\pi _1(M)`$-equivariant. So $`\mathrm{\Phi }(\stackrel{~}{D})`$ is $`\pi _1(M)`$-equivariant as well. Therefore if $`\widehat{\mathrm{\Phi }}`$ is compactly supported, then $`\mathrm{\Phi }(\stackrel{~}{D})`$ lies in $`D_G^{}(M)`$ and passes to an element of the quotient $`C_G^{}(M)`$. However, functions with compactly supported Fourier transform form a dense set in $`C_0()`$ and the functional calculus map $`ff(\stackrel{~}{D})`$ is continuous, so the result holds for all $`\mathrm{\Phi }C_0()`$. $`\mathrm{}`$
Let $`\chi :`$ be a chopping function on $``$, i.e. an odd continuous function with the property that $`\chi (x)\pm 1`$ as $`x\pm \mathrm{}`$. In addition, denote by $`B_G^{}(M)`$ the multiplier algebra of $`C_G^{}(M)`$, that is, the collection of all operators $`S`$ such that $`ST`$ and $`TS`$ belong to $`C_G^{}(M)`$ for all $`TC_G^{}(M)`$. Then $`B_G^{}(M)`$ contains $`C_G^{}(M)`$ as an ideal. If $`D`$ is a generalized elliptic operator on $`M`$ and $`\stackrel{~}{D}`$ its lift to $`\stackrel{~}{M}`$, then $`\chi (\stackrel{~}{D})`$ belongs to $`B_G^{}(M)`$. In addition, since $`\chi ^21C_0()`$, we have $`\chi (\stackrel{~}{D})^21C_G^{}(M)`$. Moreover, since the $`_2`$-grading renders the decompositions
$$\chi (\stackrel{~}{D})=\left(\begin{array}{cc}0& \chi (\stackrel{~}{D})_{}\\ \chi (\stackrel{~}{D})_+& 0\end{array}\right),\epsilon =\left(\begin{array}{cc}\hfill 1& \hfill 0\\ \hfill 0& \hfill 1\end{array}\right),$$
it follows that $`\epsilon \chi (\stackrel{~}{D})+\chi (\stackrel{~}{D})\epsilon =0`$. By the discussion in , it follows that $`F=\chi (\stackrel{~}{D})`$ is a Fredholm operator and admits an index $`\text{ Ind }FK_0(C_G^{}(M))`$. In addition, any two chopping functions $`\chi _1`$ and $`\chi _2`$ differ by an element of $`C_0()`$. By the lemma above, we have $`\chi _1(\stackrel{~}{D})\chi _2(\stackrel{~}{D})C_G^{}(M)`$, so they define the same elements of $`K`$-theory. The common value for $`\text{ Ind }F`$ is denoted $`\text{ Ind }(D)`$ and called the generalized coarse index of $`D`$. We write $`C_G^{}(M)`$ and $`\text{ Ind }(D)`$ instead of $`C_G^{}(\stackrel{~}{M})`$ and $`\text{ Ind }(\stackrel{~}{D})`$ to indicate that the construction is initiated by a generalized Dirac operator on the base space. The following statements are standard results of index theory; one may consult and for the essentially identical proof in the nonequivariant case.
Proposition 1: Let $`D`$ be a generalized elliptic operator in $`L^2(M,S)`$. If $`0`$ does not belong to the spectrum of $`\stackrel{~}{D}`$, then the generalized coarse index $`\text{ Ind }D`$ vanishes in $`K_0(C_G^{}(M))`$.
Proposition 2: Let $`\stackrel{~}{D}`$ the lift of a generalized elliptic operator in $`L^2(\stackrel{~}{M},S)`$. In the ungraded case, if there is a gap in the spectrum of $`\stackrel{~}{D}`$, then the index $`\text{ Ind }D`$ vanishes in $`K_1(C_G^{}(M))`$.
Corollary: Let $`M`$ be a complete spin manifold. If $`M`$ has a metric of uniformly positive scalar curvature in some coarse class, then the generalized coarse index of the spinor Dirac operator vanishes.
We now embark on the task of computing the $`K`$-theory of this algebra and of coarse indices.
Let $`(M,d)`$ be a proper metric space. For any subset $`UM`$ and $`R>0`$, we denote by $`\text{Pen}(U,R)`$ the open neighborhood of $`U`$ consisting of points $`xM`$ for which $`d(x,U)<R`$. Let $`A`$ and $`B`$ be closed subspaces of $`M`$ with $`M=AB`$. We then say that the decomposition $`(A,B)`$ is a coarsely excisive pair if for each $`R>0`$ there is an $`S>0`$ such that
$$\text{Pen}(A,R)\text{Pen}(B,R)\text{Pen}(AB,S).$$
We wish to analyze this decomposition in the following context.
Given general $`C^{}`$-algebras $`𝒜`$, $``$ and $``$ for which $`=𝒜+`$, we have the Mayer-Vietoris sequence
$$\mathrm{}K_{j+1}()K_j(𝒜)K_j(𝒜)K_j()K_j()\mathrm{}$$
The standard proof for the existence of such a sequence is developed from the isomorphism $`K_{}(𝒯)K_1()`$, where $`𝒯`$ is the suspension of $``$. A short discussion of this construction is given in . We are in particular interested in exploiting the boundary map $`:K_j()K_{j1}(𝒜)`$ to transfer information about the index of the Dirac operator on a complete noncompact manifold $`M`$ to information about that on some hypersurface $`V`$. For our purposes, we wish to set $``$ to be the generalized Roe algebra $`C_G^{}(M)`$ on $`M`$, while $`𝒜`$ and $``$ represent analogous operator algebras on closed subsets $`A`$ and $`B`$, where $`(A,B)`$ form a coarsely excisive decomposition of $`M`$. To construct the boundary map in question, we require a few technical lemmas and notion of equivariant operators with generalized finite propagation on a subset of $`M`$. The proof of the first lemma follows the same argument as that in and is stated without proof.
Definition: Let $`A`$ be a closed subspace of a proper metric space $`M`$. Denote by $`D_G^{}(A,M)`$ the $`C^{}`$-algebra of all operators $`T`$ in $`D_G^{}(M)`$ such that $`\text{Supp}T\text{Pen}(\pi ^1(A),R)\times \text{Pen}(\pi ^1(A),R)`$, for some $`R>0`$. Let $`C_G^{}(A,M)`$ be the quotient $`D_G^{}(A,M)/I`$.
Lemma 2: Let $`(A,B)`$ be a decomposition of $`M`$. Then
1. $`C_G^{}(A,M)+C_G^{}(B,M)=C_G^{}(M)`$.
2. $`C_G^{}(A,M)C_G^{}(B,M)=C_G^{}(AB,M)`$ if in addition we assume that $`(A,B)`$ is coarsely excisive.
Lemma 3: Suppose $`VM`$ and $`\pi _1(V)`$ injects into $`\pi _1(M)`$. There is an isomorphism $`K_{}(C_G^{}(V))K_{}(C_G^{}(V,M))`$.
Proof: Let $`\pi :\stackrel{~}{M}M`$ be the projection map. Consider the $`C^{}`$-algebra $`C^{}(\text{Pen}(\pi ^1(V),n),\pi _1(M))`$ given by the quotient by $`I`$ of the $`C^{}`$-algebra of locally compact, $`\pi _1(M)`$-equivariant operators on the $`n`$-neighborhood penumbra $`\text{Pen}(\pi ^1(V),n)`$. Then
$$C_G^{}(V,M)=\underset{}{lim}C^{}(\text{Pen}(\pi ^1(V),n),\pi _1(M)).$$
The inclusion map $`i:\pi ^1(V)\text{Pen}(\pi ^1(V),n)`$ is a coarse equivalence. Since by the construction the generalized Roe algebra its operators are defined up to their bounded parts, the map $`i`$ induces a series of isomorphisms
$$\begin{array}{ccc}K_{}(C^{}(\pi ^1(V),\pi _1(M)))& & K_{}(C^{}(\text{Pen}(\pi ^1(V),n),\pi _1(M))\hfill \\ & & K_{}(limC^{}(\text{Pen}(\pi ^1(V),n),\pi _1(M))\hfill \\ & & K_{}(C_G^{}(V,M)).\hfill \end{array}$$
Since $`\pi _1(V)\pi _1(M)`$ is an injection, the inverse image $`\pi ^1(V)\stackrel{~}{M}`$ is a disjoint union of isomorphic copies of $`\stackrel{~}{V}`$, parametrized by the coset space $`\pi _1(M)/\pi _1(V)`$. Therefore, there is a one-to-one correspondence between $`\pi _1(M)`$-equivariant operators on $`\pi ^1(V)`$ and $`\pi _1(V)`$-equivariant operators on $`\stackrel{~}{V}`$. Hence $`C^{}(\pi ^1(V),\pi _1(M))C_G^{}(V)`$. We then have $`K_{}(C_G^{}(V))K_{}(C_G^{}(V,M))`$, as desired. $`\mathrm{}`$
Let $`(A,B)`$ be a coarsely excisive decomposition of $`M`$ such that $`V=AB`$ satisfies $`\pi _1(V)\pi _1(M)`$. The boundary operator $`:K_j(C_G^{}(A,M)+C_G^{}(B,M))K_{j1}(C_G^{}(A,M)C_G^{}(B,M))`$ arising from the coarse Mayer-Vietoris sequence is by the previous lemmas truly a map
$$:K_{}(C_G^{}(M))K_1(C_G^{}(V)).$$
Theorem: (Boundary of Dirac is Dirac) Consider a coarsely excisive decomposition $`(A,B)`$ of $`M`$ and let $`V=AB`$. If $`:K_{}(C_G^{}(M))K_1(C_G^{}(V))`$ is the boundary map from the Mayer-Vietoris sequence derived above, then we have $`(\text{ Ind}_M(D))=\text{ Ind}_V(D)`$.
Remark: Here $`\text{ Ind}_M(D)`$ and $`\text{ Ind}_V(D)`$ represent the generalized coarse indices of the spinor Dirac operators on $`M`$ and $`V`$, respectively. We will continue to use a subscript if the space to which the index is related is ambiguous. The “boundary of Dirac is Dirac” principle is essentially equivalent to Bott periodicity in topological $`K`$-theory. In all cases considered here, there are commutative diagrams relating topological boundary to the boundary operator arising in the $`K`$-theory of $`C^{}`$-algebras, and on the topological side, a consideration of symbols suffices. See , , and .
Theorem 1: The $`n`$-fold product $`M`$ of punctured two-dimensional tori does not have a metric of uniform positive scalar curvature in the same coarse equivalence class as the positive hyperoctant with its standard Euclidean metric.
Proof: Consider the projection map $`p:M_0^n`$ from the multifold product $`M=\stackrel{}{T}\times \mathrm{}\times \stackrel{}{T}`$ to the positive hyperoctant, where each component $`p_i`$ is the quasi-isometric projection of the punctured torus onto the positive reals numbers. Take a hypersurface $`S_0^n`$ sufficiently far from the origin so that the inverse image of every point on $`S`$ is an $`n`$-torus, and so the space $`V`$ is coarsely equivalent to the $`(2n1)`$-dimensional noncompact manifold $`^{n1}\times T^n`$. The complement of the hypersurface $`V`$ consists of two noncompact components. Define $`A`$ to be the closure of the component containing the inverse image $`p^1(\mathrm{𝟎})`$ of the origin in $`_0^n`$. Take $`B`$ the closure of $`M\backslash A^{}`$. Then the pair $`(A,B)`$ forms a coarsely excisive decomposition of the space $`M`$ whose intersection is $`AB=V`$.
Consider the generalized coarse index $`\text{ Ind}_M(D)K_{}(C_G^{}(M))`$ of the lifted classical Dirac operator on the pullback spinor bundle of the universal cover $`\stackrel{~}{M}`$. Note that $`\pi _1(M)`$ is the $`n`$-fold product $`F_2\times \mathrm{}\times F_2`$ of free groups, and that $`\pi _1(V)\pi _1(^{n1}\times T^n)^n`$. Hence there is an injection $`\pi _1(V)\pi _1(M)`$ and the $`K`$-theoretic Mayer-Vietoris sequence applies. The boundary map $``$ of this sequence satisfies $`(\text{ Ind}_M(D))=\text{ Ind}_V(D)K_{}(C_G^{}(V))`$. However, $`V`$ is coarsely equivalent to the hypersurface $`^{n1}\times T^n`$, so the index $`\text{ Ind}_V(D)`$ can be taken to live in $`K_1(C_G^{}(^{n1}\times T^n))`$. Note that $`n`$ will be taken to be at least $`2`$. There is yet another boundary map $`K_1(C_G^{}(^{n1}\times T^n))K_{n+1}(C_G^{}(\times T^n))`$ by peeling off $`n2`$ copies of the real line. This boundary map (or composition of $`n2`$ boundary maps) preserves index.
Recall that $`D_G^{}(M)`$ is the norm closure of the $`C^{}`$-algebra of all locally compact, $`\pi _1(M)`$-equivariant, generalized finite propagation operators on $`L^2(\stackrel{~}{M})`$, and $`ID_G^{}(M)`$ is the closure of the ideal of such operators $`T`$ that satisfies the condition that $`(\pi \times \pi )(\text{Supp}T)`$ is bounded in $`M\times M`$. The short exact sequence $`0ID_G^{}(\times T^n)C_G^{}(\times T^n)0`$ gives rise to the six-term exact sequence in $`K`$-theory:
Notice that the map $`K_{}(I)K_{}(D_G^{}(\times T^n))`$ induced by the inclusion is the zero map by an Eilenberg swindle argument. Hence both maps $`K_{}(D_G^{}(\times T^n))K_{}(C_G^{}(\times T^n))`$ are injections.
If $`n`$ is even, the generalized coarse index $`\text{ Ind}_{\times T^n}(D)`$ of the Dirac operator $`D`$ resides in $`K_1(D_G^{}(\times T^n))`$. Certainly the image of this index under the boundary map $`K_1(D_G^{}(\times T^n))K_0(D_G^{}(T^n))`$ is the index $`\text{ Ind}_{T^n}(D)`$ of $`D`$ on the $`n`$-torus. Since $`T^n`$ does not have a metric of positive scalar curvature at all, the obstruction $`\alpha (T^n,f)K_{}(C^{}(^n))`$, where $`f:T^nB^n`$ is the classifying map, is nonvanishing. This “index” is constructed by Rosenberg in . This special case of the Gromov-Lawson-Rosenberg conjecture holds for the group $`^n`$. This index maps to our generalized coarse index $`\text{ Ind }_{T^n}D`$ under the isomorphism $`K_{}(C^{}(^n))K_{}(D_G^{}(T^n))`$. Hence the index of $`D`$ in $`K_1(D_G^{}(\times T^n))`$ is nonzero, and its projection onto the group $`K_1(C_G^{}(\times T^n))`$ is nonzero as well. This argument gives us the necessary index obstruction.
If $`n`$ is odd, we apply the same argument as above with respect to the map $`K_0(D_G^{}(\times T^n))K_0(C_G^{}(\times T^n))`$. $`\mathrm{}`$
The extension of this method to multifold products of hyperbolic manifolds involves the Margulis lemma, which states that in such a space there exists a small positive constant $`\mu =\mu _n`$ such that the subgroup $`\mathrm{\Gamma }_\mu (V,v)\pi _1(V,v)`$ generated by loops of length less than or equal to $`\mu `$ based at $`vV`$ is almost nilpotent, i.e. it contains a nilpotent subgroup of finite index. It can be shown that there exist such cusps, or submanifolds $`CV`$ with compact convex boundary containing $`v`$, such that $`C`$ is diffeomorphic to the product $`C\times _+`$, where $`C`$ is diffeomorphic to an $`(n1)`$-dimensional nilmanifold with fundamental group containing $`\mathrm{\Gamma }_\mu (V,v)`$. Here a nilmanifold signifies a quotient $`N/\mathrm{\Gamma }`$ of a nilpotent Lie group by a cocompact lattice $`\mathrm{\Gamma }`$. The nilmanifolds that arise in this context as boundaries of pseudospheres will have a naturally flat structure.
Theorem 2: An $`n`$-fold product of hyperbolic manifolds has no uniform positive scalar curvature metric coarsely equivalent to the usual Euclidean metric on the positive Euclidean hyperoctant.
Proof: Without loss of generality, it suffices to consider the case when the noncompact hyperbolic spaces have only one cusp. Let $`m`$ be the dimension of this product manifold. As in the multifold product of tori, there is a positive $`b`$ such that on each hyperbolic space $`_i`$ the inverse image of each point $`xb`$ under the projection $`_i_0`$ is by Margulis’ lemma a flat compact connected Riemannian manifold of finite dimension. Consider the inverse image $`𝒱`$ under the induced product map $`p:_1\times \mathrm{}\times _m_0^m`$ of the same hypersurface as described in the previous theorem. By Bieberbach’s theorem, every flat compact connected Riemannian manifold admits a normal Riemannian covering by a flat torus of the same dimension. Hence $`𝒱`$ is covered by some product of Euclidean space and a higher-dimensional torus. Any metric of positive scalar curvature on $`𝒱`$ would certainly lift to such a metric in this covering space. Using the same induction argument as before, we show that such a metric is obstructed by the presence of a nonzero Dirac class. $`\mathrm{}`$
III. Noncompact Quotients of Symmetric Spaces: A Special Case
The Iwasawa decomposition gives a unique way of expressing the group $`\mathrm{SL}_n()`$ as a product $`\mathrm{SL}_n()=NAK`$, where $`N`$ is the subgroup of standard unipotent matrices (upper triangular matrices with all diagonal entries equal to $`1`$), $`A`$ the subgroup of $`\mathrm{SL}_n()`$ consisting of diagonal matrices with positive entries, and $`K`$ the orthogonal subgroup $`\mathrm{SO}_n()`$. The effect of taking the double quotient of $`\mathrm{SL}_n()`$ by both $`\mathrm{SO}_n()`$ and $`\mathrm{SL}_n()`$ on the Iwasawa decomposition is that we are left with classes of matrices represented by those of the form $`n^{}a^{}`$, where $`n^{}`$ is represented by a unipotent matrix and $`a^{}=\text{ diag }(a_1,\mathrm{},a_n)`$ is diagonal with weakly increasing entries $`a_1a_2\mathrm{}a_n`$. Let $`N^{}`$ be the iterated circle bundle that arises upon taking a quotient of $`N`$ by $`\mathrm{SL}_n()`$ (for example, when $`n=3`$, the group $`N`$ is a Heisenberg group). Let $`A^{}`$ be the semigroup of matrices (note that it is not closed under inversion) with such an increasing condition on the entries. Consider the map $`\rho :\mathrm{SL}_n()\backslash \mathrm{SL}_n()/\mathrm{SO}_n()A^{}`$ given by $`n^{}a^{}a^{}`$. This map denotes a fiber bundle with fiber $`N^{}`$ over every point $`a^{}A^{}`$. Notice that the Iwasawa decomposition gives another way of observing the dimension of $`\mathrm{SL}_n()/\mathrm{SO}_n()`$, the sum of $`n1`$ dimensions from $`A`$ and $`\frac{n(n1)}{2}`$ from $`N`$. The quotient $`\mathrm{SL}_n()\backslash \mathrm{SL}_n()/\mathrm{SO}_n()`$ then has $`n1`$ noncompact directions coming from $`A^{}`$ and $`\frac{n(n1)}{2}`$ compact directions from $`N^{}`$.
The space $`A^{}`$ is identifiable with the subset of $`(n1)`$-dimensional Euclidean space given by $`\{(a_1,\mathrm{},a_n):0<a_1\mathrm{}a_n,a_1\mathrm{}a_n=1\}`$. We wish to construct a hypersurface in $`A^{}`$ whose inverse image under $`\rho `$ is an iterated circle bundle over Euclidean space. We recall that a geodesic in the quotient space $`\mathrm{SL}_n()/\mathrm{SO}_n()`$ is of the form $`te^{t\mathrm{\Lambda }}\mathrm{SO}_n()`$, where $`\mathrm{\Lambda }`$ is an $`n\times n`$ symmetric matrix with zero trace. We shall construct an appropriate hypersurface in $`A^{}`$ by taking the union of sufficiently many geodesics.
Consider a geodesic $`te^{t\mathrm{\Lambda }}\mathrm{SO}_n()`$ in the quotient $`\mathrm{SL}_n()/\mathrm{SO}_n()`$ where $`\mathrm{\Lambda }=\text{ diag }(\lambda _1,\mathrm{},\lambda _n)`$ with $`\lambda _1\mathrm{}\lambda _n`$ and $`\lambda _1+\mathrm{}+\lambda _n=0`$. Then the geodesic is a map $`t\text{ diag }(e^{\lambda _1t},\mathrm{},e^{\lambda _nt})\mathrm{SO}_n()`$. Because two symmetric traceless matrices of the form $`\mathrm{\Lambda }_1=\text{ diag }(\lambda _1,\mathrm{},\lambda _n)`$ and $`\mathrm{\Lambda }_2=\text{ diag }(\alpha \lambda _1,\mathrm{},\alpha \lambda _n)`$ give the same geodesic image for $`\alpha >0`$, we can normalize the $`\lambda `$-vector so that $`\lambda _1=1`$. In addition, let $`\mu _2=1`$ and $`\mu _i[i1,i]`$ for $`i=3,\mathrm{},n`$ and let $`m=_{i=2}^n\mu _i`$. Set $`\lambda _i=\mu _i/m`$. Then $`\mathrm{\Lambda }=\text{ diag }(\lambda _1,\mathrm{},\lambda _n)`$ is symmetric and traceless with weakly increasing entries.
Lemma 4: Each ordered $`(n1)`$-tuple $`(\mu _2,\mathrm{},\mu _n)[\mathrm{\hspace{0.17em}2},3]\times \mathrm{}\times [n,n+1]`$ gives rise to a unique geodesic $`te^{t\mathrm{\Lambda }}\mathrm{SO}_n()`$, up to reparametrization .
Proof: Suppose $`(\mu _2^{(1)},\mathrm{},\mu _n^{(1)}),(\mu _2^{(2)},\mathrm{},\mu _n^{(2)})[\mathrm{\hspace{0.17em}2},3]\times \mathrm{}\times [n,n+1]`$ give rise to the same geodesic. The two vectors correspond to the traceless matrices $`\mathrm{\Lambda }_1=\text{ diag }(1,\lambda _2^{(1)},\mathrm{},\lambda _n^{(1)})`$ and $`\mathrm{\Lambda }_2=\text{ diag }(1,\lambda _2^{(2)},\mathrm{},\lambda _n^{(2)})`$, respectively. Let $`\nu _1=(1,\lambda _2^{(1)},\mathrm{},\lambda _n^{(1)})`$ and $`\nu _2=(1,\lambda _2^{(2)},\mathrm{},\lambda _n^{(2)})`$. Obtain the normalized matrices $`\mathrm{\Lambda }_1^{}`$ and $`\mathrm{\Lambda }_2^{}`$ by dividing the entries by the Euclidean norms $`\nu _1`$ and $`\nu _2`$. By assumption, $`\mathrm{\Lambda }_1^{}=\mathrm{\Lambda }_2^{}`$; in other words, $`\frac{\nu _1}{\nu _1}=\frac{\nu _2}{\nu _2}`$. But then $`\theta =\mathrm{cos}^1\frac{\nu _1\nu _2}{\nu _1\nu _2}=0`$, so $`\nu _1`$ and $`\nu _2`$ are parallel. Since their first coordinates coincide, they are identical. Hence $`\frac{(\mu _2^{(1)},\mathrm{},\mu _n^{(1)})}{\mu _2^{(1)}+\mathrm{}+\mu _n^{(1)}}=\frac{(\mu _2^{(2)},\mathrm{},\mu _n^{(2)})}{\mu _2^{(2)}+\mathrm{}+\mu _n^{(2)}}`$. Since each $`\mu _i^{(j)}`$ is positive, it can be written as the square of some other positive number. Arguing as before, we see that $`(\mu _2^{(1)},\mathrm{},\mu _n^{(1)})=\beta (\mu _2^{(2)},\mathrm{},\mu _n^{(2)})`$. for some $`\beta >0`$. Since $`\mu _2^{(1)}`$ and $`\mu _2^{(2)}`$ both equal one, the vectors are coincident. $`\mathrm{}`$
The space $`A^{}`$ of diagonal matrices in $`\mathrm{SL}_n()`$ with weakly increasing entries is itself simply-connected of dimension $`n1`$, and its space at infinity is an $`(n2)`$-dimensional simplex $`P`$. Let $`W`$ be the union of geodesics constructed above, with $`(\mu _2,\mathrm{},\mu _n)`$ ranging in the product $`[\mathrm{\hspace{0.17em}2},3]\times \mathrm{}\times [n,n+1]`$ of intervals. The subset $`W^{}A^{}`$ is an $`(n1)`$-dimensional space with boundary $`V^{}=W^{}`$. The $`(n2)`$-dimensional hypersurface $`V^{}`$ is also simply-connected whose space at infinity is homeomorphic to an $`(n3)`$-sphere that is disjoint with $`P`$. The space $`V^{}`$ itself is coarsely equivalent to $`^{n2}`$.
Theorem 3: The double quotient group $`\mathrm{SL}_n()\backslash \mathrm{SL}_n()/\mathrm{SO}_n()`$ does not have a uniform positive scalar curvature metric that is coarsely equivalent to the natural one inherited from $`\mathrm{SL}_n()`$.
Proof: Let $`M=\mathrm{SL}_n()\backslash \mathrm{SL}_n()/\mathrm{SO}_n()`$ and recall the projection map $`\rho :MA^{}`$ given by $`n^{}a^{}a^{}`$. Since the fiber over each point is an arithmetic quotient of the group of unipotent matrices, the inverse image $`V\rho ^1(V^{})`$ is coarsely equivalent to an iterated circle bundle over Euclidean space. Moreover, $`V`$ partitions the space into a coarsely excisive pair whose closures $`(A,B)`$ satisfy the equalities $`AB=M`$ and $`AB=V^{}`$ (note that $`B`$ can be taken as $`\rho ^1(W^{})`$ and $`A`$ the closure of its complement). If $`\text{ Ind}_M(D)`$ denotes the generalized coarse index of the classical spinor Dirac operator on $`\stackrel{~}{M}`$, then the Mayer-Vietoris map $`:K_{}(C_G^{}(M))K_1(C_G^{}(V))`$ defined in the previous chapter satisfies $`(\text{ Ind}_M(D))=\text{ Ind}_V^{}(D)=\text{ Ind}_{^{n2}\times U^m}(D)`$, where $`U^m`$ is the compact fiber of the iterated circle bundle of dimension $`m=\frac{n(n1)}{2}`$. Applying the same argument as before, it suffices to show that the index of the Dirac operator in $`K_{}(D_G^{}(U^m))`$ is nonzero. However, $`U^m`$ is a quotient of a nilpotent group with a cocompact lattice, and hence by Gromov and Lawson has no metric of positive scalar curvature at all. As with the theorem for punctured tori, there is a nonvanishing Rosenberg index $`\alpha (U^m)K_{}(C^{}(\pi ))`$, where $`\pi =\pi _1(U^m)`$, which maps to the generalized coarse index in $`K_{}(D_G^{}(U^m))`$, as desired. Here the Gromov-Lawson-Rosenberg conjecture is true since $`U^m`$ is a nilmanifold. $`\mathrm{}`$
IV. The General Noncompact Arithmetic Case
To understand how we might achieve a similar result for general double quotient spaces $`\mathrm{\Gamma }\backslash G/K`$, we appeal to the following.
Picture from Reduction Theory: There is a compact polyhedron $`P`$ and a Lipschitz map $`\pi :McP`$, where $`cP`$ is the open cone on $`P`$ so that (1) every point inverse deform retracts to an arithmetic manifold, (2) $`\pi `$ respects the radial direction, and (3) all point inverses have uniformly bounded size.
Indeed, the polyhedron $`P`$ is the geometric realization of the category of proper $``$-parabolic subgroups of $`G`$, modulo the action of $`\mathrm{\Gamma }`$, and $`\pi ^1`$ of the barycenter of a simplex is the arithmetic symmetric space associated to that parabolic. Concretely, for $`\text{SL}_n()\text{SL}_n()`$, the space $`P`$ is an $`n2`$ simplex, the parabolics correspond to flags, and the associated arithmetic groups have a unipotent normal subgroup with quotient equal to a product of $`\text{SL}_{m_i}()`$, where the $`m_i`$ are sizes of the blocks occurring in the flag. As one goes to infinity, the unipotent directions shrink in diameter and are responsible for the finite volume property of the lattice quotient, while the other parabolic directions remain of bounded size.
Alternatively, for any choice of basepoint in the homogeneous space, there are constants $`C`$ and $`D`$ that satisfy the following condition: if $`x`$ is a given point and $`Q_x`$ is the largest parabolic subgroup associated with a simplex whose cone contains $`x`$ within its $`C`$-neighborhood, then the orbit of $`x`$ under $`Q_x`$ has diameter less than $`D`$. Note the empty simplex means that there is a compact core which is stabilized by the whole group.
This picture can essentially be ascertained from , ; the fact that $`K\backslash G/\mathrm{\Gamma }`$ has finite Gromov-Hausdorff distance from $`cP`$ is asserted in . However, one needs a key estimate about the “coarse isotropy.” Details are given in unpublished work of Eskin ; some are given below.
As a guide the reader should think through the picture suggested by a product of hyperbolic manifolds. Each hyperbolic manifold contributes to $`cP`$ either a point, in the compact case, or the open cone on a finite set of points, in the case of cusps. Thus $`P`$ is a join of some number of finite sets. Using this model, we find that the inverse image of any point in the interior of any simplex is exactly a product of closed hyperbolic manifolds, cores of hyperbolic manifolds, and flat manifolds.
Let us now consider the unique decomposition of a semisimple Lie group $`G=N_xA_xK`$, where $`x`$ is a point on the space of infinity of $`G/K`$. If $`\mathrm{\Gamma }`$ is an arithmetic lattice in $`G`$, we are interested in knowing how $`\mathrm{\Gamma }`$ acts on $`G/K`$. In other words, we ask how $`\mathrm{\Gamma }`$ acts on this particular coordinate system. Let $`g=nakG`$. If $`\gamma \mathrm{\Gamma }`$, let $`\gamma g=n^{}a^{}k^{}`$. Notice that $`N_x\mathrm{\Gamma }`$ is a lattice in $`N_x`$, and acts cocompactly on $`N_x`$. Consider the projection $`\rho :\mathrm{\Gamma }\backslash G/KA_x^\mathrm{\Gamma }`$ given by $`n^{}a^{}a^{}`$, where $`A_x^\mathrm{\Gamma }`$ is a fundamental domain of $`A_x`$ under the action of $`\mathrm{\Gamma }`$. The fiber above each point is a compact manifold arising from the action of $`N_x\mathrm{\Gamma }`$ on $`N_x`$. Let $`𝔣+𝔭`$ be the usual Cartan decomposition and let $`X𝔭`$ be the element such that $`dp(X)=\gamma _{px}^{}(0)`$. If $`Z(X)`$ is defined by $`Z(X)=\{Y𝔤:[X,Y]=0\}`$, then $`𝔞=Z(X)𝔭`$ is the unique maximal abelian subspace of $`𝔭`$ that contains $`X`$. By definition $`A_x=\mathrm{exp}(𝔞)`$.
We are interested in constructing a hypersurface $`V^{}`$ in the space $`A_x^\mathrm{\Gamma }`$ whose inverse image $`V\rho ^1(V^{})`$ under the projection map $`\rho `$ provides us with an appropriate excisive decomposition of $`\mathrm{\Gamma }\backslash G/K`$ for which the $`K`$-theoretic Mayer-Vietoris sequence is applicable. Consider the chamber decomposition of $`𝔞`$. The Weyl group $`W`$ acts on these chambers via the hyperplanes. Consider the Bruhat decomposition $`G=_{wW}BwB`$ of $`G`$, and let $`\gamma BwB`$ for some $`wW`$. Recall that, if $`g=nak`$, we write $`\gamma g=n^{}a^{}k^{}`$. Denote by $`^+`$ the set of positive roots of $`𝔞^{}`$ and $`^{}`$ the set of negative roots. Let $`=^{}w^+`$ be the set of roots that start out positive but are made negative under the action of $`w`$. For some positive reals constants $`c_\alpha `$, the following equation holds ,:
$$a^{}=wa\underset{\alpha }{}c_\alpha \alpha (a)+O(1),$$
where $`a`$ and $`a^{}`$ are viewed as elements of the Lie algebra $`𝔞`$. The implications of this equation are as follows. Consider an element $`a`$ in the positive Weyl chamber $`𝒞(a)`$. Then the intersection $`\mathrm{\Gamma }(a)𝒞(a)`$ of the orbit $`\mathrm{\Gamma }(a)`$ of $`a`$ under $`\mathrm{\Gamma }`$ and the Weyl chamber $`𝒞(a)`$ containing $`a`$ has a bounded diameter, uniformly in $`a`$. In other words, if $`\gamma (a)`$ stays in the same positive Weyl chamber, then $`w`$ is the identity and $``$ is empty. Hence $`a^{}=a+o(1)`$, implying that $`a^{}`$ can be found at a uniformly bounded distance from $`a`$ itself. In this event, the action of $`\gamma `$ corresponds to a translation of $`a`$ to (possibly) the compact fiber direction of $`\mathrm{\Gamma }\backslash G/K`$. It is also a general fact that the action of any $`g\mathrm{\Gamma }`$ will take the vertex of any subsector (as drawn in the previous figure) to the vertex of an analogous subsector. With this machinery, we are able to prove the following.
Theorem 4: The double quotient space $`M=\mathrm{\Gamma }\backslash G/K`$ has no metric of uniform positive scalar curvature in the same coarse class as the natural metric inherited from $`G`$.
Proof: The picture from reduction theory provides a polyhedron $`P`$ in the space of infinity of the positive Weyl chamber and an open cone $`W=cP`$ on $`P`$ in $`𝒞^+`$ oriented so that the distance from $`W`$ to any hyperplane $`\alpha =0`$ will exceed the quantity $`sup_{a𝒞^+}\text{diam }(\mathrm{\Gamma }(a)𝒞^+)`$. This quantity is finite by , . Let $`V^{}=W`$ and $`V=\rho ^1(V^{})`$ in $`M=\mathrm{\Gamma }\backslash G/K`$. Then $`V`$ is a hypersurface that induces a decomposition $`(A,B)`$ of $`M`$. The fundamental group $`\pi _1(V)=N\mathrm{\Gamma }`$ injects into $`\pi _1(\mathrm{\Gamma }\backslash G/K)=\mathrm{\Gamma }`$, satisfying the requirement needed in the construction of the Mayer-Vietoris sequence. In the most general case, the space $`V`$ is coarsely equivalent to a bundle over Euclidean space whose fiber consists of two components: a nilmanifold $`N`$ and (possibly) a compact homogeneous manifold $`H`$. In the absence of such an $`H`$, the argument follows exactly as it does for $`\mathrm{SL}_n()\backslash \mathrm{SL}_n()/\mathrm{SO}_n()`$. In the presence of a compact homogeneous manifold, we may pass the coarse index of the Dirac operator to $`\times H`$ and use the usual Rosenberg obstruction on $`H`$ as in Theorem 1 to obtain our desired result. $`\mathrm{}`$
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# SPACE-TIMES ADMITTING A THREE-DIMENSIONAL CONFORMAL GROUP
## 1 Introduction
We shall study perfect fluid space-times admitting a three-dimensional Lie group of conformal motions $`C_3`$, which contains a two-dimensional Abelian Lie subgroup of isometries $`G_2`$. The precise assumptions are:
1. The $`G_2`$ acts orthogonally transitively on two-dimensional spacelike orbits diffeomorphic to $`\mathrm{I}\mathrm{R}^2`$.
2. One of the Killing vectors (KV) in $`G_2`$ is hypersurface orthogonal, hence the metric can be written in diagonal form .
3. The conformal Killing vector (CKV) is proper (non-homothetic).
4. The orbits of the $`C_3`$ are non-null.
Space-times of this type have been studied recently by Kramer and Carot and Castejón-Amenedo and Coley . Here we shall systematically classify geometrically all possible space-times according to their (inequivalent) group structures and the nature of their corresponding orbits, and thereafter attempt to obtain the perfect fluid models in each class. In this classification scheme the (stationary and axisymmetric) models studied by Kramer and Kramer and Carot admit a $`C_3`$ acting on three-dimensional timelike hypersurfaces. The models studied by Castejón-Amenedo and Coley admit a $`C_3`$ acting on two-dimensional timelike hypersurfaces containing an Abelian $`C_2`$. Mars and Senovilla are looking at models in which $`C_3`$ is Abelian and trying to extend their study to other non-diagonal cases and Coley and Czapor are studying space-times in which the CKV is inheriting .
This work generalizes previous work in several ways. First, solutions admitting a CKV in space-times with additional symmetry have been studied; for example, spatially homogeneous , spherically symmetric (see, for instance and references cited therein) and plane symmetric models have been investigated. In a sense the present work is the natural generalization of this research in that CKV are studied in space-times with the next highest degree of symmetry. Indeed, it is for this reason that $`G_2`$ space-times have begun to attract much attention. Cosmological models with an Abelian $`G_2`$ acting on spacelike hypersurfaces have been studied by Hewitt et al., Hewitt and Wainwright , Hewitt et al., Ruiz and Senovilla and Van den Berg and Skea . Models with an Abelian $`G_2`$ acting on timelike hypersurfaces, including the astrophysically relevant stationary and axisymmetric models, which have been studied for many years , have also attracted renewed attention . Second, the cases in which the CKV degenerates to either a KV or a homothetic vector field (HVF) have been studied previously. Space-times admitting a $`G_3`$ have received much attention . Space-times admitting a HVF in addition to two KV have been studied by Hewitt and Wainwright and Carot et al. . Third, this work generalizes that of Coley and Czapor in which the CKV is inheriting, and complements that of Senovilla in which the Abelian $`G_2`$ acts on two-dimensional timelike hypersurfaces (these models are, in fact, stationary and axisymmetric perfect fluid solutions admitting a proper CKV). Finally, we shall assume that the perfect fluid matter satisfies $`\mu >0`$ (all vacuum space-times admitting a proper CKV are known, ).
For a very precise and accurate study of the general properties of CKVs, their Lie algebra and fixed point structure, we refer the reader to .
The outline of this paper is as follows. In section 2 we shall describe the models under consideration in detail, defining all relevant terms. We shall then classify the space-times into different inequivalent classes, in each of the two cases in which the CKV is everywhere spacelike or not (i.e., the CKV can be timelike, see Eqns. (9) and (10)), according to the structure of their corresponding three-dimensional conformal Lie group $`C_3`$ (each containing an Abelian $`G_2`$). We shall display the corresponding metrics (and symmetry vectors) in coordinates adapted to the CKV. These results are completely independent of the Einstein field equations and the assumed energy-momentum tensor. We note parenthetically that not all the classes may be possible (again, irrespective of the assumed matter content) in some situations of physical interest; for example, Mars and Senovilla have shown that in axially symmetric space-times admitting a maximal three-dimensional (or two-dimensional) conformal Lie algebra (such as in a cylindrically symmetric space-time with one CKV), the axial KV must commute with the other two (in fact, they show that this is true regardless of whether the additional KV is spacelike, such as, for example, in the case of a stationary and axially symmetric space-time with one CKV).
In section 3 we study those cases where one KV is hypersurface orthogonal, hence the metric can be written in diagonal form. In sections 4 and 5 we then utilize Einstein’s field equations in an attempt to find perfect fluid solutions representative of the classes in each of the two cases corresponding to the form of the CKV. We shall assume that the matter satisfies the weak and dominant energy conditions ($`\mu >0`$, $`\mu p\mu `$). In the final section all of the perfect fluid solutions obtained are summarized and briefly discussed.
## 2 Space-times admitting conformal Killing vectors
As we have already pointed out, we shall concern ourselves with space-times $`(M,g)`$ that admit a three-parameter conformal group $`C_3`$ containing an Abelian two-parameter subgroup of isometries $`G_2`$, whose orbits $`S_2`$ are spacelike, diffeomorphic to $`\mathrm{I}\mathrm{R}^2`$ and admit orthogonal two-surfaces; furthermore, we shall assume that the $`C_3`$ acts transitively on non-null orbits $`V_3`$. As is customary, we shall denote the KV spanning $`𝒢_2`$ (the Lie algebra associated with $`G_2`$) by $`\xi `$ and $`\eta `$, and the proper CKV in $`𝒞_3`$ (the Lie algebra of the conformal group $`C_3`$) by $`X`$. Since, by hypothesis, $`\xi `$ and $`\eta `$ commute, we can (locally) adapt two coordinates, say $`y`$ and $`z`$, so that
$$\xi =\frac{}{y},\eta =\frac{}{z}.$$
(1)
Taking now two more coordinates $`t`$ and $`r`$, it follows from the above assumptions that the line element associated with the metric $`g`$ can be written as
$$ds^2=e^{2F}\{dt^2+dr^2+Q[h^2(dy+wdz)^2+h^2dz^2]\},$$
(2)
where $`F`$, $`Q`$, $`h`$ and $`w`$ are all functions of $`t`$ and $`r`$ alone.
The (proper) CKV $`X`$ will satisfy
$$(_Xg)_{ab}=2\mathrm{\Psi }g_{ab},$$
(3)
where $``$ stands for the Lie derivative operator and $`\mathrm{\Psi }=\mathrm{\Psi }(x^c)`$ is the so-called conformal factor (the particular cases $`\mathrm{\Psi }=0`$ and $`\mathrm{\Psi }=const0`$ correspond, respectively, to $`X`$ being a KV and a proper HVF). Assuming that no further CKV exist on $`M`$ (i.e., the $`C_3`$ is maximal) one has the following two families of Lie algebra structures for $`𝒞_3`$:
$`(\mathrm{A})`$ $`[\xi ,\eta ]=0,[\xi ,X]=\alpha _1\xi +\alpha _2\eta ,[\eta ,X]=\beta _1\xi +\beta _2\eta ,`$ (4)
$`(\mathrm{B})`$ $`[\xi ,\eta ]=0,[\xi ,X]=a_1\xi +a_2\eta +a_3X,[\eta ,X]=b_1\xi +b_2\eta +b_3X,`$ (5)
where the $`\alpha _i`$, $`\beta _i`$, $`a_i`$ and $`b_i`$ are constants.
The family $`𝒞_3(\mathrm{A})`$ can, in turn, be classified into the following seven Bianchi types (see for instance ):
$`(I)`$ $`[\xi ,\eta ]=[\xi ,X]=[\eta ,X]=0,`$
$`(II)`$ $`[\xi ,\eta ]=[\xi ,X]=0[\eta ,X]=\xi ,`$
$`(III)`$ $`[\xi ,\eta ]=0[\xi ,X]=\xi [\eta ,X]=0,`$
$`(IV)`$ $`[\xi ,\eta ]=0[\xi ,X]=\xi [\eta ,X]=\xi +\eta ,`$ (6)
$`(V)`$ $`[\xi ,\eta ]=0[\xi ,X]=\xi [\eta ,X]=\eta ,`$
$`(VI)`$ $`[\xi ,\eta ]=0[\xi ,X]=\xi [\eta ,X]=q\eta (q0,1),`$
$`(VII)`$ $`[\xi ,\eta ]=0[\xi ,X]=\eta [\eta ,X]=\xi +q\eta (q^2<4),`$
whereas family $`𝒞_3`$(B) can be brought to the following form by means of appropriate re-definitions of $`X`$ and the KV’s $`\xi `$ and $`\eta `$:
$$[\xi ,\eta ]=0[\xi ,X]=X[\eta ,X]=0.$$
(7)
Notice that in the case of $`X`$ being a HVF the above algebraic structure (7) is forbidden, since the Lie bracket of a HVF and a KV must necessarily be a KV (for further information on this case, see ).
Assuming now for the CKV $`X`$ an expression of the form
$$X=X^a(x^c)_a,$$
(8)
in the coordinate chart $`\{t,r,y,z\}`$, and specializing the equation (3) to the metric given in (2) and the CKV given in (8), it is easy to see that one can always, by means of a coordinate transformation in the $`t,r`$ plane, bring $`X`$ to one of the following forms:
$`(a)`$ $`X=_t+X^y(y,z)_y+X^z(y,z)_z,`$
$`(b)`$ $`X=_r+X^y(y,z)_y+X^z(y,z)_z,`$ (9)
$`(c)`$ $`X=_t+_r+X^y(y,z)_y+X^z(y,z)_z,`$
if the conformal algebra $`𝒞_3`$ belongs to the family (A), whence $`X^y(y,z)`$ and $`X^z(y,z)`$ are linear functions of their arguments to be determined from the commutation relations of $`X`$ with $`\xi `$ and $`\eta `$ (see (6-$`(I)`$) to (6-$`(VII)`$); and to the forms:
$`(a)`$ $`X=e^y(_t+X^y(t)_y+X^z(t)_z),`$
$`(b)`$ $`X=e^y(_r+X^y(r)_y+X^z(r)_z),`$ (10)
$`(c)`$ $`X=e^y(_t+_r+X^y(t,r)_y+X^z(t,r)_z),`$
if $`𝒞_3`$ is that given by (7) (family (B)).
The forms (9-$`a`$) and (10-$`a`$) are easily seen to correspond to the case of three-dimensional timelike conformal orbits $`T_3`$, whereas (9-$`b`$) and (10-$`b`$) correspond to three-dimensional spacelike conformal orbits $`S_3`$. The remaining possibilities, (9-$`c`$) and (10-$`c`$), imply null conformal orbits $`N_3`$ and will not be considered in the present paper.
Assuming now the form (9-$`a`$) for the CKV $`X`$ (i.e., family (A), timelike conformal orbits), one then has for each possible case ($`I`$) to ($`VII`$) the following forms for $`X`$ and the metric functions:
$`(I)`$ $`Q=\widehat{Q}(r),h^2=\widehat{h}^2(r),w=\widehat{w}(r),`$ (11)
$`X=_t.`$
$`(II)`$ $`Q=\widehat{Q}(r),h^2=\widehat{h}^2(r),w=\widehat{w}(r)t,`$ (12)
$`X=_t+z_y.`$
$`(III)`$ $`Q=e^t\widehat{Q}(r),h^2=e^t\widehat{h}^2(r),w=e^t\widehat{w}(r),`$ (13)
$`X=_t+y_y.`$
$`(IV)`$ $`Q=e^{2t}\widehat{Q}(r),h^2=\widehat{h}^2(r),w=\widehat{w}(r)t,`$ (14)
$`X=_t+(y+z)_y+z_z.`$
$`(V)`$ $`Q=e^{2t}\widehat{Q}(r),h^2=\widehat{h}^2(r),w=\widehat{w}(r),`$ (15)
$`X=_t+y_y+z_z.`$
$`(VI)`$ $`Q=e^{(1+q)t}\widehat{Q}(r),h^2=e^{(1q)t}\widehat{h}^2(r),w=e^{(1q)t}\widehat{w}(r),`$ (16)
$`X=_t+y_y+qz_z(q0,1).`$
$`(VII)`$ $`Q=e^{qt}{\displaystyle \frac{\sqrt{4q^2}}{2}}a(r),`$ (17)
$`h^2={\displaystyle \frac{\frac{\sqrt{4q^2}}{2}a(r)}{\sqrt{a(r)^2+c(r)^2+g(r)^2}+c(r)\mathrm{cos}(\sqrt{4q^2}t)+g(r)\mathrm{sin}(\sqrt{4q^2}t)}},`$
$`w={\displaystyle \frac{q}{2}}+{\displaystyle \frac{\frac{\sqrt{4q^2}}{2}[c(r)\mathrm{sin}(\sqrt{4q^2}t)g(r)\mathrm{cos}(\sqrt{4q^2}t)]}{\sqrt{a(r)^2+c(r)^2+g(r)^2}+c(r)\mathrm{cos}(\sqrt{4q^2}t)+g(r)\mathrm{sin}(\sqrt{4q^2}t)}},`$
$`X=_tz_y+(y+qz)_z(q^2<4).`$
In all of these cases $`F=F(t,r)`$ and the conformal factor, $`\mathrm{\Psi }`$, is given by
$$\mathrm{\Psi }=F_{,t}$$
(18)
The form (9-$`b`$) for the CKV $`X`$ (i.e., family A, spacelike conformal orbits) would yield similar results with the role of the coordinates $`t`$ and $`r`$ reversed.
Similarly, if we take the form (10-$`a`$) for $`X`$ (Family B, timelike conformal orbits) we get, assuming the canonical form (7) for $`𝒞_3`$, the following possibilities for the metric and the CKV:
$`(1)`$ $`ds^2=e^{2F}\{dt^2+dr^2+[\mathrm{\Phi }^2(r)(c\mathrm{cosh}t+d)^2+\mathrm{sinh}^2t]dy^2+`$ (19)
$`+2\mathrm{\Phi }^2(r)(c\mathrm{cosh}t+d)dydz+\mathrm{\Phi }^2(r)dz^2\},`$
$`X=e^y\left(_t\mathrm{coth}t_y+{\displaystyle \frac{c+d\mathrm{cosh}t}{\mathrm{sinh}t}}_z\right),`$
$`(2)`$ $`ds^2=e^{2F}\{dt^2+dr^2+[\mathrm{\Phi }^2(r)(ct^2+d)^2+t^2]dy^2+`$ (20)
$`+2\mathrm{\Phi }^2(r)(ct^2+d)dydz+\mathrm{\Phi }^2(r)dz^2\},`$
$`X=e^y\left(_t{\displaystyle \frac{1}{t}}_y{\displaystyle \frac{ct^2d}{|t|}}_z\right),`$
$`(3)`$ $`ds^2=e^{2F}\{dt^2+dr^2+[\mathrm{\Phi }^2(r)(dc\mathrm{cos}t)^2+\mathrm{sin}^2t]dy^2+`$ (21)
$`+2\mathrm{\Phi }^2(r)(dc\mathrm{cos}t)dydz+\mathrm{\Phi }^2(r)dz^2\},`$
$`X=e^y\left(_t\mathrm{cot}t_y+{\displaystyle \frac{d\mathrm{cos}tc}{\mathrm{sin}t}}_z\right),`$
where $`c`$ and $`d`$ are constants, $`F=F(t,r)`$ is an arbitrary function of $`t`$ and $`r`$, and the conformal factor $`\mathrm{\Psi }`$ is given by
$$\mathrm{\Psi }=e^yF_{,t}.$$
(22)
The case of spacelike conformal orbits for the family B (i.e., $`X`$ of the form (10-$`b`$)) leads to just one possibility, namely
$`ds^2=e^{2F}\{dt^2+dr^2+[\mathrm{\Phi }^2(t)(c\mathrm{sinh}r+d)^2+\mathrm{cosh}^2r]dy^2+`$ (23)
$`+2\mathrm{\Phi }^2(t)(c\mathrm{sinh}r+d)dydz+\mathrm{\Phi }^2(t)dz^2\},`$
$`X=e^y\left(_r\mathrm{tanh}r_y+{\displaystyle \frac{cd\mathrm{sinh}r}{\mathrm{cosh}r}}_z\right),`$
where $`c`$ and $`d`$ are again constants, $`F=F(t,r)`$ is an arbitrary function of its variables and the conformal factor is in this case
$$\mathrm{\Psi }=e^yF_{,r}.$$
(24)
## 3 Diagonal Case.
Next we will study those cases where two hypersurface orthogonal KVs exist in $`𝒢_2`$; this implies that the two KVs must be mutually orthogonal and therefore, by means of a linear change of coordinates in the Killing orbits $`S_2`$, one can always set $`w`$ in (2) to zero and the metric then becomes diagonal .
It is worth noticing that this is not possible for types $`II`$ and $`IV`$ in family A, as one can see by simply inspecting the form of $`w`$ in these cases (see (12) and (14)). As for the case $`VII`$ (also belonging to family A), $`w=0`$ implies the existence of a third KV tangent to the Killing orbits $`S_2`$, in which case they are of constant curvature and the conformal algebra becomes four dimensional; therefore, we shall not consider this case further. For those metrics in family B we have:
$`ds^2`$ $`=`$ $`e^{2F}\{dt^2+dr^2+\mathrm{sinh}^2tdy^2+\mathrm{\Phi }^2(r)dz^2\},`$
$`X`$ $`=`$ $`e^y\left(_t\mathrm{coth}t_y\right),`$ (25)
$`ds^2`$ $`=`$ $`e^{2F}\{dt^2+dr^2+t^2dy^2+\mathrm{\Phi }^2(r)dz^2\},`$
$`X`$ $`=`$ $`e^y\left(_t{\displaystyle \frac{1}{t}}_y\right),`$ (26)
$`ds^2`$ $`=`$ $`e^{2F}\{dt^2+dr^2+\mathrm{sin}^2tdy^2+\mathrm{\Phi }^2(r)dz^2\},`$
$`X`$ $`=`$ $`e^y\left(_t\mathrm{cot}t_y\right),`$ (27)
if the conformal orbits are timelike, and, for spacelike conformal orbits:
$`ds^2`$ $`=`$ $`e^{2F}\{dt^2+dr^2+\mathrm{cosh}^2rdy^2+\mathrm{\Phi }^2(t)dz^2\},`$
$`X`$ $`=`$ $`e^y\left(_r\mathrm{tanh}r_y\right).`$ (28)
There is still another possibility for metrics belonging to this family; namely,
$`ds^2`$ $`=`$ $`e^{2F}\{dt^2+dr^2+\mathrm{cos}^2tdy^2+\mathrm{sin}^2tdz^2\},`$
$`X`$ $`=`$ $`e^{y+z}\left(_t+\mathrm{tan}t_y\mathrm{cot}t_z\right),`$ (29)
but this metric admits a further proper CKV, $`_r`$, therefore the maximal conformal group is four-dimensional and again we shall not consider this case any further.
## 4 Perfect fluid space-times: Fluid flow tangent or orthogonal to the conformal orbits.
In this section we will study diagonal perfect fluid solutions admitting a maximal $`C_3`$ of conformal motions. The energy-momentum tensor for a perfect fluid is given by
$$T_{ab}=(\mu +p)u_au_b+pg_{ab},$$
(30)
where $`\mu `$ and $`p`$ are, respectively, the energy density and the pressure as measured by observers comoving with the fluid, and $`u^a`$ ($`u^au_a=1`$) is the four-velocity of the fluid. Since the metric admits two KVs, $`\xi `$ and $`\eta `$, it follows that the Lie derivative of $`p`$, $`\mu `$ and $`u_a`$ with respect to each must vanish identically; in our coordinate chart this is equivalent to:
$$\mu =\mu (t,r)p=p(t,r)u_a=u_a(t,r).$$
(31)
Since the isometry orbits admit orthogonal two-surfaces one has that, in these coordinates, both the metric and the Ricci tensor have block-diagonal forms; thus, the Einstein field equations specialized to a perfect fluid, as given by (30), reduce simply to
$$u_y=u_z=0,$$
(32)
$$\frac{G_{yy}}{g_{yy}}\frac{G_{zz}}{g_{zz}}=0,$$
(33)
$$G_{tr}^2\left(G_{tt}\frac{G_{yy}}{g_{yy}}g_{tt}\right)\left(G_{rr}\frac{G_{yy}}{g_{yy}}g_{rr}\right)=0$$
(34)
(where we have already taken into account that the metric is diagonal, otherwise we would have had an extra equation, namely: $`G_{yy}/g_{yy}G_{yz}/g_{yz}=0`$).
We shall focus on those solutions such that the perfect fluid four-velocity is either tangent or orthogonal to the conformal orbits; the general case (or “tilted” case) will be the subject of the next section.
### 4.1 Fluid flow tangent to the conformal orbits.
This case can only arise when the conformal orbits are timelike, and then one has:
$$u_t=e^Fu_r=0$$
(35)
(i.e., the coordinates are comoving). The metric must be one of those given by (11), (13), (15) and (16) (with $`w=0`$) if belongs to family A (Bianchi types $`I`$, $`III`$, $`V`$ and $`VI`$, respectively), or one of those given by (25), (26) and (27) if it belongs to family B.
We next summarize the results obtained under these hypotheses for metrics of both families A and B.
Family A:
Type $`I`$: In this case, the fluid velocity is parallel to the CKV $`X`$. This case has already been studied in full generality (see ) showing that the only such space-times are locally FRW models (assuming $`\mu +p0`$ and $`p=p(\mu )`$) and therefore the $`C_3`$ is not maximal. Also, without making any assumptions on the equation of state, it can be easily seen that either the $`C_3`$ is not maximal or $`X`$ is not a proper CKV.
Type $`III`$: The only solution of this type which admits no further CKVs (including KVs and HVFs) is
$$ds^2=e^{2f}\left\{dt^2+dr^2+\frac{e^{2t}}{(\mathrm{cosh}r)^{2(A+1)}}dy^2+\frac{\mathrm{sinh}^2r}{(\mathrm{cosh}r)^{2(A+1)}}dz^2\right\},$$
(36)
where
$$f=\frac{A}{2}U\mathrm{ln}(\mathrm{cosh}kU)+A\mathrm{ln}(\mathrm{cosh}r)+\alpha ,$$
(37)
$`A`$ and $`\alpha `$ are constants, $`U=U(t,r)`$ is given by
$$U=t+\mathrm{ln}(\mathrm{cosh}r),$$
(38)
the constant $`k`$ takes the value,
$$k=\frac{\sqrt{(A+1)^2+1}}{2},$$
(39)
and $`\mu `$ and $`p`$ are given by,
$`\mu e^{2f}`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{cosh}^2r}}\{1+A+3(M1)^2\},`$ (40)
$`pe^{2f}`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{cosh}^2r}}\{1+2A3M^2+(42A)M\},`$ (41)
where $`M=M(U)`$ is defined by
$$M\frac{A}{2}k\mathrm{tanh}(kU).$$
(42)
The imposition of the energy conditions ($`\mu >0`$, $`\mu p\mu `$) results in restrictions on the values that the parameter $`A`$ can take, namely $`A(1,5)`$.
Type $`V`$: In this case $`X`$ cannot be a proper CKV. However, there are solutions belonging to this type for which $`X`$ is a HVF (see for instance ).
Type $`VI`$: No perfect fluid solutions of this type exist which admit a proper CKV, as can easily be seen from the field equations.
Family B:
All the solutions in this family which satisfy the conditions set up at the beginning of this subsection are either such that $`X`$ becomes a KV (metrics of the form (27)), or they admit two further KV which together with $`\xi `$ and $`\eta `$ generate a four parameter group of isometries $`G_4`$ acting on three-dimensional spacelike orbits, the space-time thus becoming spatially homogeneous (metrics of the form (25) and (26)). Furthermore, the subgroup $`G_3`$ that $`G_4`$ necessarily contains acts on two-dimensional spacelike orbits, coordinated by $`r`$ and $`z`$, which are then of positive constant curvature; hence the models are spherically symmetric Kantowski-Sachs space-times.
### 4.2 Fluid flow orthogonal to the conformal orbits.
The conformal orbits must be spacelike in this case, and the coordinates are comoving, i.e.
$$u_t=e^Fu_r=0.$$
(43)
Family A:
Type $`I`$: The only solutions belonging to this type which admit no further symmetry are
$$ds^2=\frac{f_o^2}{\mathrm{cosh}^2r}\left\{dt^2+dr^2+\frac{(\mathrm{cosh}t)^{\sigma +1}}{(\mathrm{sinh}t)^{\sigma 1}}dy^2+\frac{(\mathrm{sinh}t)^{\sigma +1}}{(\mathrm{cosh}t)^{\sigma 1}}dz^2\right\},$$
(44)
where $`f_o`$ is an arbitrary constant and the parameter $`\sigma `$ must be such that $`|\sigma |<1`$; the energy density and pressure are then given by
$`\mu `$ $`=`$ $`{\displaystyle \frac{1}{f_o^2}}\left\{(1\sigma ^2){\displaystyle \frac{\mathrm{cosh}^2r}{\mathrm{sinh}^22t}}+3\right\},`$ (45)
$`p`$ $`=`$ $`{\displaystyle \frac{1}{f_o^2}}\left\{(1\sigma ^2){\displaystyle \frac{\mathrm{cosh}^2r}{\mathrm{sinh}^22t}}3\right\}.`$ (46)
The equation of state is simply $`p=\mu 6/f_o^2`$ as can be easily seen from (45) and (46). Notice that this solution is only valid in the region $`t>0`$ and that it is separable in the variables $`t`$ and $`r`$; and hence it must be contained in the solutions given by Ruiz and Senovilla .
Type $`III`$: No perfect fluid solutions admitting a proper CKV exist, since $`\mu +p=0`$.
Type $`V`$: There are no solutions of this type for perfect fluids with $`\mu >0`$ and the conditions set up above, as can easily be seen from the field equations.
Type $`VI`$: In this case, the line element, density and pressure take the following forms:
$$ds^2=e^{2f(t,r)}\left\{dt^2+dr^2+Q(t)\left[\frac{e^{2r}}{h^2(t)}dy^2+e^{2qr}h^2(t)dz^2\right]\right\},$$
(47)
$$\mu e^{2f}=(\dot{\beta }^21)\left\{\frac{3m^2}{(1+e^{mu})^2}+qn^2\frac{1}{\dot{\beta }^2}\right\},$$
(48)
$$pe^{2f}=(\dot{\beta }^21)\left\{qn^2\frac{1}{\dot{\beta }^2}\frac{3m^2}{(1+e^{mu})^2}+4q\frac{1+q^2}{(1+q)^2}\frac{1}{1+e^{mu}}\right\},$$
(49)
where $`q0,1`$ and
$$m\frac{1+q^2}{1+q}n\frac{1q}{1+q},$$
(50)
and
$$f=H(u)+\lambda (t)ur+\beta (t),$$
(51)
where $`H(u)`$ is given by
$$H=\mathrm{ln}(1+e^{mu}),$$
(52)
and there are two possibilities for $`\beta (t)`$, namely:
$$(\mathrm{i})\beta =\frac{1}{k}\mathrm{ln}|\mathrm{sinh}kt|,$$
(53)
with $`k\frac{(q1)^2}{q+1}`$, whence
$$Q=\mathrm{cosh}kth^2=(\mathrm{cosh}kt)^{1/n},$$
(54)
$$\lambda =\frac{1+q^2}{(1q)^2}\mathrm{ln}|\mathrm{sinh}kt|+\lambda _o,\lambda _o=const.$$
(55)
$$(\mathrm{ii})\beta =\frac{1}{k}\mathrm{ln}(\mathrm{cosh}kt),$$
(56)
whence
$$Q=\mathrm{sinh}kth^2=(\mathrm{sinh}kt)^{1/n},$$
(57)
$$\lambda =\frac{1+q^2}{(1q)^2}\mathrm{ln}(\mathrm{cosh}kt)+\lambda _o,\lambda _o=const.$$
(58)
Again, these solutions are only valid for $`t>0`$. It is also easy to see that the choice of a negative value for the parameter $`q`$ would imply that $`(\mu +p)`$ and $`(\mu p)`$ have opposite signs, thus violating one of the energy conditions (furthermore, $`\mu `$ could only be non-negative in a certain region of the space-time). On the other hand, assuming $`q>0`$ immediately leads to the choice (i) for $`\beta (t)`$ in order to have $`\mu >0`$, whence $`\mu +p`$ is always non-negative but there will always be some region of the space-time in which $`\mu p<0`$; consequently these solutions are only valid in a certain region of the space-time (in particular: $`(t,r)`$ satisfying $`\mathrm{sinh}|k|t<[(q^22q/3+1)/2q]^{k/m}e^{kr}`$).
Family B:
It is easy to see from the field equations that there are no perfect fluid solutions in this family.
## 5 Perfect fluid space-times: Tilted case.
The field equations in this case are (33) and (34), but we no longer have the additional condition $`u_r=0`$ as in the previous cases. However, it should be noted that for a perfect fluid solution admitting an Abelian $`G_2`$ of isometries acting on spacelike orbits and such that they admit orthogonal surfaces, it is always possible to perform a change of coordinates so as to bring the coordinates into a comoving form with respect to $`u^a`$, while leaving the metric diagonal . Such a coordinate change in the $`t,r`$ plane (in our coordinates) would dramatically change the form of the proper CKV $`X`$; for example, we would no longer be able to integrate out the conformal equations (3) nor provide simple expressions for the metric functions. Roughly speaking, one must choose between the coordinate chart adapted to the CKV and the one adapted to the four-velocity of the fluid. It is interesting to notice that this is not the case for HVFs; since one may change there to comoving coordinates the HVF changing then in a “controlled” way, so that one can still integrate out the homothetic equations (see for instance ). This difference is mainly due to the fact that the four-velocity field of a perfect fluid and a HVF are always surface forming (or in a more physical language, the fluid “inherits” the symmetry, see ), whereas this is not necessarily so in the case of a CKV.
We shall deal separately with the cases of timelike and spacelike conformal orbits and, in each case, we shall distinguish between the two families of metrics A and B.
### 5.1 Timelike conformal orbits.
Family A:
We shall discuss here some general features of the solutions belonging to different types in this family.
The metric can be written as:
$$ds^2=e^{2F(t,r)}\left\{dt^2+dr^2+e^{at}\frac{q(r)}{h^2(r)}dy^2+e^{bt}q(r)h^2(r)dz^2\right\},$$
(59)
where
| Type | $`a`$ | $`b`$ |
| --- | --- | --- |
| $`I`$ | 0 | 0 |
| $`III`$ | -2 | 0 |
| $`V`$ | -2 | -2 |
| $`VI`$ | -2 | -2q |
and $`q0,1`$.
For these cases Eqn. (33) gives:
$$\frac{a^2b^2}{4}2\left(\frac{h^{}}{h}\right)^2+2\frac{h^{}q^{}}{hq}+2\frac{h^{\prime \prime }}{h}+4\frac{h^{}}{h}F_{,r}+(ab)F_{,t}=0,$$
(60)
where a dash indicates differentiation with respect to $`r`$. At this point we must distinguish between two cases depending on whether $`a=b`$ or not.
Case (i): Assume $`a=b`$ (types $`I`$ and $`V`$). Equation (60) reads now:
$$0=\left(\frac{h^{}}{h}\right)^{}+\frac{h^{}q^{}}{hq}+2\frac{h^{}}{h}F_{,r}$$
(61)
since $`h^{}0`$ (otherwise the metric admits a further KV), this equation immediately gives
$$e^{2F}=f(t)\frac{h}{qh^{}}$$
(62)
Case (ii): If $`ab`$ (types $`III`$ and $`VI`$). Differentiating (60) with respect to $`t`$, we have
$$0=4\frac{h^{}}{h}F_{,rt}+(ab)F_{,tt},$$
(63)
that integrates to give
$$F=H(t+\beta (r))+\lambda (r),$$
(64)
with
$$\beta ^{}\frac{h^{}}{h}=\frac{ba}{4},$$
(65)
and the whole equation (60) reads
$$\frac{a^2b^2}{4}+2\left(\frac{h^{}}{h}\right)^{}+2\frac{h^{}q^{}}{hq}+4\frac{h^{}}{h}\lambda ^{}=0.$$
(66)
Let us deal next with the first case. For types $`I`$ and $`V`$ we can use an alternative form of the metric that satisfies (33) identically, in order to simplify the resulting expressions:
$$ds^2=\frac{e^{f(t)q(r)}}{h^{}(r)}\left\{dt^2+dr^2\right\}+\frac{e^{f(t)+at}}{h^{}(r)}\left\{e^{h(r)}dy^2+e^{h(r)}dz^2\right\},$$
(67)
the conformal factor now being:
$$\mathrm{\Psi }=\frac{f_{,t}}{2}.$$
(68)
Equation (34) can be written as
$$0=\mathrm{\Sigma }_0+a\mathrm{\Sigma }_1f_{,t}+\mathrm{\Sigma }_2f_{,t}^2+\frac{a}{4}f_{,t}^3+\mathrm{\Sigma }_3f_{,tt}\frac{a}{2}f_{,t}f_{,tt},$$
(69)
where $`\mathrm{\Sigma }_i`$, $`i=0,\mathrm{},3`$, denote functions depending only on $`r`$
$`\mathrm{\Sigma }_0`$ $`=`$ $`{\displaystyle \frac{a^2}{4}}(q^{})^2+{\displaystyle \frac{a^2}{4}}{\displaystyle \frac{h^{\prime \prime }}{h^{}}}q^{}+{\displaystyle \frac{1}{4}}h^{\prime \prime }h^{}q^{}{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{h^{\prime \prime }}{h^{}}}\right)^2(q^{})^2+{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{h^{\prime \prime }}{h^{}}}\right)^3q^{}{\displaystyle \frac{a^2}{4}}q^{\prime \prime }`$ (70)
$`{\displaystyle \frac{1}{4}}(h^{})^2q^{\prime \prime }{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{h^{\prime \prime }}{h^{}}}\right)^2q^{\prime \prime }+{\displaystyle \frac{1}{4}}(q^{\prime \prime })^2{\displaystyle \frac{1}{2}}{\displaystyle \frac{h^{\prime \prime \prime }h^{\prime \prime }}{(h^{})^2}}q^{}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{h^{\prime \prime \prime }}{h^{}}}q^{\prime \prime },`$
$`\mathrm{\Sigma }_1`$ $`=`$ $`{\displaystyle \frac{a^2}{4}}+{\displaystyle \frac{1}{4}}(h^{})^2+{\displaystyle \frac{1}{2}}(q^{})^2+{\displaystyle \frac{1}{2}}{\displaystyle \frac{h^{\prime \prime }}{h^{}}}q^{}+{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{h^{\prime \prime }}{h^{}}}\right)^2{\displaystyle \frac{1}{2}}q^{\prime \prime }{\displaystyle \frac{1}{2}}{\displaystyle \frac{h^{\prime \prime \prime }}{h^{}}},`$ (71)
$`\mathrm{\Sigma }_2`$ $`=`$ $`{\displaystyle \frac{a^2}{2}}+{\displaystyle \frac{1}{4}}(h^{})^2+{\displaystyle \frac{1}{4}}(q^{})^2+{\displaystyle \frac{1}{4}}{\displaystyle \frac{h^{\prime \prime }}{h^{}}}q^{}+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{h^{\prime \prime }}{h^{}}}\right)^2{\displaystyle \frac{1}{4}}q^{\prime \prime }{\displaystyle \frac{1}{2}}{\displaystyle \frac{h^{\prime \prime \prime }}{h^{}}},`$ (72)
$`\mathrm{\Sigma }_3`$ $`=`$ $`{\displaystyle \frac{a^2}{2}}{\displaystyle \frac{1}{2}}(h^{})^2+{\displaystyle \frac{1}{2}}{\displaystyle \frac{h^{\prime \prime }}{h^{}}}q^{}{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{h^{\prime \prime }}{h^{}}}\right)^2+{\displaystyle \frac{1}{2}}q^{\prime \prime }+{\displaystyle \frac{h^{\prime \prime \prime }}{h^{}}},`$ (73)
Since $`f_{,tt}0`$ (otherwise $`X`$ would become a HVF) the following possibilities arise from Eqn. (69) namely:
($`I.a`$) $`\mathrm{\Sigma }_3=0`$ then $`\mathrm{\Sigma }_2=\mathrm{\Sigma }_0=0`$ (for all $`f(t)`$ ).
($`I.b`$) $`\mathrm{\Sigma }_30`$ then $`\mathrm{\Sigma }_2=c\mathrm{\Sigma }_3`$, $`\mathrm{\Sigma }_0=\alpha \mathrm{\Sigma }_3`$,
$$0=\alpha c(f_{,t})^2+f_{tt}$$
(74)
where $`\alpha `$ and $`c`$ are constants.
($`V.a`$) $`\mathrm{\Sigma }_i=\alpha _i`$ (constants).
($`V.b`$)
$$0=e+df_{,t}+\frac{1}{2}(f_{,t})^2f_{,tt}$$
(75)
$$0=2\mathrm{\Sigma }_2+\mathrm{\Sigma }_3+2d,$$
(76)
$$0=\mathrm{\Sigma }_1+d\mathrm{\Sigma }_3+e,$$
(77)
$$0=\mathrm{\Sigma }_0+e\mathrm{\Sigma }_3,$$
(78)
where $`e`$, $`d`$ are constants.
For all these cases one has:
$`G_{tr}`$ $`=`$ $`{\displaystyle \frac{a}{2}}q^{}{\displaystyle \frac{1}{2}}f_{,t}\left(q^{}+{\displaystyle \frac{h^{\prime \prime }}{h^{}}}\right),`$ (79)
$`\mu g_{rr}`$ $`=`$ $`{\displaystyle \frac{3}{4}}a^2+{\displaystyle \frac{3}{2}}af_{,t}+{\displaystyle \frac{3}{4}}(f_{,t})^2+{\displaystyle \frac{1}{4}}(h^{})^2{\displaystyle \frac{3}{4}}\left({\displaystyle \frac{h^{\prime \prime }}{h^{}}}\right)^2{\displaystyle \frac{1}{2}}q^{\prime \prime },`$ (80)
$`(\mu p)g_{rr}`$ $`=`$ $`a^2+2af_{,t}+(f_{,t})^2+f_{,tt}2\left({\displaystyle \frac{h^{\prime \prime }}{h^{}}}\right)^2+{\displaystyle \frac{h^{\prime \prime \prime }}{h^{}}},`$ (81)
$`(\mu +p)g_{rr}`$ $`=`$ $`{\displaystyle \frac{a^2}{2}}+af_{,t}+{\displaystyle \frac{1}{2}}(f_{,t})^2f_{,tt}+{\displaystyle \frac{1}{2}}(h^{})^2+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{h^{\prime \prime }}{h^{}}}\right)^2q^{\prime \prime }{\displaystyle \frac{h^{\prime \prime \prime }}{h^{}}},`$ (82)
$`(\mu +p)u_t^2`$ $`=`$ $`{\displaystyle \frac{a}{2}}f_{,t}+{\displaystyle \frac{1}{2}}(f_{,t})^2f_{,tt}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{h^{\prime \prime }}{h^{}}}q^{}{\displaystyle \frac{1}{2}}q^{\prime \prime }.`$ (83)
Subcase $`I.a`$ $`a=0`$ and $`\mathrm{\Sigma }_0=\mathrm{\Sigma }_2=\mathrm{\Sigma }_3=0`$. It can be trivially checked that
$$q^{}+\frac{h^{\prime \prime }}{h^{}}=0,$$
(84)
and that $`G_{tr}=0`$, so the solution has a comoving fluid flow and, as was pointed out before, it has more symmetry.
Subcase $`I.b`$ $`\mathrm{\Sigma }_2=c\mathrm{\Sigma }_3`$ yields
$$0=(12c)(h^{})^2+(q^{})^2+(1+2c)\frac{h^{\prime \prime }}{h^{}}q^{}+2(1c)\left(\frac{h^{\prime \prime }}{h^{}}\right)^2(12c)q^{\prime \prime }2(12c)\frac{h^{\prime \prime \prime }}{h^{}}.$$
(85)
For the particular case $`c=\frac{1}{2}`$ the above equation reads
$$\left(q^{}+\frac{h^{\prime \prime }}{h^{}}\right)^2=0,$$
(86)
and then $`\mathrm{\Sigma }_0=\alpha \mathrm{\Sigma }_3`$ implies
$$\frac{1}{2}\frac{h^{\prime \prime \prime }}{h^{}}\left(\frac{h^{\prime \prime }}{h^{}}\right)^2=\alpha ,$$
(87)
and one can see from Eqn. (83) that $`(\mu +p)u_t^2=0`$, so, this solution cannot represent a perfect fluid. For $`c\frac{1}{2}`$, we obtain, after some calculations involving (85)
$`\mathrm{\Sigma }_0`$ $`=`$ $`{\displaystyle \frac{1}{4(12c)}}\left[q^{\prime \prime }q^{}{\displaystyle \frac{h^{\prime \prime }}{h^{}}}\right]\left[q^{}+{\displaystyle \frac{h^{\prime \prime }}{h^{}}}\right]^2,`$ (88)
$`\mathrm{\Sigma }_3`$ $`=`$ $`{\displaystyle \frac{1}{2(12c)}}\left[q^{}+{\displaystyle \frac{h^{\prime \prime }}{h^{}}}\right]^2.`$ (89)
The remaining equation $`\mathrm{\Sigma }_0=\alpha \mathrm{\Sigma }_3`$ can be written as
$$q^{\prime \prime }q^{}\frac{h^{\prime \prime }}{h^{}}=2\alpha .$$
(90)
Defining now
$$\sigma ^{}\frac{1}{h^{}},$$
(91)
equation (90) can be integrated and (85) rewritten to give
$$q^{}=2\alpha \frac{\sigma }{\sigma ^{}}+\frac{k}{\sigma ^{}},k=\mathrm{const},$$
(92)
$$0=\sigma ^{}\sigma ^{\prime \prime \prime }+\frac{1+3c}{12c}(\sigma ^{\prime \prime })^2\frac{2c}{12c}\sigma ^{\prime \prime }(2\alpha \sigma +k)\alpha (\sigma ^{})^2+\frac{1}{2(12c)}(2\alpha \sigma +k)^2+\frac{1}{2},$$
(93)
and one then has
$`(\mu +p)u_t^2`$ $`=`$ $`{\displaystyle \frac{12c}{2}}(f_{,t})^2,`$ (94)
$`(\mu +p)g_{rr}`$ $`=`$ $`{\displaystyle \frac{12c}{2}}(f_{,t})^2{\displaystyle \frac{1}{2(12c)}}\left[2\alpha {\displaystyle \frac{\sigma }{\sigma ^{}}}+{\displaystyle \frac{k}{\sigma ^{}}}{\displaystyle \frac{\sigma ^{\prime \prime }}{\sigma ^{}}}\right]^2,`$ (95)
$`(\mu p)g_{rr}`$ $`=`$ $`(1+c)(f_{,t})^2\alpha {\displaystyle \frac{\sigma ^{\prime \prime \prime }}{\sigma ^{}}}.`$ (96)
In order to satisfy the energy conditions in all regions of the space-time, the right-hand sides of the above equations must all be positive; necessary conditions for this are $`c<\frac{1}{2}`$ and also that $`(f_{,t})^2`$ has a strictly lower bound. The latter condition singles out one of the possible solutions to (74) namely:
$$\alpha =c\beta ^2f=\frac{1}{c}\mathrm{ln}|\mathrm{sinh}(c\beta t+C_1)|.$$
(97)
where another additive constant of integration in the expression of $`f`$ has been set equal to zero without loss of generality. Notice that $`C_1`$ can also be made zero by suitably redefining $`t`$, nevertheless we choose to maintain it in order to allow singularity-free solutions when $`t=0`$.
One particular solution to (93) is
$$\sigma =\frac{\mathrm{cosh}\beta r}{\beta ^2}\mathrm{and}k=12c.$$
(98)
Notice that in order to have the correct signature in the metric $`r`$ is restricted to positive values only. The energy condition (96) is satisfied for all values of $`r`$ but only if $`c>1`$.
Type $`V`$: After some calculations, it can be seen that the only possibility for the line element given by (67) leading to an Einstein tensor of the perfect fluid type is:
$$ds^2=e^{f(t)2ar}\{dt^2+dr^2\}+e^{f(t)2t}\{e^{mr}dy^2+e^{mr}dz^2\}$$
(99)
where $`a`$ and $`m`$ are constants. The only remaining field equation is then
$$0=4a^2(2+4a^2+\frac{m^2}{2})f_{,t}+(2+a^2+\frac{m^2}{4})(f_{,t})^2\frac{1}{2}(f_{,t})^3+f_{,tt}(f_{,t}2\frac{m^2}{2}),$$
(100)
that integrates to give
$$ke^t=(f_{,t}2)(f_{,t}p^2(1+s))^{(1s)/2s}(f_{,t}p^2(1s))^{(1+s)/2s},$$
(101)
where the constants $`p^2`$ and $`s`$ are defined in terms of $`a`$ and $`m`$ by
$$p^2=a^2+\frac{m^2}{4}+1,s^2=1\frac{4a^2}{p^4}.$$
(102)
The density, pressure and velocity field are given by
$`(\mu +p)u_t^2`$ $`=`$ $`{\displaystyle \frac{a^2(f_{,t}2)^2}{f_{,t}2(p^2a^2)}},`$ (103)
$`(\mu +p)g_{rr}`$ $`=`$ $`(f_{,t}2(p^2a^2))+{\displaystyle \frac{a^2(f_{,t}2)^2}{f_{,t}2(p^2a^2)}},`$ (104)
$`(\mu p)g_{rr}`$ $`=`$ $`{\displaystyle \frac{(f_{,t}2)(\frac{3}{2}(f_{,t})^2+(3p^2+2a^22)f_{,t}+4p^22a^2)}{f_{,t}2(p^2a^2)}}.`$ (105)
The energy conditions $`\mu \pm p0`$ will be satisfied if and only if the right-hand sides of the above equations are positive, i.e.:
$`f_{,t}2(p^2a^2)>0,`$ (106)
$`(a^21)(f_{,t})^2+(4p^28a^2)f_{,t}+4a^24(p^2a^2)^2>0,`$ (107)
$`{\displaystyle \frac{3}{2}}(f_{,t})^2+(3p^2+2a^22)f_{,t}+4p^22a^20,`$ (108)
and these, in turn, can be seen to imply (after a careful analysis):
$`a^210`$ $`f_{,t}>\beta `$ (109)
$`a^21<0`$ $`\beta <f_{,t}<2\left(1+{\displaystyle \frac{m^2}{4}}{\displaystyle \frac{1}{1|a|}}\right)`$ (110)
where
$$\beta \frac{1}{3}\left[a^2+3\frac{m^2}{4}+5+\sqrt{(a^21)^2+3\frac{m^2}{4}(2+2a^2+3\frac{m^2}{4})}\right].$$
(111)
From (103)-(105) it is easy to see that a barotropic equation of state of the form $`p=p(\mu )`$ is not possible.
The analysis of types $`III`$ and $`VI`$ turns out to be very involved, with many subcases arising, and so far, we have not been able to find any solution to the EFEs which is valid (i.e., satisfies energy conditions) all over the space-time manifold. Since we cannot provide any result in the positive we shall not discuss these cases any further here, so as to keep the present study at a reasonabe length.
Family B:
For this family, the metric and CKV $`X`$ take one of the forms given by (19)-(21), and the conformal factor $`\mathrm{\Psi }`$ is then given by (22).
From the field equations for a perfect fluid it is easy to see that the metrics (20) and (21) cannot represent a perfect fluid space-time with $`(\mu +p)0`$ and $`\mu >0`$; thus, only (19) needs be considered. The field equations (33)-(34) for this metric suggest redefinitions of the coordinates $`t=\varphi _1(t)`$ and $`r=\varphi _2(r)`$, so that the metric takes the form (in the new coordinates):
$$ds^2=\frac{e^{2f(tr)}}{\sqrt{H^{}}}\left\{\frac{dt^2}{t}+\frac{H^{}}{H}dr^2+4tdy^2+Hdz^2\right\},$$
(112)
where $`H=H(r)`$, as usual a dash indicates a derivative with respect to $`r`$ and $`f=f(tr)`$ is a function of $`trx`$. The CKV $`X`$ and conformal factor $`\mathrm{\Psi }`$ are in these coordinates:
$$X=e^y\left(t^{1/2}_t\frac{1}{2}t^{1/2}_y\right),$$
(113)
$$\mathrm{\Psi }=e^y\sqrt{t}f_{,x}.$$
(114)
The remaining field equations are now:
$$(f_{,x}^3f_{,x}f_{,xx})\mathrm{\Sigma }_0(r)+f_{,xx}\mathrm{\Sigma }_1(r)+f_{,x}^2\mathrm{\Sigma }_2(r)=0,$$
(115)
where
$$\mathrm{\Sigma }_0(r)2\frac{H^{\prime \prime }}{H^{}},\mathrm{\Sigma }_1(r)\frac{5}{4}\left(\frac{H^{\prime \prime }}{H^{}}\right)^2+\frac{H^{\prime \prime \prime }}{H^{}},$$
(116)
$$\mathrm{\Sigma }_2(r)\frac{3}{2}\left(\frac{H^{\prime \prime }}{H^{}}\right)^2\frac{H^{\prime \prime \prime }}{H^{}},$$
(117)
and they satisfy: $`\mathrm{\Sigma }_1+\mathrm{\Sigma }_2=\frac{\mathrm{\Sigma }_0^2}{16}`$.
Also notice that one of the trivial solutions to (115), namely $`f_{,x}=0`$ would correspond to $`X`$ being a KV, on account of the form of the conformal factor (114) and that the case $`f_{,x}=`$constant leads to an incorrect signature of the metric. Excluding these cases, $`\mathrm{\Sigma }_0`$, $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_2`$ must be constants, and the following possibilities then arise:
(B-1) $`\mathrm{\Sigma }_0=0`$, then necessarily $`\mathrm{\Sigma }_1=\mathrm{\Sigma }_2=0`$, and
$$H=K^2r+C,K,C=\mathrm{const},$$
(118)
where we can always make $`C=0`$ by suitably redefining $`r`$, and $`f`$ is then a completely arbitrary function of its argument. The density, pressure and velocity field can then be obtained from:
$`\mu {\displaystyle \frac{g_{yy}}{4t}}`$ $`=`$ $`3\{f_{,x}+(tr)f_{,x}^2\},`$ (119)
$`(\mu +p){\displaystyle \frac{g_{yy}}{4t}}`$ $`=`$ $`2(tr)\{f_{,x}^2f_{,xx}\},`$ (120)
$`(\mu p){\displaystyle \frac{g_{yy}}{4t}}`$ $`=`$ $`2\{(tr)f_{,xx}+2(tr)f_{,x}^2+3f_{,x}\},`$ (121)
$`(\mu +p)u_t^2`$ $`=`$ $`2\{f_{,x}^2f_{,xx}\}.`$ (122)
From the above expressions one can readily see that no function $`f(x)`$ exists such that $`\mu `$, $`\mu +p`$ and $`\mu p`$ are positive over all of the space-time.
(B-2) $`\mathrm{\Sigma }_00`$. This leads to
$$H=K^2\{e^{4sr}+h_0\},$$
(123)
where $`K`$, $`s`$ and $`h_0`$ are arbitrary constants. The function $`f`$ is then given implicitly by
$$e^{s/f_{,x}}\left\{1+\frac{s}{f_{,x}}\right\}=e^{2sx},$$
(124)
and the density and pressure can be obtained from:
$`\mu {\displaystyle \frac{g_{yy}}{4t}}`$ $`=`$ $`{\displaystyle \frac{3e^{4sr}}{4s}}\left\{e^{4sr}\left[f_{,x}^2(4st1)+2sf_{,x}+{\displaystyle \frac{s^2}{3}}\right]h_0(s+f_{,x})^2\right\},`$ (125)
$`(\mu +p){\displaystyle \frac{g_{yy}}{4t}}`$ $`=`$ $`{\displaystyle \frac{e^{4sr}}{2(2f_{,x}+s)}}\{e^{4sr}[f_{,x}^2(4st1)+2sf_{,x}+s^2]`$
$`+`$ $`h_0(s+f_{,x})^2\},`$ (126)
$`(\mu p){\displaystyle \frac{g_{yy}}{4t}}`$ $`=`$ $`{\displaystyle \frac{e^{4sr}}{s(2f_{,x}+s)}}\{3e^{4sr}f_{,x}[f_{,x}^2(4st1)+{\displaystyle \frac{4s}{3}}(2st+1)f_{,x}+s^2]`$
$``$ $`h_0[3f_{,x}^3+8sf_{,x}^2+7s^2f_{,x}+2s^3]\}.`$ (127)
The energy conditions $`\mu >0`$, $`\mu \pm p>0`$ will be satisfied if and only if the above expressions are all positive, and a careful analysis of these conditions taking into account (124) shows that they can only hold in some open domains of the space-time manifold (e.g., for $`h_0=0`$ the energy conditions can only hold in the region $`tr<0`$ if $`s<0`$, and in the region $`𝒦^2(t)<tr<0`$ if $`s>0`$, $`𝒦`$ being some function of $`t`$).
### 5.2 Spacelike conformal orbits.
Family A:
The forms of the metric functions are then (11), (13), (15) and (16), with $`w=0`$ and the coordinates $`t`$ and $`r`$ interchanged. We note that the solutions are not simply obtained from (5.1) by interchanging $`t`$ and $`r`$. We next briefly discuss a few solutions belonging to some of the Bianchi types.
Type $`I`$: The analysis of this type follows much along the same lines as its counterpart in the case of timelike conformal orbits. Thus we have
$$ds^2=\frac{e^{f(r)q(t)}}{\dot{h}(t)}\left\{dt^2+dr^2\right\}+\frac{e^{f(r)}}{\dot{h}(t)}\left\{e^{h(t)}dy^2+e^{h(t)}dz^2\right\},$$
(128)
where a dot indicates differentiation with respect to $`t`$. Then the independent Einstein equation takes the form:
$$0=\widehat{\mathrm{\Sigma }}_0+\widehat{\mathrm{\Sigma }}_2(f_{,r})^2+\widehat{\mathrm{\Sigma }}_3f_{,rr},$$
(129)
where the expressions of the $`\widehat{\mathrm{\Sigma }}_i`$s can be formally obtained from those of the $`\mathrm{\Sigma }_i`$s given by Eqns. (70), (72) and (73), by simply changing $`r`$ by $`t`$ there (i.e., changing primes into dots).
Since $`f_{,rr}`$ is non-null, the following possibilities arise:
($`I.a`$) $`\widehat{\mathrm{\Sigma }}_3=0`$, then necessarily $`\widehat{\mathrm{\Sigma }}_0=\widehat{\mathrm{\Sigma }}_1=0`$. In this subcase there are no perfect fluid solutions because the Einstein tensor do not have any timelike eigenvector.
($`I.b`$) $`\widehat{\mathrm{\Sigma }}_30`$, then $`\widehat{\mathrm{\Sigma }}_2=c\widehat{\mathrm{\Sigma }}_3`$, $`\widehat{\mathrm{\Sigma }}_0=\alpha \widehat{\mathrm{\Sigma }}_3`$ and
$$f_{,rr}=c(f_{,r})^2\alpha ,$$
(130)
where $`c`$ and $`\alpha `$ are arbitrary constants. The first equation yields
$$0=(12c)(\dot{h})^2+(\dot{q})^2+(1+2c)\frac{\ddot{h}}{\dot{h}}\dot{q}+2(1c)\left(\frac{\ddot{h}}{\dot{h}}\right)^2(12c)\ddot{q}2(12c)\frac{h^{\mathrm{}}}{\dot{h}}$$
(131)
The case $`c=\frac{1}{2}`$ can be easily seen to correspond to one of the previously studied cases, namely those with the fluid flow orthogonal to the conformal orbits.
Following now a similar procedure to the one outlined in the case of timelike conformal orbits, the field equations can then be written as
$$\dot{q}=2\alpha \frac{\sigma }{\dot{\sigma }}+\frac{k}{\dot{\sigma }},k=\mathrm{const},$$
(132)
$$0=\dot{\sigma }\sigma ^{\mathrm{}}+\frac{1+3c}{12c}(\ddot{\sigma })^2\frac{2c}{12c}\ddot{\sigma }(2\alpha \sigma +k)\alpha (\dot{\sigma })^2+\frac{1}{2(12c)}(2\alpha \sigma +k)^2+\frac{1}{2},$$
(133)
with $`\sigma (t)`$ defined through
$$\dot{\sigma }=\frac{1}{\dot{h}}.$$
(134)
Explicit expressions for $`\mu `$, $`p`$ and $`u_t`$ can be readily derived from:
$`(\mu +p)u_t^2`$ $`=`$ $`{\displaystyle \frac{1}{2(12c)}}\left[2\alpha {\displaystyle \frac{\sigma }{\dot{\sigma }}}+{\displaystyle \frac{k}{\dot{\sigma }}}{\displaystyle \frac{\ddot{\sigma }}{\dot{\sigma }}}\right]^2,`$ (135)
$`(\mu p)g_{rr}`$ $`=`$ $`(1+c)(f_{,r})^2+\alpha +{\displaystyle \frac{\sigma ^{\mathrm{}}}{\dot{\sigma }}},`$ (136)
$`(\mu +p)g_{rr}`$ $`=`$ $`{\displaystyle \frac{1}{2(12c)}}\left[2\alpha {\displaystyle \frac{\sigma }{\dot{\sigma }}}+{\displaystyle \frac{k}{\dot{\sigma }}}{\displaystyle \frac{\ddot{\sigma }}{\dot{\sigma }}}\right]^2{\displaystyle \frac{12c}{2}}(f_{,r})^2.`$ (137)
Now, from all the possible solutions for $`f(r)`$ to equation (130), the energy conditions $`\mu \pm p>0`$ single out
$$\alpha =c\beta ^2,f=\frac{1}{c}\mathrm{ln}(\mathrm{cosh}c\beta r)$$
(138)
and restrict $`c`$ to values $`c<1/2`$.
A particular solution to (133) which satisfies the energy conditions is:
$$\sigma =\frac{\mathrm{cosh}\beta t}{\beta ^2},k=(12c),$$
(139)
and we notice that the solution is valid for $`c>1`$ and in order for the metric to have the correct signature, $`t`$ must be positive.
Type $`V`$: There are no perfect fluid solutions of this type.
Family B:
An inspection of the field equations suggests, as in the case of timelike conformal orbits, a redefinition of $`t`$ and $`r`$, so that in the new coordinates the metric, CKV $`X`$ and the conformal factor $`\mathrm{\Psi }`$ take the following forms:
$$ds^2=\frac{2e^{2f(x)}e^{kt/2}}{\sqrt{\dot{H}}}\left\{\frac{1}{4}\frac{\dot{H}}{H}dt^2+\frac{k}{4}\frac{ce^{kr}}{1+ce^{kr}}dr^2+\frac{1+ce^{kr}}{k}dy^2+Hdz^2\right\},$$
(140)
$$X=e^y\left\{e^{kr/2}\sqrt{1+ce^{kr}}_r\frac{k}{2}\frac{ce^{kr/2}}{\sqrt{1+ce^{kr}}}_y\right\},$$
(141)
$$\mathrm{\Psi }=e^yf_{,x}e^{kr/2}\sqrt{1+ce^{kr}},$$
(142)
where $`k`$ and $`c`$ are constants and $`x`$ is defined as $`xtr`$. The field equations then reduce to:
$$2(f_{,x}f_{,xx}f_{,x}^3)\mathrm{\Sigma }_0+f_{,x}^2\mathrm{\Sigma }_1+(f_{,xx}+\frac{k}{2}f_{,x})\mathrm{\Sigma }_2=0,$$
(143)
where $`\mathrm{\Sigma }_0(t)`$, $`\mathrm{\Sigma }_1(t)`$ and $`\mathrm{\Sigma }_2(t)`$ are given by:
$$\mathrm{\Sigma }_0k+\frac{\ddot{H}}{\dot{H}},$$
(144)
$$\mathrm{\Sigma }_1k^2+\frac{k}{2}\frac{\ddot{H}}{\dot{H}}+\frac{1}{2}\left(\frac{\ddot{H}}{\dot{H}}\right)^2\left(\frac{\ddot{H}}{\dot{H}}\right)^{}=\dot{\mathrm{\Sigma }}_0+\frac{1}{2}\mathrm{\Sigma }_{0}^{}{}_{}{}^{2}+\frac{3}{2}k\mathrm{\Sigma }_0,$$
(145)
$$\mathrm{\Sigma }_2\frac{k^2}{4}\frac{1}{4}\left(\frac{\ddot{H}}{\dot{H}}\right)^2+\left(\frac{\ddot{H}}{\dot{H}}\right)^{}=\dot{\mathrm{\Sigma }}_0\frac{1}{4}\mathrm{\Sigma }_0(\mathrm{\Sigma }_0+2k).$$
(146)
The possibility $`f_{,x}=const0`$ leads to the wrong Segre type of the Einstein tensor (i.e., no perfect fluid solutions exist). Then, from Eqn.(143) it is immediate that either $`f_{,x}=0`$ (and then $`X`$ becomes a KV; see (142) ) or $`\mathrm{\Sigma }_0`$ and $`\mathrm{\Sigma }_1`$ are both constants, which, in turn, implies:
$$H=Me^{dt}+m,$$
(147)
where $`M`$, $`m`$ and $`d`$ are constants. Substituting this back into the field equation (143) we obtain:
$$2(d^2k^2)f_{,xx}+16(dk)(f_{,x}f_{,xx}f_{,x}^3)+4(d^2k^2+k(dk))f_{,x}^2k(d^2k^2)f_{,x}=0,$$
(148)
and we can distinguish the following cases:
(i) $`f_{,xx}=0`$; i.e., $`f_{,x}=`$constant, and from Eqn.(148) we obtain
$$f_{,x}=\frac{k+d}{4}\mathrm{or}f_{,x}=\frac{k}{4}.$$
(149)
In the first case ($`f_{,x}=\frac{k+d}{4}`$) the Segre type does not correspond to a perfect fluid (since $`u_t^2=0`$) and in the second case $`(\mu +p)u_t^2<0`$, so this cannot represent a physical perfect fluid.
(ii) $`f_{,xx}0`$; whence Eqn.(148) can be integrated once, leading to an implicit expression for $`f_{,x}`$. A careful analysis of the energy conditions $`\mu \pm p0`$ shows that in the general case (i.e., no assumptions on the values of the parameters $`k`$ and $`d`$), they can only be satisfied over certain restricted open domains of the space-time. However, two special cases arise which deserve special attention; namely, $`k=d`$ and $`k=d`$. For $`k=d`$, Eqn.(148) is identically satisfied, $`f`$ thus being a completely arbitrary function; however, this does not correspond to a perfect fluid (it is of the wrong Segre type). In the second case, Eqn.(148) can be integrated to give:
$$f=\mathrm{ln}\frac{e^{kx/4}}{\frac{4}{k}+ae^{kx/4}},$$
(150)
$`a`$ being a constant, and $`\mu `$, $`p`$ and $`u_t`$ can then be obtained from:
$`(\mu +p)`$ $`=`$ $`{\displaystyle \frac{e^{kx/2}}{c\sqrt{kM}}}(4+ake^{kx/4})(cme^{kt}Me^{kr}2cM),`$ (151)
$`(\mu p)`$ $`=`$ $`{\displaystyle \frac{e^{kx/2}}{c\sqrt{kM}}}[(16+ake^{kx/4})(cme^{kt}Me^{kr})`$ (152)
$`+ackMe^{kx/4}(2ake^{kx/4})],`$
$`(\mu +p)u_t^2`$ $`=`$ $`{\displaystyle \frac{k^2}{2(4+ake^{kx/4})}}.`$ (153)
As usual, the energy conditions will be fulfilled if and only if the above expressions are all positive and, again, this can only be possible over certain open domains of the manifold (notice that $`k`$ must be positive in order for the metric to have the correct signature).
## 6 Discussion
We have studied perfect fluid space-times admitting a three-dimensional Lie group of conformal motions containing a two-dimensional Abelian subgroup of isometries. All such space-times have been classified geometrically and in each class the metric has been explicitly given in coordinates adapted to the symmetry vectors.
In section 3 we restricted attention to the diagonal case, and in section 4 we found all such perfect fluid solutions in which the fluid four-velocity is tangential or orthogonal to the conformal orbits. In the former case the orbits are necessarily timelike and the only solutions for which $`𝒞_3`$ is maximal and $`X`$ is proper are of type $`III`$ in family A and given by Eqns. (36)-(42). In the latter case, in which $`u^a`$ is orthogonal to spacelike conformal orbits, there is a 2-parameter family of solutions given by (44)-(46) of type $`I`$ which are valid for $`t>0`$ and are separable in $`t`$ and $`r`$ (see ); and a class of solutions of type $`VI`$ \[see Eqns. (47)-(52) and Eqns. (53)-(55) or (56)-(58)\]. In this latter case, physical constraints (e.g., the energy conditions) restrict the validity of the solutions to a region of the space-time and, again, there are no solutions in family B.
In section 5 perfect fluid solutions were sought in the general (tilting) case in which the fluid four-velocity is neither tangential to nor orthogonal to the conformal orbits. We chose to work in coordinates adapted to the CKV and again the two cases in which the conformal orbits are timelike (subsection 5.1) or spacelike (subsection 5.2) were considered.
In the timelike case, solutions in family A of types $`I`$ \[metric (67), field eqns. (91)-(96)\] and $`V`$ \[metric (99), field eqns. (103)-(105)\] were obtained. In both cases solutions exist such that the energy conditions are satisfied on the whole space-time manifold; a particular solution of type $`I`$ was given by Eqns. (97) and (98). We noted that solutions of type $`V`$ cannot admit an equation of state of the form $`p=p(\mu )`$. In the spacelike case, solutions in family A of type $`I`$ were found and a particular solution, given by Eqn. (139), in which the energy conditions are always satisfied, was displayed. There are no solutions of type $`V`$ in this case.
All solutions in family B can be found in both the timelike case \[see Eqns. (112)-(124)\] and the spacelike case \[see Eqns. (140)-(150)\]. However, our analysis showed that in general there are no solutions for which the energy conditions are satisfied over the entire space-time manifold.
The case of null conformal orbits will be studied in a future paper.
Acknowledgments
The authors would like to thank M. Mars (Univ. of Barcelona) for many helpful discussions and suggestions. They would also like to thank the referees for many helpful suggestions which have greatly contributed to making the paper more precise and readable. One of the authors (AAC) would like to thank the Spanish Government and the “Dirección General de Investigación Científica y Técnica (DGICYT)” for financial assistance and the Department of Physics at the Universitat de les Illes Balears for its warm hospitality. This work is partially supported by the DGICYT Project No. PB 91-0335. Financial support from STRIDE program (Research Project No. STRDB/C/CEN/509/92) is also acknowledged.
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# 1 Introduction
## 1 Introduction
Integrable theories in two-dimensions have been an extraordinary laboratory for the understanding of basic nonperturbative aspects of physical theories and various phenomena such as dynamical mass generation, asymptotic freedom, and quark confinement, relevant in more realistic models, have been tested. In QCD<sub>2</sub> one of the open questions is related to the higher order correction to the limit $`m/e0`$ ($`m`$, quark mass and $`e`$, coupling constant). A speculation regarding this issue was set forward earlier (see, for example, Ref. ), that the low-energy effective action of QCD<sub>2</sub> might be related to some massive two dimensional integrable models, thus leading to the exact solution (not semi-classical) of the strong coupled QCD<sub>2</sub>. Although some hints toward a possible integrable structure in QCD<sub>2</sub> have been encountered the problem still remains largely open (see, e.g., ). The problem resides, for example, in a better understanding of the meson spectrum of the theory, for finite number of colors ($`N_c`$). It is believed that integrability of a theory implies stability of its bound states (mesonic spectrum). Therefore, one could expect vanishing decay amplitudes for the mesonic states of the theory. Several numerical computations indicate the contrary .
In the present paper we study the recently proposed conformal affine $`sl(2)^{(1)}`$ Toda model coupled to (Dirac) matter field (CATM) which is an example of a wide class of integrable theories presented in . The zero curvature representation, the construction of the general solution and many other properties are discussed in . This model possesses a Noether current depending only on the matter field and under some circunstances, it is possible to choose one solution in each orbit of the conformal group such that, for these solutions, the $`U(1)`$ current is equal to a topological current depending only on the Toda field. Such equivalence leads, at the classical level, to the localization of the (Dirac) matter field inside the Toda field soliton. It is well known the relevance of localized classical solutions of non-linear relativistic field equations to the corresponding quantum theories. In particular, solitons can be associated with quantum extended-particle states; this picture is valid for the soliton of our model, which at the classical level is associated to the Toda field $`\phi `$. Then one may regard the model a sort of one dimensional bag model for QCD. Besides, an additional feature is present; the masses of solitons and particles are proportional to the $`U(1)`$ Noether charge. This fact indicates the existence of a sort of duality in these models involving solitons and particles . The interest in such models comes from their integrability and duality properties , which can be used as a toy model to understand the electric-magnetic duality in four dimensional gauge theories, conjectured in and developed in . Thereby nonperturbative analysis of the spectrum and of the phase structure in SUSY Yang-Mills theory becomes possible.
The CATM theory is well defined mathematically for a set of fields which, in general, may be complex fields giving rise to a complex Lagrangian. By imposing convenient reality conditions on the fields of the model it is possible to define an integrable sub-model with a real Lagrangian. Moreover, there exists a related off-critical model, the so called affine $`sl(2)^{(1)}`$ Toda model coupled to the matter (ATM), which can be obtained at the classical or quantum mechanical level through some convenient reduction process starting from the CATM model . Recently the classical one and two soliton solutions of the model have been found, as well as the time-delays arising from the collisions of two solitons, and the implications of the reality conditions (imposed on the fields of the CATM model) on the solitonic solutions have been studied . In addition, in Ref. it was suggested the possibility of implementing a confinement mechanism of some degrees of freedom of the ATM model inside the baryons (solitons) of the sine-Gordon type reduced model. The symplectic structure of the off-critical ATM model has recently been studied . It was performed in the context of Faddeev-Jackiw and (constrained) symplectic methods (Barcelos-Neto, Montani and Wotzasek); by imposing the equivalence between the Noether and topological currents as a constraint, the authors have been able to obtain either, the sine-Gordon model or the massive Thirring model, through a Hamiltonian reduction and gauge fixing the symmetries of the model in two different ways.
Here we perform the reduction process of the CATM model to the off-critical ATM model through a BRST analysis and study the quantum spectrum of this reduced model. By bosonizing the spinor field we show that the off-critical ATM model is equivalent to a theory of a free massless scalar and a sine-Gordon field. A further quantum reduction process is implemented, this time to eliminate from the spectrum the modes associated to the free massless scalar; this amounts to imposing the equivalence between the Noether and topological currents as a constraint. Then it is shown that in the solitons (baryons) of the sine-Gordon theory the confinement of some degrees of freedom of the ATM model does take place.
Moreover, by means of a semi-classical analysis we find the “color” intercharge potential in the ATM model, obtaining a linearly rising potential for large intercharge separation, thus revealing the presence of a confining phase in the model. This behaviour leads to an interesting analogy with what one expects to happen in QCD. Then we suggest an interesting connection of the ATM model with the baryonic sector of the low-energy effective action of QCD<sub>2</sub> with one flavor ($`N_f=1`$) and two colors ($`N_c=2`$) . With that motivation we reproduce the bosonized form of QCD<sub>2</sub> in Baluni’s gauge . This gauge gives a convenient framework to study the topological realization of color symmetry in QCD<sub>2</sub>, and then reproduces the bosonized form of the theory in a form adequate to our discussion in connection to the bosonized formulation of the ATM model.
The paper is organized as follows. In section 2 we present some relevant aspects of the conformal affine $`sl(2)^{(1)}`$ Toda model coupled to matter (Dirac) field. Section 3 presents a Hamiltonian reduction of the two-loop WZNW model to obtain the equations of motion of the CATM theory. Section 4 presents a submodel in which a suitable reality conditions on the fields of CATM are imposed so as to define positive-definite kinetic term in the scalar sector and a usual spinorial kinetic term of the Lagrangian. These conditions will give rise to an off-critical model (ATM) with real Lagrangian, and the relevant Hamiltonian bounded from below. Also for the sake of completeness, we reproduce some relevant steps of the classical reduction process: CATM $``$ ATM . Moreover we present the general soliton (antisoliton) solution of the real Lagrangian submodel. In subsection 4.1 we present the masses of the fundamental particles and solitons, using a linear field approximation of the equations of motion. Section 5 deals with the quantum aspects of the model. In subsection 5.1 through a BRST analysis the conformal symmetry of the CATM is spontaneously broken, then defining the off-critical ATM model. In subsection 5.2, using bosonization techniques we arrive at the sine-Gordon model (massive Thirring) plus a free massless scalar field. A further quantum reduction is performed imposing the currents (Noether and topological) equivalence, then obtaining the sine-Gordon (massive Thirring) model only. In section 6 we present the realization of color symmetry in QCD<sub>2</sub> bosonizing the theory such that the topological confinement mechanism is more transparent. In section 7 a confinement mechanism of the ‘color’ sector of the ATM model inside the soliton (baryon) is explained observing the behaviour of the static intercharge potential.
Finally, section 8 summarise the main point and contains some remarks about the results. The appendix A provides the relevant notations and conventions, and appendix B presents some useful results concerning the $`sl(2)^{(1)}`$ affine Lie algebra.
## 2 The conformal affine Toda model coupled to the matter (CATM): the principal gradation of $`sl(2)^{(1)}`$ case
We will be interested in the so called conformal affine Toda system coupled to the matter field, whose construction, in the particular case of the algebra $`sl(2)^{(1)}`$, using the parlance of the original reference, will be summarized in the following paragraphs.
We discuss the example associated with the principal gradation of the untwisted affine Kac-Moody algebra $`sl(2)^{(1)}`$ <sup>1</sup><sup>1</sup>1In Ref. the terminology “affine Lie algebra” is preferred. I would like to thank Professor M.B. Halpern for correspondence, and for pointing out to me this terminology used in physics.. This belongs to a special class of models introduced in possessing a $`U(1)`$ Noether current depending only on the matter fields. It is then possible, under some circunstances, to choose one solution in each orbit of the conformal group, such that for these solutions, that $`U(1)`$ current is proportional to a topological current depending only on the (gauge) zero grade field.
The zero curvature condition in light-cone coordinates $`x_\pm =t\pm x`$ takes the form
$`_+A_{}_{}A_++[A_+,A_{}]=0.`$ (2.1)
The connections are of the form
$`A_+=BF^+B^1,A_{}=_{}BB^1+F^{},`$ (2.2)
where the mapping $`B`$ is parametrised as
$$B=be^{\nu C}e^{\eta Q_𝐬}=e^{\phi H^0}e^{\stackrel{~}{\nu }C}e^{\eta Q_𝐬},\text{with}b=e^{\phi \stackrel{~}{H}^0},$$
(2.3)
and so $`\stackrel{~}{\nu }=\nu \frac{1}{2}\phi .`$
The special class of models with the $`U(1)`$ Noether current proportional to a topological current, occurs for those models where the grades of the first constant terms of the nonvanishing grade potentials $`F^\pm `$ are equal to $`\pm N_𝐬`$, respectively. So, the potentials $`F^\pm `$ are of the form
$`F^+E_2+F_1^+,F^{}E_2+F_1^{},`$ (2.4)
where
$`F_1^+=2\sqrt{im}(\psi _RE_+^0+\stackrel{~}{\psi }_RE_{}^1),F_1^{}=2\sqrt{im}(\psi _LE_+^1\stackrel{~}{\psi }_LE_{}^0),`$ (2.5)
and $`E_{\pm 2}=mH^{\pm 1}`$ ($`m`$ =constant).
Substituting the gauge potentials (2.2) and (2.4) into (2.1), one gets the equations of motion
$`_+(_{}bb^1)+_+_{}\nu C`$ $`=`$ $`m^2e^{2\eta }[H^1,bH^1b^1]+e^\eta [F_1^{},bF_1^+b^1],`$ (2.6)
$`_{}F_1^+`$ $`=`$ $`me^\eta [H^1,b^1F_1^{}b],`$ (2.7)
$`_+F_1^{}`$ $`=`$ $`me^\eta [H^1,bF_1^+b^1].`$ (2.8)
Making use of the explicit form of the connections $`b`$ and $`F_1^\pm `$, Eqs.(2.3) and (2.5), respectively, into Eqs.(2.6)-(2.8), one can get the equations of motion
$`^2\phi `$ $`=`$ $`4m_\psi \overline{\psi }\gamma _5e^{\eta +2\phi \gamma _5}\psi ,`$ (2.9)
$`^2\stackrel{~}{\nu }`$ $`=`$ $`2m_\psi \overline{\psi }(1\gamma _5)e^{\eta +2\phi \gamma _5}\psi {\displaystyle \frac{1}{2}}m_\psi ^2e^{2\eta },`$ (2.10)
$`^2\eta `$ $`=`$ $`0,`$ (2.11)
$`i\gamma ^\mu _\mu \psi `$ $`=`$ $`m_\psi e^{\eta +2\phi \gamma _5}\psi ,`$ (2.12)
$`i\gamma ^\mu _\mu \stackrel{~}{\psi }`$ $`=`$ $`m_\psi e^{\eta 2\phi \gamma _5}\stackrel{~}{\psi },m_\psi 4m.`$ (2.13)
The equations (2.9)-(2.13) can be derived from the Lagrangian
$`{\displaystyle \frac{1}{k}}_{CATM}={\displaystyle \frac{1}{4}}_\mu \phi ^\mu \phi +{\displaystyle \frac{1}{2}}_\mu \nu ^\mu \eta {\displaystyle \frac{1}{8}}m_\psi ^2e^{2\eta }+i\overline{\psi }\gamma ^\mu _\mu \psi m_\psi \overline{\psi }e^{\eta +2\phi \gamma _5}\psi .`$ (2.14)
It is real (for $`\eta `$ $`=`$real) if $`\stackrel{~}{\psi }`$ is proportional to the complex conjugate of $`\psi `$, and if $`\phi `$ is pure imaginary. This is true for the particular solutions of (2.9)-(2.13) such as the 1-soliton (1-antisoliton), soliton-soliton (antisoliton-antisoliton) . The general reality conditions, imposed on the fields in order to define a real Lagrangian with its corresponding Hamiltonian bounded from below, will be discussed in section 4.
## 3 Affine two-loop WZNW models and reduction
We provide a brief presentation of the Hamiltonian reduction to obtain the conformal affine Toda model coupled to the matter field (CATM) starting from the two-loop WZNW theory. The action for the two-loop WZNW model is
$`S_{WZNW}(\widehat{g})={\displaystyle \frac{\kappa }{8\pi }}{\displaystyle _𝐁}Tr(\widehat{g}^1_\mu \widehat{g}\widehat{g}^1^\mu \widehat{g})+{\displaystyle \frac{\kappa }{12\pi }}{\displaystyle _𝐁}ϵ^{\mu \nu \sigma }Tr(\widehat{g}^1_\mu \widehat{g}\widehat{g}^1_\nu \widehat{g}\widehat{g}^1_\sigma \widehat{g})`$ (3.1)
with the fields being mappings from $`𝐁`$ ($`𝐁=D\text{x}𝐑`$, with $`D`$ being a disc) to the affine KM group $`\widehat{G}`$.
The equations of motion read
$`_+(_{}\widehat{g}\widehat{g}^1)=0,_{}(\widehat{g}^1_+\widehat{g})=0.`$ (3.2)
We now consider those group elements that can be written in the modified Gauss decomposition form
$`\widehat{g}=NBM,`$ (3.3)
where N, B and M are spanned by the subalgebras $`\widehat{𝒢}_+`$, $`\widehat{𝒢}_0`$ and $`\widehat{𝒢}_{}`$, respectively. Introducing the mappings $`K_{L/R}`$ as $`_{}\widehat{g}\widehat{g}^1=NK_LN^1`$ and $`\widehat{g}^1_+\widehat{g}=M^1K_RM`$, and taking into account (3.2) one gets
$`_{}K_R=[K_R,_{}MM^1],_+K_L=[K_L,N^1_{}N].`$ (3.4)
To recover the model (2.6)-(2.8) we impose the constraints
$`(_{}MM^1)_2=B^1E_2B,(_{}MM^1)_{<2}=0,`$ (3.5)
$`(N^1_+N)_2=BE_2B^1,(N^1_+N)_{>2}=0.`$ (3.6)
Then, the mappings $`_{}MM^1`$ and $`N^1_+N`$ have components in the subspaces $`\widehat{𝒢}_1`$ and $`\widehat{𝒢}_1`$, respectively. To give a relation to system (2.6)-(2.8), make the correspondence
$`B_{}MM^1B^1=E_2+F_1^{}F^{},B^1N^1_+NB=E_2+F_1^+F^+.`$ (3.7)
Using the definitions given for $`K_{L/R}`$ and taking into account the constraints (3.5)-(3.6), we can substitute them into (3.4), showing that $`B,F^+`$ and $`F^{}`$ satisfy (2.6)-(2.8). Then, the conformal and integrable model called affine Toda model coupled to the matter field (CATM) can be obtained from the affine (two-loop) WZNW model by means of a Hamiltonian reduction. This relation of the CATM theory to the one of WZNW model allows us to write the following relationship between their coupling constants
$`\kappa =\mathrm{\hspace{0.17em}2}\pi k.`$ (3.8)
Let us recall that in the WZNW model it is a well known fact that the coupling constant $`\kappa `$ takes integer values.
## 4 A real Lagrangian sub-model
From the point of view of their eventual quantization it would be important to distinguish those models whose kinetic terms are positive-definite and whose action is real. So, we are going to consider a sub-model of the conformal affine Toda system coupled to matter (Dirac) field (CATM) (2.14), defining the two-dimensional field theory
$`{\displaystyle \frac{1}{k}}={\displaystyle \frac{1}{4}}_\mu \phi ^\mu \phi +{\displaystyle \frac{1}{2}}_\mu \nu ^\mu \eta +{\displaystyle \frac{1}{8}}m_\psi ^2e^{2\eta }+i\overline{\psi }\gamma ^\mu _\mu \psi m_\psi \overline{\psi }e^{\eta +2i\phi \gamma _5}\psi ,`$ (4.1)
where $`\overline{\psi }\psi ^{}\gamma _0`$, and $`\phi `$, $`\eta `$ and $`\nu `$ are real fields.
The Lagrangian (4.1) differs from the Lagrangian $`_{CATM}`$ presented in (2.14) in three points:
i) the conformal affine Toda system coupled to matter (Dirac) field (CATM), Eq. (2.14), contains two Dirac spinor fields, $`\stackrel{~}{\psi }`$ and $`\psi `$, and a complex $`\phi `$ field, as well as, the real fields $`\nu `$ and $`\eta `$.
ii) we have imposed the reality conditions, $`\stackrel{~}{\psi }=\psi ^{}`$ (the star means complex conjugation) and $`\phi `$ pure imaginary (making the replacement $`\phi i\phi `$ in (2.14), we have a real $`\phi `$ in (4.1)), in order to have a real Lagrangian. Moreover, for later convenience, we have made the change $`\nu \nu `$. The implications of these reality conditions on the solitonic solutions of the model (2.14) were studied in .
iii) an overall minus sign comes out in order to construct a Hamiltonian bounded from below.
The point ii) deserves a far greater attention. The classical theory defined by (2.14) is well defined mathematically, even for a general complex nature of the field solutions (complex Lagrangian). This fact immediately prompts the reaction that the Hamiltonian, hence the energy of configurations of such a system, can not be bounded below, spelling disaster both at the classical level (unstable modes) and at the quantum level (loss of unitarity). This issue certainly has to be handled carefully before this type of model can be considered as providing a sound toy-model for aspects of dualities between particles and solitons, as is one of the motivations for studying it. Here we follow the prescription to restrict the model to a subspace of well-behaved classical solutions. For example the one and two soliton (anti-soliton) solutions satisfy the above reality conditions . Then taking into account the overall minus sign and the above reality conditions allow us to define a sensible physical Lagrangian. These kind of issues in the case of affine Toda field theories are discussed in Refs. , and for non-abelian Toda theories see, for example, Ref. .
Moreover, instead of point ii), the following equally meaninful reality conditions could have been imposed
ii) $`\stackrel{~}{\psi }=\psi ^{}`$, $`\phi i\phi i\pi /2`$ and $`\nu \nu `$, supplied with the change $`x^\mu x^\mu `$. These reality conditions could have also allowed us to end up with Lagrangian (4.1) starting from the CATM Lagrangian (2.14), provided point iii) were also taken into account.
The most general, $`\stackrel{~}{\psi }=e_\psi \psi ^{}`$ and $`\phi `$ pure imaginary, reality conditions and their implications on the solitonic solutions of this model were also presented in . See below (subsection 4.1) a discussion about the implications of these reality conditions on the classical masses of the fundamental particles and the particle/soliton dualities of the theory.
We are going to mention the relevant symmetries of (4.1). The model (4.1) is invariant under the conformal transformations
$$x_+\widehat{x}_+=f(x_+),x_{}\widehat{x}_{}=g(x_{}),$$
(4.2)
with $`f`$ and $`g`$ being analytic functions; and with the fields transforming as
$`\phi (x_+,x_{})`$ $``$ $`\widehat{\phi }(\widehat{x}_+,\widehat{x}_{})=\phi (x_+,x_{}),`$
$`e^{\nu (x_+,x_{})}`$ $``$ $`e^{\widehat{\nu }(\widehat{x}_+,\widehat{x}_{})}=\left(f^{}\right)^\delta \left(g^{}\right)^\delta e^{\nu (x_+,x_{})},`$ (4.3)
$`e^{\eta (x_+,x_{})}`$ $``$ $`e^{\widehat{\eta }(\widehat{x}_+,\widehat{x}_{})}=\left(f^{}\right)^{\frac{1}{2}}\left(g^{}\right)^{\frac{1}{2}}e^{\eta (x_+,x_{})},`$
$`\psi (x_+,x_{})`$ $``$ $`\widehat{\psi }(\widehat{x}_+,\widehat{x}_{})=e^{\frac{1}{2}(1+\gamma _5)\mathrm{log}\left(f^{}\right)^{\frac{1}{2}}+\frac{1}{2}(1\gamma _5)\mathrm{log}\left(g^{}\right)^{\frac{1}{2}}}\psi (x_+,x_{}),`$
where the conformal weight $`\delta `$, associated to $`e^\nu `$, is arbitrary.
On the other hand, (4.1) is also invariant under the commuting $`U(1)_LU(1)_R`$ left and right local gauge transformations
$`\phi \phi +\xi _+\left(x_+\right)+\xi _{}\left(x_{}\right);\nu \nu ;\eta \eta `$ (4.4)
and
$`\psi e^{i\left(1+\gamma _5\right)\xi _+\left(x_+\right)+i\left(1\gamma _5\right)\xi _{}\left(x_{}\right)}\psi .`$ (4.5)
By a special choice of $`\xi _+\left(x_+\right)=\xi _{}\left(x_{}\right)=\frac{1}{2}\theta `$, with $`\theta =\mathrm{const}.`$, one gets a global $`U(1)`$ transformation
$$\phi \phi ,\nu \nu ,\eta \eta ,\psi e^{i\theta }\psi ,$$
(4.6)
and the Noether current, associated to the choice, is given by
$$J^\mu =\overline{\psi }\gamma ^\mu \psi ,_\mu J^\mu =0.$$
(4.7)
Another choice could also be possible by taking $`\xi _+\left(x_+\right)=\xi _{}\left(x_{}\right)=\frac{1}{2}\alpha `$, with $`\alpha =\mathrm{const}.`$. In this way, one has the global chiral symmetry
$$\psi e^{i\gamma _5\alpha }\psi ,\phi \phi \alpha ,\nu \nu ,\eta \eta ,$$
(4.8)
and the corresponding Noether current is
$$J_5^\mu =\overline{\psi }\gamma ^\mu \gamma _5\psi +\frac{1}{2}^\mu \phi ,_\mu J_5^\mu =0.$$
(4.9)
Concerning the topological current, the Lagrangian (4.1) is invariant under $`\phi \phi +n\pi `$, with all the other fields unchanged. Indeed, the vacua are infinitely degenerate, and the topological charge
$$Q_{\mathrm{topol}.}𝑑xj^0,j^\mu =\frac{1}{2\pi }ϵ^{\mu \nu }_\nu \phi ,$$
(4.10)
depending only on the asymptotic values of $`\phi `$, at $`x=\pm \mathrm{}`$, can take non-zero values.
Next, for the sake of completeness, we are going to discuss the relevance of these symmetries in the reduction process to the off-critical affine Toda model coupled to the matter (ATM) . Associated to the conformal symmetry (4.2) there are two chiral currents
$$𝒥=i\psi ^{}\left(1+\gamma _5\right)\psi +i_+\phi +_+\eta ,\overline{𝒥}=i\psi ^{}\left(1\gamma _5\right)\psi +i_{}\phi +_{}\eta ,$$
(4.11)
satisfying
$$_{}𝒥=0;_+\overline{𝒥}=0.$$
(4.12)
Notice, from (4.3), that the currents $`𝒥`$ and $`\overline{𝒥}`$ have conformal weights $`(1,0)`$ and $`(0,1)`$ respectively. Under the conformal transformations (4.2), the chiral currents transform as
$`𝒥(x_+)`$ $``$ $`[\mathrm{ln}f^{}(x_+)]^1\left(𝒥(x_+)[\mathrm{ln}f^{}(x_+)]^{}\right),`$ (4.13)
$`\overline{𝒥}(x_{})`$ $``$ $`[\mathrm{ln}g^{}(x_{})]^1\left(\overline{𝒥}(x_{})[\mathrm{ln}g^{}(x_{})]^{}\right).`$ (4.14)
Then, given a solution of the model, one can always map it, under a conformal transformation, into a solution where
$$𝒥=0,\overline{𝒥}=0.$$
(4.15)
Such a procedure amounts to “gauging” away the free field $`\eta `$. In other words, (4.15) are constraints implementing a Hamiltonian reduction. This resembles the connection between the affine (AT) and conformal affine Toda (CAT) models performed in through a gauge fixing of the conformal symmetry. The quantum version of the reduction CAT $``$ AT was performed in .
Now, let us explain the gauge fixing procedure. Notice from (4.3) that $`\phi `$ is a scalar under conformal transformations and if we set $`\delta `$ to zero, $`e^\nu `$ is also scalar. On the other hand $`e^\eta `$ is a $`(1/2,1/2)`$ primary field. Let us perform a conformal transformation (4.3) with
$`f^{}(x_+)=e^{2\eta _+(x_+)},g^{}(x_{})=e^{2\eta _{}(x_{})},`$ (4.16)
where $`\eta _\pm (x_\pm )`$ are solutions of the $`\eta `$ free field ( $`\eta (x_+,x_{})=\eta _+(x_+)+\eta _{}(x_{})`$; with $`\eta _\pm (x_\pm )`$, being arbitrary functions), one gets
$`\widehat{\phi }(\widehat{x}_+,\widehat{x}_{})\phi (x_+,x_{}),e^{\widehat{\eta }_+(\widehat{x}_+,\widehat{x}_{})}1,e^{\widehat{\nu }_+(\widehat{x}_+,\widehat{x}_{})}e^{\nu _+(x_+,x_{})}.`$ (4.17)
Therefore, we are choosing one solution in each orbit of the conformal group. Then for every regular solution of the $`\eta `$ field the CATM defined on a space-time $`(x_+,x_{})`$ corresponds to an off-critical sub-model which we call the affine Toda model coupled to the matter (ATM), with the extra field $`\nu `$ defined on a space time $`(\widehat{x}_+,\widehat{x}_{})`$. For the particular solution $`\eta =0`$ the CATM and the ATM are defined on the same space-time.
Thus setting the $`\eta `$ field to zero in the equations of motion of the CATM model (2.9)-(2.13) we recover the equations (for the spinor $`\psi `$ and scalar $`\phi `$ fields),
$`^2\phi `$ $`=`$ $`4m_\psi \overline{\psi }\gamma _5e^{2\phi \gamma _5}\psi ,`$ (4.18)
$`i\gamma ^\mu _\mu \psi `$ $`=`$ $`m_\psi e^{2\phi \gamma _5}\psi ,`$ (4.19)
$`i\gamma ^\mu _\mu \stackrel{~}{\psi }`$ $`=`$ $`m_\psi e^{2\phi \gamma _5}\stackrel{~}{\psi },m_\psi 4m,`$ (4.20)
which define the off-critical ATM model
$`{\displaystyle \frac{1}{k}}_{ATM}(\phi ,\psi ,\overline{\psi })={\displaystyle \frac{1}{4}}_\mu \phi ^\mu \phi +i\overline{\psi }\gamma ^\mu _\mu \psi m_\psi \overline{\psi }e^{2i\phi \gamma _5}\psi .`$ (4.21)
This Lagrangian is not conformal invariant and defines the ATM model. Let us now explain the role played by the fields $`\nu `$ and $`\eta `$. In the construction of these fields correspond to the extension of the $`sl(2)`$ loop algebra and are responsible for making the system (4.1) conformally invariant. The field $`\eta `$ is a kind of conformal “gauge connection” and by choosing a particular constant solution ($`\eta =`$ const. is a solution of the equations of motion) the conformal symmetry is broken, thus obtaining a massive theory. Notice that the above gauge fixing ($`\eta 0`$ ) holds, so far, only at the clasical level. In the next section, we will give a quantum version of this reduction. On the other hand, the soliton solutions are obtained in the orbit of the vacuum solution, $`\eta =0`$, and their corresponding masses are determined solely by the behaviour at $`x=\pm \mathrm{}`$ of the space derivative of the auxiliary field $`\nu `$ .
Let us notice that the constraints (4.15), once we have set $`\eta =0`$, are equivalent to
$$\frac{1}{2\pi }ϵ^{\mu \nu }_\nu \phi =\frac{1}{\pi }\overline{\psi }\gamma ^\mu \psi .$$
(4.22)
Therefore, in the gauge fixed model, Eq. (4.21), the Noether current (4.7) is proportional to the topological current (4.10). This equivalence between the currents leads to some remarkable properties of the off-critical ATM model. For instance, it implies that the charge density $`\psi ^{}\psi `$ is proportional to the space derivative of $`\phi `$; consequently, the matter field get confined inside the $`\phi `$ field solitons. It can be seen, independently of the reality conditions imposed above, that the one-soliton and two-soliton solutions of (2.14) satisfy the relationship (4.22) between the Noether and topological currents .
The one-soliton solution for $`\phi `$ is a sine-Gordon type soliton, and the corresponding $`\psi `$ solution is of the massive Thirring model type . In addition, one can check that these solutions satisfy (4.22), and so they become solution of the gauge fixed model, implying that the Dirac field must be confined inside the solitons. In fact, the one-soliton (antisoliton) type solutions of the system of equations (2.9)-(2.13) can be written in the compact form
$`\phi `$ $`=`$ $`2\mathrm{arctan}\left(\mathrm{exp}\left(2m_\psi \text{sign}(e_\psi )\left(xx_0vt\right)/\sqrt{1v^2}\right)\right)`$
$`\psi `$ $`=`$ $`e^{i\theta }\sqrt{m_\psi }e^{m_\psi \text{sign}(e_\psi )\left(xx_0vt\right)/\sqrt{1v^2}}\left(\begin{array}{c}\text{sign}(e_\psi )\left(\frac{1v}{1+v}\right)^{1/4}\frac{1}{1+\text{sign}(e_\psi )ie^{2m_\psi \left(xx_0vt\right)/\sqrt{1v^2}}}\\ \text{sign}(e_\psi )\left(\frac{1+v}{1v}\right)^{1/4}\frac{1}{1\text{sign}(e_\psi )ie^{2m_\psi \left(xx_0vt\right)/\sqrt{1v^2}}}\end{array}\right)`$ (4.25)
$`\nu `$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{log}\left(1+\mathrm{exp}\left(4m_\psi \text{sign}(e_\psi )\left(xx_0vt\right)/\sqrt{1v^2}\right)\right)+{\displaystyle \frac{1}{8}}m_\psi ^2x_+x_{}`$
$`\eta `$ $`=`$ $`0`$ (4.26)
and the solution for $`\stackrel{~}{\psi }`$ is proportional to the complex conjugate of $`\psi `$, i.e., $`\stackrel{~}{\psi }=e_\psi \psi ^{}`$ ($`e_\psi `$ being a real number). According to the reality conditions ii) or ii) imposed above, we have $`e_\psi =\pm 1`$ (for any real number $`e_\psi `$, one can make a convenient rescaling of the $`\psi `$ fields in order to normalize to $`e_\psi =\pm 1`$ ). In such case, the topological charge is<sup>2</sup><sup>2</sup>2We are using the definition presented in for the topological charge (4.10) as being twice that of the reference , in order to make it integer.
$`Q_{\mathrm{topol}.}=\text{sign}e_\psi =\pm 1.`$ (4.27)
Hence, one has both the soliton and the anti-soliton solutions of the real Lagrangian sub-model, Eq. (4.1). Indeed, one can construct solutions of system (4.1) starting from solutions of system (2.14) if the condition ii) or ii) is taken into account. In particular this will be true for the solitonic solutions. Let us mention that in the authors have discarded one of the solutions (soliton or antisoliton) since they have used only one of the possible reality conditions. Since for the study of the full ATM quantum spectrum (subsection 5.2) it could be desirable to know all the soliton type classical solutions (solitons and breathers), let us mention that a carefull analysis reveals the absence of pure imaginary breather (doublet) solutions for the field $`\phi `$ of (2.14) ; i.e., the system (4.1) does not possess bound soliton-antisoliton pairs as a solution of its field equations. However, such type of solutions exist for a general complex but asymptotically pure imaginary $`\phi `$ field of (2.14) .
Observe that the condition (4.22) together with the equations of motion for the spinor fields (2.12)-(2.13) imply the equation of motion for $`\phi `$; namely, Eq. (2.9). Therefore in the gauge fixed model, defined by (4.21), one can replace a second order differential equation, i.e. (4.18), by two first order equations, i.e. (4.22).
The vaccum solution, $`\eta =0`$, was used in to perform the dressing transformation in order to obtain soliton solutions, which are in the orbit of a vacuum solution. Let us emphasize that these transformations do not excite the field $`\eta `$ and the solitonic solutions are solutions of the gauge fixed model; i.e., the off-critical ATM model.
### 4.1 Classical masses of the fundamental particles and solitons
The CATM model is conformally invariant and, therefore, their fundamental particles are massless. Following the arguments presented in , the masses of the model are generated by the spontaneous breakdown of the conformal symmetry of the model by the choice of a particular vaccuum configuration $`\eta =\eta _0=`$constant. Then, considering only the linear field approximations in (2.6)-(2.8) and writting the mapping $`B`$ as $`B=e^T`$, one gets
$`_+_{}T=v_\eta [E_2,[E_2,T]],`$ (4.28)
$`_+_{}F_1^\pm =v_\eta [E_2,[E_2,F_1^\pm ]],`$ (4.29)
where $`v_\eta =e^{2\eta _0}.`$ Thus, the masses of the fundamental particles are given by an eigenvalue equation
$`[E_2,[E_2,X]]=\lambda X,`$ (4.30)
where $`X\widehat{𝒢}_n`$, $`n=0,\pm 1`$. Since $`_+_{}=\frac{1}{4}(_t^2_x^2)`$ we obtain the masses from the Klein-Gordon type equations, then we get
$`m_\phi =m_{\stackrel{~}{\nu }}=m_\eta =0,m_\psi ^{}=m_{\stackrel{~}{\psi }}^{}=4me^{\eta _0}.`$ (4.31)
Now, following the reasoning presented in Refs. , we may say that the masses of the solitons are generated by the spontaneous breakdown of the conformal symmetry. So, the soliton mass is given by
$`{\displaystyle \frac{M_{\text{sol}}}{\sqrt{1u^2}}}=k_x(\nu \nu _0)_{\mathrm{}}^+\mathrm{}.`$ (4.32)
We simply give the mass of the one-soliton (one-antisoliton) computed in (for $`\eta =0`$); viz.,
$`M_{\text{sol}}=8km.`$ (4.33)
The soliton solutions are created by the eigenvalues $`V`$ of $`E_{\pm 2}`$ . Expanding $`V`$ as $`V=_nV^{(n)}`$, where $`[Q_s,V^{(n)}]=nV^{(n)}`$, one observes that $`[E_2,[E_2,V^{(n)}]]V^{(n)}`$. Then, if some $`V^{(n)},n=0,\pm 1`$, does not vanish, it implies that $`V^{(n)}`$ must be one of the eigenvectors $`X`$ in (4.30). In this way, we associate a soliton with a fundamental particle. In addition, the masses of the corresponding soliton and fundamental particle are determined by the same eigenvalue and they are proportional to a $`U(1)`$ charge . It is then possible to have more than one fundamental particle associated to a one-soliton solution if the expansion of $`V`$ contains more than one non-vanishing $`V^{(n)},n=0,\pm 1`$. It is just this argument that allowed us to impose the reality conditions above, identifying $`\stackrel{~}{\psi }`$ as being the complex conjugate of $`\psi `$ (up to a factor $`e_\psi `$), since both elementary particles ( $`\psi `$ and $`\stackrel{~}{\psi }`$) have the same mass $`4me^{\eta _o}`$, Eq. (4.31), and are associated to the same soliton(antisoliton) solution (4.26) for the $`\phi `$ field. Therefore, the imposition of the reality condition $`\stackrel{~}{\psi }=e_\psi \psi ^{}`$ on the CATM spinors does not spoil the particle-soliton duality, i.e., the off-critical and physically well defined ATM model inherits from the CATM the remarkable particle-soliton duality property. The above arguments indicate some sort of particle-soliton duality in the theory similar to the electromagnetic duality of some four dimensional gauge theories .
To see more closely the role played by the sine-Gordon and massive Thirring models in describing some aspects of the soliton (particle) sector of our model, let us write the following suggestive relationship between the one-(anti)soliton solution for the scalar $`\phi `$ and the corresponding classical Dirac field $`\psi `$ solution
$$\psi _R\stackrel{~}{\psi }_L=\frac{m_\psi }{4i}(e^{2\phi }1),\psi _L\stackrel{~}{\psi }_R=\frac{m_\psi }{4i}(e^{2\phi }1).$$
(4.34)
These relations are a good example of the classical correspondence between the sine-Gordon equation and the massive Thirring model . Substituting conveniently the relations (4.34) in the equations of motion (4.18)-(4.19) (written with $`\eta =\eta _0=`$const.) one gets
$$^2\phi =2m_\psi ^2e^{\eta _0}\mathrm{sin}2\phi $$
(4.35)
and
$$i\gamma ^\mu _\mu \psi =2m_\psi \psi 4(\overline{\psi }\gamma _\mu \psi )\gamma ^\mu \psi .$$
(4.36)
The equation (4.35) is the sine-Gordon equation and the one-(anti)soliton solution (4.26) satisfies this equation for $`\eta _0=\mathrm{log}2,`$ and (4.36) is the equation of motion of the massive Thirring model with coupling constant $`g=4`$ and mass M$`{}_{Th}{}^{}=m_\psi ^{}=2m_\psi `$; of course, this $`\psi `$ field mass coincides with the fundamental particle mass (4.31).
## 5 Quantum aspects of the model
Our aim in this section is to consider the quantum aspects of the real Lagrangian sub-model defined by the Lagrangian (4.1). In the next subsection we consider the quantum version of the reduction CATM $``$ ATM through the BRST analysis. In subsection 5.2, we use bosonization techniques to study the model, and perform a further reduction imposing the currents (Noether and topological) equivalence as a constraint. Moreover, we discuss the quantum spectrum of the successively reduced models and outline some remarkable properties concerning the duality soliton/particle, as well as some confinement mechanism present in the ATM.
### 5.1 BRST symmetry and spontaneous breakdown of conformal symmetry
We are interested in studying how the classical reduction of the model, by setting $`\eta =0`$, is recovered at the quantum level. This procedure resembles the conformal affine Toda (CAT)$``$ affine Toda (AT) reduction . We present firstly a naive version of the reduction procedure based on the path integral
$`𝒵`$ $`=`$ $`{\displaystyle 𝒟\phi 𝒟\overline{\psi }𝒟\psi 𝒟\nu 𝒟\eta \mathrm{exp}[iS(\phi ,\psi ,\overline{\psi },\nu ,\eta )]},`$ (5.1)
with $`S`$ the corresponding action of the model defined by the lagrangian (4.1). In the above equation we have not written the gauge fixing term (this term should fix the left-right continuous symmetry (4.4)-(4.5)), the relevant ghost field and its integration measure. In a term of the form $`\lambda _\mu (ϵ^{\mu \nu }_\nu \phi 2\overline{\psi }\gamma ^\mu \psi )`$ ($`\lambda _\mu `$ are Lagrange multipliers), which is the equivalence between the Noether and topological currents, has been considered as a gauge fixing of the symmetry (4.4)-(4.5).
We observe that $`\nu `$ appears as a Lagrange multiplier. Integrating over $`\nu `$ and $`\eta `$ successively, we get
$`𝒵={\displaystyle 𝒟\phi 𝒟\overline{\psi }𝒟\psi \frac{1}{det^2}\mathrm{exp}\left(iS(\phi ,\psi ,\overline{\psi })\right)},`$ (5.2)
where
$`S(\phi ,\psi ,\overline{\psi })={\displaystyle \frac{1}{k}}{\displaystyle d^2x\left\{\frac{1}{4}_\mu \phi ^\mu \phi +i\overline{\psi }\gamma ^\mu _\mu \psi m_\psi \overline{\psi }e^{2i\phi \gamma _5}\psi \frac{1}{8}m_\psi ^2\right\}}.`$ (5.3)
Since the determinant $`\text{det}^2`$ is a constant we have derived an effective theory which is, nothing but, the affine Toda model coupled to the matter field (ATM).
In Ref. , to end up with the quantum off-critical affine Toda model coupled to the matter (ATM), the ideas presented in have been used to perform a reduction process by eliminating the degrees of freedom associated to the fields $`\eta `$ and $`\nu `$. The breaking of the conformal invariance has been made by choosing a specific vacuum in the framework of perturbative Lagrangian approach. The main idea in this process was hamiltonian reduction and spontaneous symmetry breaking of conformal invariance.
A more rigorous analysis of the reduction CATM $``$ ATM can be made by means of BRST analysis. Following similar steps presented in the case of the $`sl(2)`$ affine Toda model , we add to the action (4.1) the following ghost term
$`S_{ghost}=i{\displaystyle d^2x\left\{\frac{1}{2}_\mu \overline{c}^\mu c\right\}}`$ (5.4)
where $`c(\overline{c})`$ is an anticommuting field. One can show that $`S_{tot}=S+S_{ghost}`$ is invariant under the BRST transformation
$`\delta \nu =ic,\delta \overline{c}=\eta ,\delta c=\delta \eta =0,\delta \mathrm{\Phi }(i)=0.`$ (5.5)
where $`\mathrm{\Phi }(i)`$ denotes collectively the fields $`\phi ,\psi `$ and $`\overline{\psi }`$. Then
$`\delta S_{tot}=\mathrm{\hspace{0.17em}0}`$ (5.6)
In addition to the equations of motion for the $`\mathrm{\Phi }(i)`$ fields, we have
$`\mathrm{}c=\mathrm{}\overline{c}=\mathrm{\hspace{0.17em}0}.`$ (5.7)
The conjugate momenta of the (anti)ghost are
$`\pi _c={\displaystyle \frac{i}{2}}_t\overline{c},\pi _{\overline{c}}={\displaystyle \frac{i}{2}}_tc,`$ (5.8)
and the relevant canonical comutation relations are
$`\{\mathrm{\Pi }_\mathrm{\Psi },\mathrm{\Psi }\}_{}=i\delta (xy),`$ (5.9)
where $`\mathrm{\Psi }`$ denotes collectively the set of fields $`\{\mathrm{\Phi }(i),\eta ,\nu ,c,\overline{c}\}`$, and the “$`+()`$” signs are valid for the set of fields $`\{\psi ,\overline{\psi },c,\overline{c}\}`$ and $`\{\mathrm{\Phi }(i),\eta ,\nu \}`$, respectively.
The BRST charge is
$`Q_{BRST}=i{\displaystyle 𝑑x(ic\pi _\nu +\eta \pi _{\overline{c}})}={\displaystyle 𝑑x\left(\eta (_tc)(_t\eta )c\right)},`$ (5.10)
and generates the BRST transformations (5.5)
$`\delta \mathrm{\Psi }=[Q_{BRST},\mathrm{\Psi }]_{},`$ (5.11)
and satisfies the nilpotency $`Q_{BRST}^2=\mathrm{\hspace{0.17em}0}`$. Moreover, considering the hermicity property of the fields $`\mathrm{\Psi }^{}=\mathrm{\Psi }`$, we have $`Q_{BRST}^{}=Q_{BRST}`$.
Introduce a suitable wave packet system of massless particles
$`\mathrm{}f_k(x)=0,`$ (5.12)
$`i{\displaystyle 𝑑xf_k^{}(x)\stackrel{}{}_tf_l(x)}=\delta _{kl},`$ $`{\displaystyle \underset{k}{}}f_k(x)f_k^{}(y)=\delta (xy),`$ (5.13)
where $`f\stackrel{}{}_tg=f(_tg)(_tf)g`$. Therefore, since $`\eta ,c`$ and $`\overline{c}`$ are free fields, we can expand them as follows
$`\eta (x)={\displaystyle \underset{k}{}}\eta _kf_k(x)+\eta _k^{}f_k^{}(x),`$
$`c(x)={\displaystyle \underset{k}{}}c_kf_k(x)+c_k^{}f_k^{}(x),`$ (5.14)
$`\overline{c}(x)={\displaystyle \underset{k}{}}\overline{c}_kf_k(x)+\overline{c}_{k}^{}{}_{}{}^{}f_k^{}(x).`$
Besides, the field $`\nu `$ is not a simple pole field, as can be seen from its equation of motion
$`\mathrm{}\nu ={\displaystyle \frac{1}{2}}m_\psi ^2e^{2\eta }2m_\psi \overline{\psi }e^{\eta +2i\phi \gamma _5}\psi `$ (5.15)
That is to say, the field $`\nu `$ is multipole field and thus can not be expanded in the simple form (5.14). Notice that, in section 4.1, precisely the asymptotic behavior of this field has been used to obtain the classical soliton masses. Nevertheless, we can write the multipole field in the form
$`\nu (x)={\displaystyle \underset{k}{}}\nu _kf_k(x)+\nu _k^{}f_k^{}(x)+\mathrm{}`$ (5.16)
where the ellipsis corresponds to the possible modes of the fields $`\mathrm{\Phi }(i)`$ and $`\eta `$ which from the BRST transformations Eq. (5.5) commute with $`Q_{BRST}`$, and then they will not be important in our considerations.
In terms of the creation an annihilation operators, we have
$`Q_{BRST}=i{\displaystyle \underset{k}{}}(c_k^{}\eta _k\eta _k^{}c_k).`$ (5.17)
From the canonical commutation relations (5.9) we have
$`[\eta _k,\nu _l^{}]=[\nu _k,\eta _l^{}]=\delta _{kl}`$ (5.18)
$`\{c_k,\overline{c}_l^{}\}=\{\overline{c}_k,c_l^{}\}=i\delta _{kl}.`$ (5.19)
Using these relations in (5.11) give us
$`\{Q_{BRST},\overline{c}_k\}=\eta _k,[Q_{BRST},\nu _k]=ic_k`$
$`\{Q_{BRST},c_k\}=\mathrm{\hspace{0.17em}0},[Q_{BRST},\eta _k]=\mathrm{\hspace{0.17em}0}.`$ (5.20)
In addition, $`Q_{BRST}`$ commutes with the modes associated with the fields $`\mathrm{\Phi }(i)`$.
From the above considerations we realize that the fields {$`\eta ,\nu `$} and {$`c,\overline{c}`$} form a pair of BRST doublets of Kugo-Ojima’s and then we may use their quartet mechanism. The free property of the ghost fields $`c`$ and $`\overline{c}`$ implies that the total state vector space $`𝒱`$ can be decomposed persistently into a direct product $`_{\mathrm{\Phi }(i)}_{\eta ,\nu }_{c,\overline{c}}`$, where $`_{\mathrm{\Phi }(i)}`$ correponds to the vector space spanned by the modes associated with the fields $`\mathrm{\Phi }(i)`$, $`_{c,\overline{c}}`$ is the Fock space spanned by $`c`$ and $`\overline{c}`$ alone, and similarly for $`_{\eta ,\nu }`$ . The proof that the ghosts decouple from the physical subspace goes as follows. Consider the following projection operator $`P^{(0)}`$,
$`P^{(0)}:_{\mathrm{\Phi }(i)}_{\eta ,\nu }_{c,\overline{c}}_{\mathrm{\Phi }(i)}|0>_{\eta ,\nu }|0>_{c,\overline{c}}`$ (5.21)
where the vacua $`|0>_{\eta ,\nu }`$ and $`|0>_{c,\overline{c}}`$ are defined by
$`\eta _k|0>_{\eta ,\nu }=\nu _k|0>_{\eta ,\nu }=\mathrm{\hspace{0.17em}0},c_k|0>_{c,\overline{c}}=\overline{c}_k|0>_{c,\overline{c}}=\mathrm{\hspace{0.17em}0}.`$ (5.22)
Then the projection operator commutes with $`Q_{BRST}`$ trivially because of the commutativity of $`Q_{BRST}`$ with the $`\mathrm{\Phi }(i)`$ modes. Here we follow closely the procedure presented in Refs. . We introduce a set of operators $`P^{(n)}(n1)`$ defined inductively as
$`P^{(n)}={\displaystyle \frac{1}{n}}{\displaystyle \underset{k}{}}\left(\nu _k^{}P^{(n1)}\eta _k\eta _k^{}P^{(n1)}\nu _kic_k^{}P^{(n1)}\overline{c}_k+i\overline{c}_k^{}P^{(n1)}c_k\right).`$ (5.23)
These operators $`P^{(n)}(n0)`$ commute with $`Q_{BRST}`$. In addition, it can be seen that they are complete
$`{\displaystyle \underset{n0}{}}P^{(n)}=\mathbf{\hspace{0.17em}1},`$ (5.24)
and for $`n1`$, $`P^{(n)}`$ is BRST exact
$`P^{(n)}`$ $`=`$ $`\{Q_{BRST},R^{(n)}\},`$ (5.25)
$`R^{(n)}`$ $`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{k}{}}\left(\overline{c}_k^{}P^{(n1)}\nu _k+\nu _k^{}P^{(n1)}\overline{c}_k\right).`$ (5.26)
Let $`|\psi >`$ be a physical state in the Hilbert space $``$. It must satisfy the physical condition
$`Q_{BRST}|\psi >=\mathrm{\hspace{0.17em}0}.`$ (5.27)
Therefore any physical state $`|\psi >`$ annihilated by $`Q_{BRST}`$ is written as
$`|\psi >={\displaystyle \underset{n0}{}}P^{(n)}|\psi >=P^{(0)}|\psi >+Q_{BRST}\left({\displaystyle \underset{n1}{}}R^{(n)}|\psi >\right)`$ (5.28)
This means that the physical state is equivalent to its projection onto $`_{\mathrm{\Phi }(i)}|0>_{\eta ,\nu }|0>_{c,\overline{c}}`$ modulo the BRST operator. That is to say, we end up with a theory in which only the modes of the fields $`\mathrm{\Phi }(i)`$ are present (up to the zero modes of the {$`\nu ,\eta `$} and {$`c,\overline{c}`$} fields). Therefore the physical Hilbert space of the CATM (4.1) becomes exactly the one of the off-critical ATM model (4.21). Then we can consider the latter as a Hamiltonian reduced and spontaneously conformal-symmetry-broken version of the CATM.
A perturbative Lagrangian viewpoint has also been considered in order to understand the above reduction process . If one considers scattering processes whose external legs consist only of $`\phi ,\psi `$ and $`\overline{\psi }`$ particles, then one can easily see from the structure of the propagators and the interaction terms, the $`\nu `$ and $`\eta `$ modes decouple and therefore we can ignore the fields $`\nu `$ and $`\eta `$ in the Lagrangian (4.1) by setting these fields to zero in the corresponding action. Then we end up with the action associated to the Lagrangian (4.21) which is just the off-critical model (ATM). However, this analysis relies upon the use of perturbation theory around the trivial vaccum, i.e., $`\phi =\eta =\psi =\overline{\psi }=0,\nu =1/8m_\psi ^2x_+x_{}`$. Then the analysis of the soliton sector is missing. We believe that for this sector the zero modes of the pair of BRST doublets and the nontrivial boundary conditions of the $`\nu `$ and $`\phi `$ fields must be taken into account.
In connection to this reduction process, notice that the argument used in section 4.1 to define the masses of the solitons was to consider the field $`e^{2\eta }`$ as a kind of Higgs field, since it not only spontaneously breaks the conformal symmetry, but also because its vacuum expectation value sets the mass scale of the theory . The role played by the non-trivial boundary conditions of the field $`\nu `$ in connection to the ATM model was pointed out in section 4.1. It was used to provide a relation for the classical soliton masses of the theory.
### 5.2 Bosonization approach and quantum reduction
A startling property which was exploited in the study of two-dimensional field theories is related to the possibility of transforming Fermi fields into Bose fields, and vice versa (for a complete review of the most important references in the field see ). The existence of such a transformation, called bosonization, provided in the last years a powerful tool to obtain nonperturbative information of two-dimensional field theories.
Even though our aim is to solve the model (5.3) exactly, one may use the well known semiclassical methods to verify the validity of the conservation laws, such as (4.7) and (4.9), at the quantum level; as well as to quantize the static classical solutions. In particular, the charge fractionization phenomenon for the fermions of the model (5.3), interacting with the external soliton field $`\phi ,`$ has been considered in , and the conservation laws (4.7) and (4.9) have been verified in the semi-classical approximation.
In the author derived the path-integral version of Coleman’s proof of the equivalence between the massive Thirring model and sine-Gordon models . A Lagrangian of the type (5.3) (plus a free massless scalar field) appears in that process, as a total effective Lagrangian, which gives an equivalent generating functional, after suitable field redefinitions, of the massive Thirring model. Instead of presenting details of that calculation, we will proceed directly bosonizing the Dirac fields of (5.3), which is more suited for our purposes.
Let us write the relevant action of (5.3) in a form which is more convenient for the bosonization of the fermion bilinears
$`S_{ATM}`$ $`=`$ $`{\displaystyle }d^2x\{{\displaystyle \frac{1}{2}}_\mu \phi ^\mu \phi +i\overline{\psi }\gamma ^\mu _\mu \psi m_\psi [\overline{\psi }{\displaystyle \frac{(1+\gamma _5)}{2}}\psi e^{i\sqrt{8/k}\phi }+`$ (5.29)
$`\overline{\psi }{\displaystyle \frac{(1\gamma _5)}{2}}\psi e^{i\sqrt{8/k}\phi }\},`$
where the field rescalings $`\phi \sqrt{2/k}\phi `$ and $`\psi \sqrt{1/k}\psi `$ have been made. The model (5.3), originally proposed by Kogut and Sinclair , has been discussed in Refs. and some of the points that follow have been treated in those papers.
Introduce a new boson field representation of fermion bilinears as
$$i\overline{\psi }\gamma ^\mu _\mu \psi =\frac{1}{2}(_\mu \varphi )^2,$$
(5.30)
$$\overline{\psi }(1\pm \gamma _5)\psi =\mu \mathrm{exp}(\pm i\sqrt{4\pi }\varphi ),$$
(5.31)
$$\overline{\psi }\gamma ^\mu \psi =\frac{1}{\sqrt{\pi }}ϵ^{\mu \nu }_\nu \varphi ,$$
(5.32)
then the action now becomes
$`S_B`$ $`=`$ $`{\displaystyle }d^2x\{{\displaystyle \frac{1}{2}}_\mu \phi ^\mu \phi +{\displaystyle \frac{1}{2}}(_\mu \varphi )^2`$ (5.33)
$`{\displaystyle \frac{1}{2}}\mu m_\psi [\mathrm{exp}i[\sqrt{8/k}\phi +\sqrt{4\pi }\varphi ]+\mathrm{exp}i[\sqrt{8/k}\phi +\sqrt{4\pi }\varphi ]]\}.`$
Introducing new fields
$$\stackrel{~}{c}\frac{\phi /a+\sqrt{4\pi }\varphi }{\sqrt{4\pi +1/a^2}},\stackrel{~}{\sigma }\frac{\varphi /a+\phi \sqrt{4\pi }}{\sqrt{4\pi +1/a^2}},$$
(5.34)
then the action takes the form
$$S_{sG+\stackrel{~}{\sigma }}=d^2x\left[\frac{1}{2}(_\mu \stackrel{~}{c})^2+\frac{1}{2}(_\mu \stackrel{~}{\sigma })^2+\mu m_\psi \mathrm{cos}(\sqrt{4\pi +1/a^2}\stackrel{~}{c})\right],$$
(5.35)
where $`a^2k/8`$.
Therefore we obtain a sine-Gordon model for the field $`\stackrel{~}{c}`$ and a free, massless scalar $`\stackrel{~}{\sigma }`$ field. Let us mention that an analog of the field $`\stackrel{~}{\sigma }`$ will be associated in section 6 to the color charge sector of the full $`QCD_2`$ Lagrangian, and that of the field $`\stackrel{~}{c}`$ to color singlet states (baryons and mesons) in that model.
The sine-Gordon spectrum is known exactly (see, e.g., Refs. ). Thus, since in (5.35) $`\sqrt{4\pi +1/a^2}`$ is grater than $`\sqrt{4\pi }`$, it has only a massive fermion (the soliton) and a massive antifermion (the antisoliton). There are no bound states (see below a discussion regarding this issue).
Moreover, the currents (4.11) (setting $`\eta =0`$, as required in the ATM case) can be expressed as
$`𝒥={\displaystyle \frac{i}{\beta \pi }}_+\stackrel{~}{\sigma },\overline{𝒥}={\displaystyle \frac{i}{\beta \pi }}_{}\stackrel{~}{\sigma }.`$ (5.36)
Let us consider a physical state $`\psi >`$ such that, $`_\pm \stackrel{~}{\sigma }`$ $`\psi >=0,`$ or equivalently,
$`𝒥\psi >=\overline{𝒥}\psi >=0.`$ (5.37)
In terms of the original fields the last relations become
$`(ϵ^{\mu \nu }_\nu \phi 2i\overline{\psi }\gamma ^\mu \psi )\psi >=0.`$ (5.38)
Then, Eq. (5.38) represents the quantum version of the equivalence between the Noether and topological currents. Therefore, their expectation values for any physical states are proportional. This is a kind of “first quantize and then reduce” method (see, e.g., ) and becomes a quantum version of the classical reduction process, ATM $``$ sine-Gordon model (massive Thirring) performed in , imposing the equivalence of those currents as a constraint. Thus one can argue that the relations (5.37) or equivalently (5.38) will hold, in particular, for the energy states in the soliton sector, i.e., the quantum solitons. It corresponds to having none of the $`\stackrel{~}{\sigma }`$ field modes excited . This reduction of the extended “Hilbert” space by the conditions (5.37) forbids states with certain quantum numbers. As we will explain in the next section, a $`\stackrel{~}{\sigma }`$ type field is associated to the “color” degrees of freedom in the bosonized QCD<sub>2</sub> (with one flavour and two colors).
The presence of the physical fermions (zero chirality particles) can be clarified introducing a new fermion field that has the quantum numbers of the physical particles. We simply introduce a new fermion field $`\chi `$ with
$`{\displaystyle \frac{1}{\sqrt{\pi }}}ϵ^{\mu \nu }_\nu \stackrel{~}{c}=\overline{\chi }\gamma ^\mu \chi .`$ (5.39)
According to the standard rules which identify the sine-Gordon theory to the charge-zero sector of the massive Thirring model<sup>3</sup><sup>3</sup>3It has been shown that there is no full equivalence between them. They contain a subset of fields with identical correlation functions, but each model also has fields the other one does not., (5.35) can now be written as
$$S_{Th+\stackrel{~}{\sigma }}=d^2x[i\overline{\chi }\gamma ^\mu _\mu \chi m_F\overline{\chi }\chi \frac{1}{2}g(\overline{\chi }\gamma ^\mu \chi )^2+\frac{1}{2}(_\mu \stackrel{~}{\sigma })^2],$$
(5.40)
where
$$\frac{\beta ^2}{4\pi }=\frac{1}{1+g/\pi },\beta \sqrt{4\pi +1/a^2}$$
(5.41)
In this form, it is obvious that our theory consists of a massive fermion with self-interaction, and a free massless scalar. These are the physical fermion $`\chi `$ and antifermion $`\overline{\chi }`$. Then the original fermions have acquired masses. For the sine-Gordon theory it is thought that the leading order quantum correction of the soliton mass is exact and there are no higher order corrections. It would be interesting to compute the relevant quantum corrections to the soliton mass (4.33), for example using the method of , of course in this case that computation should take into account the fermion fluctuations.
The breather masses in terms of the soliton mass $`M_{\text{sol}}`$ are
$$M_N=2M_{\text{sol}}\mathrm{sin}\left(\frac{\pi }{2}\frac{N}{\frac{8\pi }{\beta ^2}1}\right),N=1,2,\mathrm{}<\frac{8\pi }{\beta ^2}1.$$
(5.42)
These quantum corrections were found by Dashen, Hasslacher and Neveu via a semiclassical quantization of the two soliton solutions, and is thought to be exact. Let us point out that in Refs. the authors have not found classical (real) breather solutions for the field $`\phi `$. This fact reflects at the quantum level and the absence of the corresponding breather spectrum is deduced from (5.42); i.e., since $`\beta ^2`$ is greater than $`4\pi `$ then there are no breathers.
In the sine-Gordon theory it is well known that the fundamental particle is, in the quantum theory, the ground state of the breather spectrum . It desappears from the spectrum when the coupling constant becomes strong enough that the elementary particle mass reaches the kink-antikink threshold. Then one can argue that the ‘elementary’ boson associated to the field $`\phi `$ of our model is missing in the quantum spectrum<sup>4</sup><sup>4</sup>4However, the origin of solitonic excitations due to purely quantum effects, with no soliton at the classical level, has recently been reported .
Another point which deserves a remark is the following. The theory (5.29) presents the chiral symmetry (4.8) and despite this symmetry the physical fermion particles have zero chirality.
In fact, the original chiral current, in terms of the fields $`\stackrel{~}{c}`$ and $`\stackrel{~}{\sigma }`$, takes the form
$$A_\mu =\frac{i}{2}_\mu \stackrel{~}{\sigma }.$$
(5.43)
Thus the chiral current (4.9) involves only $`\stackrel{~}{\sigma }`$ and not $`\stackrel{~}{c}`$. This means that the field $`\stackrel{~}{c}`$ and therefore also the physical fermion and antifermion associated with this field $`\chi `$ and $`\overline{\chi }`$ of (5.40) are neutral under chirality. Then even though the elementary fermion field $`\psi `$ has nonzero chirality, the physical fermion particle has zero chirality.
## 6 Realization of color symmetry in the bosonized QCD<sub>2</sub>
It is reasonable to consider the ATM model a sort of $`1+1`$ dimensional bag model for QCD. This connection can be made evident if we introduce the two-by-two matrix $`U=\text{exp}\{2\phi \gamma _5\}`$; thus, we may rewrite the Lagrangian (5.3) as
$`{\displaystyle \frac{1}{k}}_{ATM}(\phi ,\psi ,\overline{\psi })`$ $`=`$ $`{\displaystyle \frac{1}{16}}\left\{\text{tr}\left[U^1_\mu U{\displaystyle \frac{1+\gamma _5}{2}}U^1^\mu U\right]{\displaystyle \frac{1}{2}}\text{tr}\left[U^1_\mu U\right]\text{tr}\left[U^1^\mu U\right]\right\}+`$ (6.1)
$`i\overline{\psi }\gamma _\mu ^\mu \psi \overline{\psi }U\psi ,`$
which is a two-dimensional analog of the low energy effective Lagrangian for QCD (see, e.g. ).
In the following steps we will undertake the study of the relationship between the affine Toda model coupled to the matter (ATM) and the low-energy affective Lagrangian of QCD in $`1+1`$ dimensions. Although QCD<sub>2</sub> is not exactly soluble still some interesting results may be obtained (see e.g. and references therein). A relation of the dynamics of QCD<sub>2</sub> at large distances and small energies to some massive integrable models would lead to the exact (not semi-classical) solution of the strong coupled QCD<sub>2</sub> as anticipated, for example, in Ref. . Concerning the last point, in QCD<sub>2</sub> with $`N_c`$ colors and $`N_f`$ flavors it is possible to study the low-energy spectrum of hadrons in various approximation schemes. The $`1/N_c`$ expansion reveals a spectrum valid in the weak coupling regime, while the strong coupling offers the posibility of understanding the baryon as a generalized sine-Gordon soliton. In this regard, it is appealing to consider the bosonized form of the ATM model, Eq. (5.35), as a low energy effective Lagrangian for QCD<sub>2</sub>. In Ref. , it was outlined a confining mechanism of some degrees of freedom in the ATM model inside the baryons of its spectrum, and it was suggested an analogous behaviour with what one expects to happen in QCD.
With that motivation, we shall now study the bosonized QCD<sub>2</sub> with two internal symmetries, one flavor ($`N_f=1`$) and two colors ($`N_c=2`$). Our treatment closely follows that of Ref. and, therefore, we only sketch a few essential steps. First of all, let us consider a single massless Dirac theory with two internal symmetries ($`N_f`$=1 and $`N_c`$=2) defined by
$`_{}={\displaystyle \frac{1}{2}}\overline{\psi }_\alpha i\gamma ^\mu _\mu \psi _\alpha ,\alpha =1,2.`$ (6.2)
This Lagrangian can be written in terms of the bosonic variables, $`\phi _\alpha ,\alpha =1,2`$, one for each internal fermionic degree of freedom. The bosonised lagrangian is
$`_{}={\displaystyle \frac{1}{2}}_\mu \phi _\alpha ^\mu \phi _\alpha ,\alpha =1,2.`$ (6.3)
with the fermionic bilinear verifying
$`:\overline{\psi }_\alpha \gamma ^\mu \psi _\alpha :={\displaystyle \frac{1}{\sqrt{\pi }}}ϵ^{\mu \nu }_\nu \phi _\alpha .`$ (6.4)
Considering the linear combinations $`\stackrel{~}{\sigma }=(1/2)(\phi _1\phi _2)`$ and $`\stackrel{~}{c}=(1/2)(\phi _1+\phi _2)`$ allow us to write the $`U(1)`$ flavor and the third component of color isospin as topological currents,
$`J_{U(1)}^\mu ={\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha }{}}:\overline{\psi }_\alpha \gamma ^\mu \psi _\alpha :={\displaystyle \frac{1}{\sqrt{2\pi }}}ϵ^{\mu \nu }_\nu \stackrel{~}{c},`$ (6.5)
$`J_3^\mu ={\displaystyle \underset{\alpha }{}}\text{Tr}:\overline{\psi }_\alpha \gamma ^\mu {\displaystyle \frac{\tau _3}{2}}\psi _\alpha :={\displaystyle \frac{1}{\sqrt{2\pi }}}ϵ^{\mu \nu }_\nu \stackrel{~}{\sigma },`$ (6.6)
and the corresponding conserved charges as
$`B={\displaystyle \frac{1}{\sqrt{2\pi }}}[\stackrel{~}{c}(+\mathrm{})\stackrel{~}{c}(\mathrm{})],`$ (6.7)
$`T_3={\displaystyle \frac{1}{\sqrt{2\pi }}}[\stackrel{~}{\sigma }(+\mathrm{})\stackrel{~}{\sigma }(\mathrm{})].`$ (6.8)
The bosonized Lagrangian presents an explicit decoupling of flavor and color degrees of freedom when written in terms of the new fields
$`_{}={\displaystyle \frac{1}{2}}_\mu \stackrel{~}{\sigma }^\mu \stackrel{~}{\sigma }+{\displaystyle \frac{1}{2}}_\mu \stackrel{~}{c}^\mu \stackrel{~}{c}.`$ (6.9)
This Lagrangian shows explicitly the $`U(1)_FU(1)_C`$ symmetry. In order to realize the global $`SU(2)`$ symmetry, the $`\stackrel{~}{\sigma }`$ must take non-trivial boundary conditions (NTBC), these NTBC determine the vacuum structure of the theory and the quantum numbers associated to the physical states determined from (6.8) .
One can consider a mass term in (6.2), which after bosonization takes the form
$`S_M`$ $`=`$ $`{\displaystyle d^2x\widehat{k}^2\mu M(\text{cos}\sqrt{2\pi }\stackrel{~}{c})(\text{cos}\sqrt{2\pi }\stackrel{~}{\sigma })},`$ (6.10)
$`\mu `$ is a renormalization mass and $`\widehat{k}`$ is a constant appearing in the bosonization formula . Since the classical potential has a minimum at $`V_M(\stackrel{~}{c}_o,\stackrel{~}{\sigma }_o)=\widehat{k}^2\mu M`$, one has a degenerate vacuum formed by the $`(\stackrel{~}{c}_o,\stackrel{~}{\sigma }_o)`$ lattice
$`\mathrm{}`$ $``$ $`\{(n,m);n,m\mathrm{Z}\}{\displaystyle \{(n+\frac{1}{2},m+\frac{1}{2});n,m\mathrm{Z}\}}.`$ (6.11)
The minima are related to the possible boundary conditions, and therefore give us precise information about the solutions of the bosonized theory, viz.,
$`\underset{x\pm \mathrm{}}{lim}(\stackrel{~}{c}(x),\stackrel{~}{\sigma }(x))`$ $`=`$ $`(\stackrel{~}{c}_o,\stackrel{~}{\sigma }_o)\mathrm{}.`$ (6.12)
Moreover, the full QCD<sub>2</sub> bosonized Lagrangian acquires a term due to the quark-gluon interaction
$`V(\stackrel{~}{\sigma })`$ $`=`$ $`{\displaystyle \frac{e^2}{32\pi }}\stackrel{~}{\sigma }^2+\sqrt{\pi }\widehat{k}^4\mu ^2[1{\displaystyle \frac{\text{sin}\sqrt{2\pi }\stackrel{~}{\sigma }}{\sqrt{2\pi }\stackrel{~}{\sigma }}}],`$ (6.13)
where $`e`$ is the quark-gluon interaction coupling constant. This potential is positive definite and has an unique absolute minimum $`V(\stackrel{~}{\sigma })=0`$, at $`\stackrel{~}{\sigma }=0`$. Therefore the interaction with the gauge field reduces the possible quantum numbers of the physical states to
$`\mathrm{}_\text{b}`$ $``$ $`\{(n,0);n\mathrm{Z}\}.`$ (6.14)
In Ref. , it was verified, by computing $`<T^2>`$, that these states are color singlets. It is possible to give a topological interpretation to this mechanism. Once the confinement interaction is added, the vacuum structure $`\mathrm{}`$ collapses into a colorless subset of the lattice points, $`\mathrm{}_\text{b}`$; i.e., the topology associated to color becomes trivial.
Since there is a relation between the minima structure and NTBC’s, Eq. (6.12), let us study the constraints induced by $`V(\stackrel{~}{\sigma })`$ into the full QCD<sub>2</sub> action
$`S`$ $`=`$ $`{\displaystyle \left[\frac{1}{2}_\mu \stackrel{~}{\sigma }^\mu \stackrel{~}{\sigma }+\frac{1}{2}_\mu \stackrel{~}{c}^\mu \stackrel{~}{c}+\widehat{k}^2\mu M(\text{cos}\sqrt{2\pi }\stackrel{~}{c})(\text{cos}\sqrt{2\pi }\stackrel{~}{\sigma })V(\stackrel{~}{\sigma })\right]}.`$ (6.15)
Notice that the full potential has a lower bound
$`V(\stackrel{~}{c},\stackrel{~}{\sigma })`$ $`=`$ $`V_M(\stackrel{~}{c},\stackrel{~}{\sigma })+V(\stackrel{~}{\sigma })\widehat{k}^2\mu M,(\stackrel{~}{c},\stackrel{~}{\sigma }).`$ (6.16)
The equality is saturated only for the set of points defined by the lattice, Eq. (6.14), which correspond to the color singlet states. The presence of the color interaction has transformed the free vacuum, Eq. (6.11), in an essential way. The topologically non-trivial solutions associated to the color field $`\stackrel{~}{\sigma }`$ are pushed to infinite energy (infinite mass) and thus the relevant Fock space has been reduced to that of only $`T_3=0`$ states (only the solutions satisfying $`\stackrel{~}{\sigma }(\pm \mathrm{})=0`$ can remain with finite energy , and color solitons associated to non-trivial topology are given infinite energy, see Refs. for details). The singlet spectrum is constituted by particles of baryon number $`B=n`$, i.e., by baryons $`(n=1,2,\mathrm{})`$ or antibaryons $`(n=1,2,\mathrm{})`$ and by mesons $`(n=0)`$.
In order to relate these results to the bosonized ATM model, Eq. (5.35), we must consider the dynamics of QCD<sub>2</sub> at large distances and small energies. The procedure followed is very simple: we eliminate the color field $`\stackrel{~}{\sigma }`$ from the Lagrangian via its classical equation of motion for solutions, in an infinite mass scale limit , and for solutions of vanishing “color” charge, as discussed above; which in physical terms means that the mass scale ($`M^{}`$) must be much bigger than the color singlet state baryon (soliton) mass $`M_{\text{sol}}`$ <sup>5</sup><sup>5</sup>5In solving for the $`\stackrel{~}{\sigma }`$ field from the equations of motion, we are considering the cos$`(\sqrt{2\pi }\stackrel{~}{\sigma })`$ expansion up to the second order term, then our approximation is meaningful in the weak field approximation. This amounts to cosidering soliton like solutions, where the field $`\stackrel{~}{\sigma }`$ asymptotically vanishes as $`x\pm \mathrm{}`$ and makes a lump within a range $`|x|<ϵ`$, for a given $`ϵ`$. . The heavy color field $`\stackrel{~}{\sigma }`$ decouples; as argued above, it freezes down at the value $`\stackrel{~}{\sigma }=\stackrel{~}{\sigma }_o=`$constant (the leading order term in the $`1/M^{}`$ expansion), leaving the physical singlet states described by the “hadron” field $`\stackrel{~}{c}`$. Thus the system is described by the effective Lagrangian involving only the field $`\stackrel{~}{c}`$
$`_{\text{eff}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}_\mu \stackrel{~}{c}^\mu \stackrel{~}{c}+\widehat{k}^2\mu M(\text{cos}\sqrt{2\pi }\stackrel{~}{\sigma }_o)(\text{cos}\sqrt{2\pi }\stackrel{~}{c})`$ (6.17)
The spectrum of this effective Lagrangian involves baryons (a soliton and an antisoliton) and some number of soliton-antisoliton bound states (for the case we are interested in there are just two bound states or mesons). The full QCD<sub>2</sub> Lagrangian (6.15) has a “hadron” spectrum which according to recent studies would not be stable in the mesonic sector (see, e.g., Refs. ), revealing a non-integrable character of QCD<sub>2</sub>. We believe that, in spite of this fact (the breakdown of integrability in the mesonic sector), it is reasonable to describe the baryonic sector of (6.15) by a QCD-motivated integrable field theory, as is evident in its low energy effective Lagrangian, Eq. (6.17), whose color singlet baryonic sector and some of its features may be described by the ATM model, Eq. (5.35), supplied with a condition on its state space; namely, the equivalence between the topological and Noether currents, Eq. (5.37), imposed as a constraint on the physical states. The above idea is further supported considering that in QCD the integrable sector and the mesonic sector are seemingly decoupled . The reduction process ATM $``$ sine-Gordon model, which is performed imposing the equivalence between the Noether and topological currents as a constraint, is the mechanism required to describe the confining phase in the ATM model, as we shall explain in the next section.
## 7 Topological Confinement in the ATM model
We want to examine the question of confinement of the “color” degrees of freedom associated to the field $`\stackrel{~}{\sigma }`$ in the ATM model. We will resort to a “semi-classical” analysis . This goes as follows. We put a pair of external probe color charges $`e^{}`$ and $`\overline{e}^{}=e^{}`$ at $`L/2`$ and $`L/2`$. In order to evaluate the inter-charge potential, we compare the ground state expectation value of the Hamiltonian associated to (5.35) in the presence of an external $`\overline{e}^{}e^{}`$ source, with the corresponding obtained in the absence of such a source
$`V(L)=<\mathrm{\Omega }_Q|H(L)|\mathrm{\Omega }_Q><\mathrm{\Omega }_0|H(0)|\mathrm{\Omega }_0>,`$ (7.1)
where $`H(0)(H(L))`$ and $`|\mathrm{\Omega }_0>(|\mathrm{\Omega }_Q>)`$ are the Hamiltonian and ground state in the absence (presence) of the probe charges, respectively. Let us consider an additional term in (5.29) such that an external spinor field couples to $`\phi `$ in the same form, viz.,
$`m_\psi ^{}[\overline{\psi }_{\text{ext}}{\displaystyle \frac{(1+\gamma _5)}{2}}\psi _{\text{ext}}e^{i\sqrt{8/k}\phi }+\overline{\psi }_{\text{ext}}{\displaystyle \frac{(1\gamma _5)}{2}}\psi _{\text{ext}}e^{i\sqrt{8/k}\phi }],`$ (7.2)
the coupling constant, taken as $`m_\psi ^{}`$, will be associated to the mass of the external fermion. Next, one can bosonize the external fermion bilinears in the usual form,
$`\left(\overline{\psi }\gamma ^\mu \psi \right)_{\text{ext}}={\displaystyle \frac{1}{\sqrt{\pi }}}ϵ^{\mu \nu }_\nu Q`$ (7.3)
where
$`Q={\displaystyle \frac{ke^{}}{2}}\sqrt{\pi }\mathrm{\Xi },`$ (7.4)
where $`e^{}`$ is the external color charge<sup>6</sup><sup>6</sup>6Notice that the factor $`k/2`$ is introduced in relation (7.4) in order to match the definition of the charges presented in section 4 in terms of the unrescaled ATM fields of the Lagrangian (4.21) or (5.3).. The field $`\mathrm{\Xi }`$ above is taken as
$`\mathrm{\Xi }=[\mathrm{\Theta }(xL/2)\mathrm{\Theta }(x+L/2)],`$ (7.5)
this function will represent the external color charge potential of the probe charges $`e^{}\overline{e}^{}`$. This will give a vanishing total “color” charge.
Introducing the field $`\stackrel{~}{c}_Q`$ as
$`\stackrel{~}{c}_Q={\displaystyle \frac{\sqrt{4\pi }Q+\phi /a}{\sqrt{4\pi +1/a^2}}}`$ (7.6)
we can write the total bosonized action as
$`S_{\text{tot}}={\displaystyle d^2x[\frac{1}{2}(_\mu \stackrel{~}{c})^2+\frac{1}{2}(_\mu \stackrel{~}{\sigma })^2+\lambda \mathrm{cos}(\beta \stackrel{~}{c})+\lambda ^{}\mathrm{cos}(\beta \stackrel{~}{c}_Q)]},`$ (7.7)
where $`\lambda \mu m_\psi `$ and $`\lambda ^{}\mu m_\psi ^{}`$.
The bosonized Lagrangian (7.7) will be the starting point for our future analysis and from this point on our treatment will be purely classical. It should, however, be noted that this is better than treating the original fermionic model, since bosonization capture some of the quantum nature of the model. Therefore, one may argue that semiclassically, the constraints (5.37) are equivalent to $`\stackrel{~}{\sigma }=\stackrel{~}{\sigma }_o=`$constant, and consequently the degrees of freedom eliminated by them correspond to the $`\stackrel{~}{\sigma }`$ field.
Since we are looking for a static potential between heavy “quarks” in a color singlet state, the field $`\stackrel{~}{\sigma }`$ may be set equal to a constant (zero) in accordance with the reduction process considered in section 5, Eq. (5.37) and the discussion above. This procedure is the analog of the QCD confinement mechanism, in which the topology associated to an analogous field is dynamically collapsed, as discussed in the last section. Then the static Hamiltonian, in terms of the field $`\stackrel{~}{c}`$ and incorporating the external color charge potential, becomes
$`H(L)`$ $`=`$ $`{\displaystyle 𝑑x[\frac{1}{2}(\stackrel{~}{c}^{})^2\lambda \mathrm{cos}(\beta \stackrel{~}{c})\lambda ^{}\mathrm{cos}(\sqrt{4\pi }Q+\frac{\stackrel{~}{c}}{a^2\beta })]}.`$
The classical potential has a minimum at
$`V_{min}=\lambda \lambda ^{},`$ (7.8)
and this implies a degenerate minima formed by the set of points
$`{\displaystyle \frac{2\pi m}{\beta }}\text{and}\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}2}\pi na^2\beta ;m,n\mathrm{Z}.`$ (7.9)
Then the parameters must satisfy the relationship, $`a^2\beta ^2n=m`$, which implies
$`({\displaystyle \frac{\pi k}{2}}+1)n=m.`$ (7.10)
Therefore $`\pi k`$ should be rational. The coupling of the external field $`Q`$, related to the probe charges, has not affected in an essencial way the nature of the coupling constant $`k`$ (a quantization of $`k`$ is of course expected, since the CATM is related to the one of WZNW ; see section $`3`$ and relation (3.8)). The minima provide us important information about the solutions of the bosonized theory, thus allowed us, for example, to determine the correct quantum numbers ($`B`$ or $`T_3`$) in QCD<sub>2</sub>. In the same way one can consider the boundary conditions in the ATM case
$`\underset{x\pm \mathrm{}}{lim}\stackrel{~}{c}(x)=\stackrel{~}{c}(\pm \mathrm{}),\underset{x\pm \mathrm{}}{lim}\stackrel{~}{\sigma }(x)=\stackrel{~}{\sigma }(\pm \mathrm{})(=0,\text{for “color” singlet states})`$ (7.11)
The inter-charge potential energy $`V(L)`$ is associated to the field configurations that minimize the Hamiltonian. Then the equation of motion of the field $`\stackrel{~}{c}`$ which follows from the bosonized Lagrangian (7.7) is, in the static case
$`\stackrel{~}{c}^{\prime \prime }{\displaystyle \frac{\lambda }{\beta }}\mathrm{sin}(\beta \stackrel{~}{c}){\displaystyle \frac{\lambda ^{}}{a^2\beta }}\mathrm{sin}(\sqrt{4\pi }Q+{\displaystyle \frac{\stackrel{~}{c}}{a^2\beta }})=0.`$ (7.12)
It is well known that the solution of the above classical equation of motion in the bosonic version of the theory contains quantum mechanical information from the fermionic theory. Then we need to solve (7.12) imposing some boundary conditions. The nonlinear equation (7.12) is not exactly solvable and can only be solved considering the first order term in the sine expansions<sup>7</sup><sup>7</sup>7This approximation is only meaningful if the conditions $`\beta |\stackrel{~}{c}(\pm \mathrm{})|<<1`$ and $`4\pi e^{}<<k`$ are satisfied.. In this approximation the equation reduces to
$`\stackrel{~}{c}^{\prime \prime }w^2\stackrel{~}{c}=f\mathrm{\Xi },`$ (7.13)
where
$`w^2=\lambda +{\displaystyle \frac{\lambda ^{}}{a^4\beta ^2}},f={\displaystyle \frac{4\pi \lambda ^{}e^{}}{a^2\beta k}}.`$ (7.14)
This equation is exactly solvable and its solution is
$`\stackrel{~}{c}=\{\begin{array}{cc}Fe^{\alpha x}+\stackrel{~}{c}(+\mathrm{})\hfill & x>L/2,\hfill \\ Ce^{\alpha x}+De^{\alpha x}+E\hfill & |x|<L/2,\hfill \\ Ae^{\alpha x}+\stackrel{~}{c}(\mathrm{})\hfill & x<L/2.\hfill \end{array}`$ (7.18)
We will consider the non-trivial boundary conditions (NTBC): $`\stackrel{~}{c}(+\mathrm{})=\stackrel{~}{c}_+\mathrm{\hspace{0.17em}0}`$ and $`\stackrel{~}{c}(\mathrm{})=0`$. Moreover, the solutions and their derivatives should be matched at the boundaries $`\pm L/2`$ giving us a set of algebraic equations for the unknown coefficients. For example the parameter $`\alpha `$ is given by
$`\alpha ^2=\lambda +{\displaystyle \frac{\lambda ^{}}{a^4\beta ^2}}.`$ (7.19)
The potential derivative becomes
$`{\displaystyle \frac{d}{dL}}V(L)`$ $`=`$ $`{\displaystyle \frac{d}{dL}}[H(L)H(0)]`$
$`=`$ $`{\displaystyle \frac{2\lambda e^{}}{k}}\pi \left[\mathrm{sin}(\sqrt{4\pi }Q(L/2)+{\displaystyle \frac{\stackrel{~}{c}(L/2)}{a^2\beta }})\mathrm{sin}(\sqrt{4\pi }Q(L/2)+{\displaystyle \frac{\stackrel{~}{c}(L/2)}{a^2\beta }})\right].`$
Expanding the sine-terms and keeping the first order terms we have
$`{\displaystyle \frac{d}{dL}}V(L)={\displaystyle \frac{2\lambda ^{}e^{}}{ka^2\beta }}\pi [\stackrel{~}{c}(L/2)\stackrel{~}{c}(L/2)].`$ (7.20)
Finally, for the above field configurations which minimize the Hamiltonian we compute the inter-charge potential
$`V(L)={\displaystyle \frac{\stackrel{~}{c}_+\lambda ^{}e^{}\pi }{a^2k\beta }}\left[L+{\displaystyle \frac{1}{\alpha }}e^{\alpha L}\right].`$ (7.21)
Thus, we observe that the inter-charge potential has both the confining and screening type terms. However, if the non-trivial boundary condition (NTBC) for the sine-Gordon soliton (baryon), $`\stackrel{~}{c}_+0`$, is considered and the inter-charge separation increases, i.e., $`L+\mathrm{}`$, then the confining term dominates. Therefore, taking into account the particle spectrum of our model (see section 5), we can argue that the ATM model presents a permanent confining phase of the “color” degrees of freedom inside the baryons of the model (solitons and antisolitons).
## 8 Discussions
To summarise the results of this paper: it has been shown that the baryonic sector of the low energy effective Lagrangian of QCD<sub>2</sub> with one flavor and two colors, and some of its properties may be described by the integrable and off-critical affine Toda model coupled to the matter field (ATM), provided the equivalence between the Noether and topological currents of the ATM model were imposed as a constraint.
It could be worthwhile to mention that a potential of type (7.21) is present in the Toda lattice system . In the Toda lattice model a particle interacts only with its nearest neighborhoods via a potential defined by (7.21), where $`L`$ represents the distance between the interacting particles. The appearance of this type of potential and the connection of the ATM with the Toda lattice systems, if there exists any in the context studied above, deserves a further investigation.
We believe this model could be useful as a toy-model to implement the so called “Cheshire Cat Mechanism” as a gauge symmetry, used to study confinement in $`1+1`$ dimensional chiral bag models . Work in this direction is in progress .
It is natural to inquire into the extension of the above analysis to models constructed using affine Lie algebras associated to higher rank finite Lie algebras, since the same property, i.e., the equivalence between the Noether and topological currents appears in a class of special models defined in . It could be interesting to check the validity of this property for the various soliton type solutions, and whether the currents equivalence holds true at the quantum level. As the general CATM, for any affine Lie algebra, is related to the two-loop WZNW model, the analysis could also be performed using the WZNW fields, in terms of which the matter fields can be written locally but in a somewhat cumbersome way.
Note Added: After the completion of this work we became aware of the Refs. where the authors were able to recover similar models to the CATM considered above. In the case of the $`sl(2)`$ CATM studied by us, it is the special $`sl(2)`$ case of a family of their so-called bosonic superconformal affine Toda models based on arbitrary affine Lie algebra. The authors considered the construction and the integrability properties, as well as the conformal correspondences of such models. It is also worth to point out that one of the main points of Ref. was the uncovering of the special class of models in which the equivalence between the Noether and topological currents holds true, at least, at the classical level.
Acknowledgements
The author is grateful to Professors L.A. Ferreira, G.M. Sotkov, and A.H. Zimerman for valuable discussions. I thank Professor L. Zhao for correspondence and letting me know the Refs. . R. Bentín and C. Tello are also akcnowledged for fruitful conversations. I would like to thank FAPESP for financial support.
## Appendix A Notations and Conventions
We use the following conventions in two dimensions. The metric tensor is $`g_{\mu \nu }=\text{diag}(1,1)`$ and the antisymmetric tensor $`ϵ_{\mu \nu }`$ is defined so that $`ϵ_{01}=ϵ^{01}=1`$. $`_\pm `$ are derivatives w.r.t. to the light cone variables $`x_\pm =t\pm x`$. The gamma matrices are in the following representation:
$$\gamma _0=i\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\gamma _1=i\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\gamma _5=\gamma _0\gamma _1=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),$$
(A.1)
satisfying anticommutation relations
$`\{\gamma _\mu ,\gamma _\nu \}=2g_{\mu \nu }\mathrm{𝟏},`$ (A.2)
so the spinors $`\psi `$ and $`\overline{\psi }`$ are of the form
$`\psi =\left(\begin{array}{c}\psi _R\\ \psi _L\end{array}\right),\stackrel{~}{\psi }=\left(\begin{array}{c}\stackrel{~}{\psi }_R\\ \stackrel{~}{\psi }_L\end{array}\right),\overline{\psi }=\left(\begin{array}{cc}\stackrel{~}{\psi }_R\stackrel{~}{\psi }_L& \end{array}\right)\gamma _0.`$ (A.8)
## Appendix B The $`sl(2)^{(1)}`$ affine Lie algebra notations
Let us work with the Chevalley basis generators $`H^n,E_\pm ^n`$, $`D`$ and $`C`$ of $`sl(2)^{(1)}`$. The commutation relations are
$`[H^m,H^n]`$ $`=`$ $`2mC\delta _{m+n,0},`$
$`[H^m,E_\pm ^n]`$ $`=`$ $`\pm 2E_\pm ^{m+n},`$
$`[E_+^m,E_{}^n]`$ $`=`$ $`H^{m+n}+mC\delta _{m+n,0},`$
$`[D,T^n]`$ $`=`$ $`nT^n,T^nH^m,E_\pm ^n;`$ (B.1)
all other commutation relations are trivial. The grading operator for the principal gradation defined by the vector $`𝐬=(1,1)`$ is
$`Q_𝐬H^0+N_𝐬D{\displaystyle \frac{1}{2N_𝐬}}Tr(H_𝐬^2)C(N_𝐬𝐬_1+𝐬_2=2).`$ (B.2)
Then, the corresponding eigensubspaces are
$`\widehat{𝒢}_0=\{H^0,C,Q_𝐬\};`$ (B.3)
$`\widehat{𝒢}_{2n+1}=\{E_+^n,E_{}^{n+1}\}n\mathrm{Z};`$ (B.4)
$`\widehat{𝒢}_{2n}=\{H^n\},n\{\mathrm{Z}0\}.`$ (B.5)
Under the above gradation the affine Lie albebra decomposes as
$`\widehat{𝒢}={\displaystyle \underset{s}{}}\widehat{𝒢}_s`$ (B.6)
with
$`[\widehat{𝒢}_s,\widehat{𝒢}_s]\widehat{𝒢}_{s+r}.`$ (B.7)
The grades $`s`$ take zero, positive and negative values, i.e.,
$`\widehat{𝒢}=\widehat{𝒢}_+{\displaystyle \widehat{𝒢}_0\widehat{𝒢}_{}}`$ (B.8)
with
$`\widehat{𝒢}_+{\displaystyle \underset{s>0}{}}\widehat{𝒢}_+,\widehat{𝒢}_{}{\displaystyle \underset{s<0}{}}\widehat{𝒢}_s`$ (B.9)
In addition, we will use a special basis for the generators of $`\widehat{𝒢}_0`$ such that they are all orthogonal to $`Q_𝐬`$ and $`C.`$ Thus, shifting the Cartan elements as $`\stackrel{~}{H}^0`$ = $`H^0\frac{1}{2}C,`$ one gets
$`Tr(C^2)`$ $`=`$ $`Tr(C\stackrel{~}{H}^0)=Tr(Q_𝐬^2)=Tr(Q_𝐬\stackrel{~}{H}^0)=0,`$
$`Tr(Q_𝐬C)`$ $`=`$ $`2,Tr(\stackrel{~}{H}^0\stackrel{~}{H}^0)=Tr(H^0H^0)=2.`$ (B.10)
Here we have used $`Tr(CD)=1,`$ and the normalization for the root, $`\alpha ^2=2`$. For more details of such a special basis, see appendix $`C`$ of Ref. .
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# Escape-Rate Crossover between Quantum and Classical Regimes in Molecular Magnets: A Diagonalization Approach
## Abstract
We have studied numerically the quantum-classical crossover in the escape-rate for an uniaxial spin system with an arbitrarily directed field. Using the simple quantum transition-state theory, we have obtained the boundary separating the first- and the second-order crossover and the escape-rate in the presence of the transverse and longitudinal field. The results apply to the molecular nanomagnet, $`\mathrm{Mn}_{12}`$.
We consider an easy-axis ferromagnetic nanoparticle, or a molecular cluster, that has metastable or degenerate classical spin states. The direction of the magnetization may change due to two mechanisms. At sufficiently high temperature the rate of the magnetization reversal $`\mathrm{\Gamma }`$ obeys the Arrhenius law, $`\mathrm{\Gamma }\mathrm{exp}(\mathrm{\Delta }U/k_BT)`$, with $`\mathrm{\Delta }U`$ being the height of the energy barrier. At a temperature low enough to ignore the thermal activation, quantum tunneling comes into play with $`\mathrm{\Gamma }\mathrm{exp}(\mathrm{\Delta }U/\mathrm{}\omega )`$ where $`\omega `$ is some temperature-independent frequency related to the shape of the metastable potential well. The crossover between thermal and quantum regimes has been intensively studied in nanospin systems.
This issue was first raised by Chudnovsky and Garanin, who observed that the crossover in the spin Hamiltonian $`=DS_z^2H_xS_x`$ becomes sharp (first order) for $`h_x(H_x/(2DS))<0.25`$ and smooth (second order) for $`0.25h_x<1`$. In the exponential approximation, when only the transition exponent is concerned, the first- (second-) order crossover of the escpe-rate is characterized by the discontinuity (continuity) of $`d\mathrm{\Gamma }(T)/dT`$ at the crossover temperature, $`T_0`$. Subsequent calculations rendered the boundary between the first- and the second-order crossover for the uniaxial model with a transverse and longitudinal field. It was also pointed out that the nonanalyticity of the rate for the first-order crossover disappears when quantum corrections to the exponential approximation are computed. The purpose of this Letter is to investigate the “sharpness” of the first-order crossover. This question is especially important in the light of a recent experimental evidence of the first-order crossover in $`\mathrm{Mn}_{12}`$ .
We shall focus on the crossover in the molecular magnet $`\mathrm{Mn}_{12}`$ ($`S=10`$) when the external magnetic field has both transverse and longitudinal components. In order to calculate the splitting of the excited states, we shall perform a numerical diagonaliztion of the Hamiltonian. Using numerical results, we then obtain the group of levels which make the dominant contribution to the thermally assisted tunneling. The full problem of the escape rate will be solved by mapping the spin problem onto a particle one . Summing contributions of all excited levels with account of quantum corrections, we will show that the escape rate becomes analytic for the first-order crossover, but changes sharply around $`T_0`$, in contrast with the second-order crossover. We also obtain the boundary between the two types of the crossover and find that the first-order regime is greatly suppresed by the longitudinal field in accordance with experiment .
According to the simple quantum transition-state theory, the escape-rate in the temperature range $`T\mathrm{\Delta }U`$ is given by
$`\mathrm{\Gamma }(T)={\displaystyle \frac{1}{Z_0}}{\displaystyle \underset{m}{}}\mathrm{\Gamma }_m\mathrm{exp}\left({\displaystyle \frac{E_mU_{\mathrm{min}}}{k_BT}}\right),`$ (1)
where $`\mathrm{\Gamma }_m=\omega (E_m)W(E_m)/(2\pi )`$, $`\omega (E_m)`$ is the frequency of oscillations at the energy $`E_m`$, $`W(E_m)`$ quantum transition probabilities, and $`Z_0`$ the partition function in the well. It is evident from Eq. (1) that the nonanalyticity of the rate is not expected around $`T_0`$ for any escape process, even though the crossover can be sharp because of the exponential dependences of $`\mathrm{\Gamma }_m`$ and thermal populations on the parameters. The rate for the first-order crossover, $`\mathrm{\Gamma }_I(T)`$ is not analytic in the exponential approximation, when the escape rate is approximately given by the dominant term in the summation. However, the summation over all energy levels in Eq. (1) smoothens this non-analyticity.
The model with an arbitrarily directed magnetic field is described by the Hamiltonian
$`=DS_z^2H_zS_zH_xS_x,`$ (2)
The zero field Hamiltonian has uniaxial symmetry with easy axis along $`z`$ and hard plane, $`xy`$. $`H_z`$ is the longitudinal field which affects the height of the energy barrier, and $`H_x`$ is the transverse field which is responsible for quantum tunneling as well as for the reduction of the barrier. In the first approximation, this spin model describes the magnetic molecule $`\mathrm{Mn}_{12}`$.
Within the thermally assisted model, the magnetization reversal occurs by quantum tunneling from thermally excited magnetic levels at magnetic fields which bring into resonance the levels $`m`$ and $`m^{}`$ belonging to different potential wells. Denoting $`m`$ to be the escape level from the metastable well, the resonance condition is that the levels $`m`$ and $`m^{}`$ have the same energy when $`H_z=kD`$ where $`m^{}=mk`$ and $`k`$ is the bias index. The escape rate from any level is proportional to the product of the probability of the thermal occupation of that level and the probability of quantum tunneling from the level. In this respect, the dominat level for a given temperature is determined by the function
$`f(m)={\displaystyle \frac{\pi (\mathrm{\Delta }E_{mm^{}})^2}{2\omega (E_m)}}\mathrm{exp}\left({\displaystyle \frac{E_mU_{\mathrm{min}}}{k_BT}}\right),`$ (3)
where it is assumed that the sum of the linewidths of the $`m`$-th and $`m^{}`$-th levels substantially exceeds the level spacing $`\omega _m^{}(=E_{m^{}+1}E_m^{})`$. Here $`E_m(=Dm^2H_zm)`$ is the energy level of the spin system without transverse field, and $`\mathrm{\Delta }E_{mm^{}}`$ is the splitting of the pair of in-resonance levels $`m`$ and $`m^{}`$ on the opposite side of the anisotropy barrier. It is seen from Eq. (3) that, since the escape rate decreases exponentially with decreasing temperature, larger longitudinal fields are necessary at lower temperature to produce an observable tunneling rate.
Now, let us calculate the dominant level $`m_d`$ which maximizes $`f(m)`$. In order to do that, we first need to find the range of $`m`$ in the metastable well for a given transverse and longitudinal field, i.e., $`Smm_t0`$, when $`m_t`$ is the level which is near the top of the barrier and inside the metastable well. For $`H_x=0`$ and $`H_z=kD`$, simple analysis shows that $`m_t^0=[k/2]1`$ where $`[x]`$ gives the integer part of $`x`$. Since the height of barrier decreases with increasing $`H_x`$, one expects $`m_t<m_t^0`$. To find the value of $`m_t`$ we express Eq. (2) in the spherical coordinate and study the energy in the easy plane given by
$`E(\theta ,\varphi =0)=DS^2(\mathrm{cos}^2\theta +2h_x\mathrm{sin}\theta +2h_z\mathrm{cos}\theta ),`$ (4)
where $`h_{x,z}=H_{x,z}/(2DS)`$. Writing the height of the barrier as $`\mathrm{\Delta }UDS^2(\mathrm{\Delta }u)`$, $`m_t`$ is determined by the relation
$`m_t=\left[{\displaystyle \frac{k}{2}}+S\sqrt{(1h_z)^2\mathrm{\Delta }u}\right]1,`$ (5)
where it is noted that, since $`\mathrm{\Delta }u=(1h_z)^2`$ in the absence of the transverse field, we obtain $`m_t=m_t^0`$. Also, $`m_t=S\sqrt{h_x(2h_x)}`$ at $`h_z=0`$. Numerical calculation of $`\mathrm{\Delta }u`$ of Eq. (4) leads to the results for $`m_t(h_x,h_z)`$ shown in Fig. 1.
Next, we consider the frequency of the real-time oscillations, $`\omega (E_m)`$ at the energy $`E_m`$. This quantity cannot be calculated with the use of the energy (4) in the spherical coordinate system in which the physical quantity which is equivalent to the mass of the system in a one dimensional case is unknown. Thus, for the mapping of the spin problem onto a particle one, the corresponding energy-dependent frequency is of the form
$`\omega (E)=2\pi \left({\displaystyle \frac{1}{\sqrt{D}}}{\displaystyle _{x_1(E)}^{x_2(E)}}{\displaystyle \frac{dx}{\sqrt{EU(x)}}}\right)^1`$ (6)
where $`x_{1,2}`$ are turning points in the particle potential $`U(x)`$ for a given energy $`E`$. Introducing the parameter $`p=(U_{\mathrm{sad}}E_m)/\mathrm{\Delta }U`$, the specific form of $`p`$ at a small value of $`h_x`$ becomes,
$`p=\left({\displaystyle \frac{\frac{m}{S}+h_z}{1h_z}}\right)^2,`$ (7)
where $`p=0`$ at the top of the barrier. This gives $`m/S=h_z`$, i.e., $`m=k/2`$ which is related to $`m_t`$ discussed previously. After some trivial manipulation of Eq. (6), we obtain the dependence of $`\omega `$ on $`h_x`$, $`h_z`$, and $`m`$. As a result, the frequency of oscillation in $`f(m)`$ can be numerically deduced by taking $`h_x0`$. In this limit the frequency also can be computed for the spin model having the energy levels $`E_m=Dm^2H_zm`$ as the inversed density of states, i.e., the energy difference between neighboring levels at energy $`E`$. This is simply given by $`\omega _mD(2m+k)`$ for $`S1`$. Now, in order to calculate the level splitting $`\mathrm{\Delta }E_{mm^{}}`$, we first consider the formula of the perturbation theory
$`\mathrm{\Delta }E_{mm^{}}`$ $`=`$ $`{\displaystyle \frac{2DS^{m^{}m}}{[(m^{}m1)!]^2}}`$ (9)
$`\times \left[{\displaystyle \frac{(S+m^{})!(Sm)!}{(Sm^{})!(S+m)!}}\right]^{1/2}h_x^{m^{}m}.`$
This formula is compared with the results from the direct numerical diagonalization. As is illustrated in Table I, there is a disagreement between them for levels with $`m6`$. However, noting that $`m_t(h_x=0,h_z=0.1)=2`$ is shifted to $`m_t(h_x=0.05,h_z=0.1)=6`$, the level $`m=6`$ or $`7`$ is important to study the type of the crossover and the escape-rate around $`T_0`$. Accordingly, we will perform the numerical diagonalization for the level splitting of $`f(m)`$.
Now, we are in a position to calculate the dominant level $`m_d`$ which is determined by the maximal value of the function (3) for a given temperature. Within the thermally assisted tunneling model the system tunnels through the levels between the bottom and the top. In this process $`m_d(T)`$ behaves in two different ways. One way is that $`m_d`$ changes continuously from $`S`$ to $`m_t`$, whereas the other way is that it performs some discontinuous jump. We call the former - the second-order crossover and the latter - the first-order crossover. As is shown in Fig. 2, for the resonant field $`h_z=0.1`$ we have the first-order crossover for $`h_x=0.05,0.1`$ and the second-order crossover for $`h_x=0.15`$. Also, for $`h_z=0.4`$ ( $`H_z3.28`$ Tesla in $`\mathrm{Mn}_{12}`$ ), the abrupt shift occurs at $`h_x=0.04`$ and the corresponding dominant level changes by 2 ($`m_d=10`$ and $`m_d=8`$) in the range of temperature ($``$ 0.1 K – $``$ 1 K ). Strikingly, this feature is observed in a recent experiment, in which the step positions shift abruptly at low temperature and high magnetic field. Employing these schemes in the whole range of $`h_z`$, we obtain the phase boundary for the values of the transverse field, which is shown by the symbols in Fig. 3. In the quasiclassical method, the order of the quantum-classical escape-rate crossover was determined by the sign of the coefficient in the expansion of the imaginary-time action near the top of the barrier. In the perturbation method, the behavior of the result (9) is inserted into the calculation of the dominant level $`m_d`$ in the metastable well, whose behavior determines the type of the order. Using this method, the first-order crossover is found to be suppresed as compared to the quasiclassical method. This was noticed in Ref. . However, as discussed previously, the corresponding level splitting is not quite correct, especially, near the top of the barrier, in which range the perturbation fails. The correct calculation based on the diagonalization method shows that the first-order regime is even more suppresed as compared to the perturbative results. For example, in the unbiased case, the phase boundary becomes at $`h_x=`$ 0.114, 0.139, and 0.25 for the diagonalization, perturbation, and quasiclassical method, respectively. In Fig. 3, there is no data point beyond $`h_z=0.8`$. The reason is that $`m_t=9`$ in this region and thereby it is meaningless to ask whether the shift is continuous or discontinuous in this region.
The values of the crossover temperature $`T_0^{(c)}`$ at the phase boundary between first- and second-order crossover, which have been obtained by the diagonalization method described above, are shown in Fig. 4. Comparing this with two other methods, the result based on the diagonalization method is the smallest in the whole range of the bias field. For example, in the unbiased case, we have $`T_0^{(c)}/(DS)=`$ 0.117, 0.124, and 0.137 for the diagonalization, perturbation, and quasiclassical method, respectively.
Now, for the relaxation at resonance we will present the results of numerical calculation for the escape rate, $`\mathrm{\Gamma }(T)`$. As is shown in Fig. 5 (as, e.g., for $`h_z=0.1`$), the escape rate in the first-order region ($`h_x=0.05`$ or 0.1) changes sharply, as shows the comparison with the escape rate in the second-order one ($`h_x=0.15`$ or 0.2). Also, for $`h_z=0.4`$ we can clearly distinguish the behavior of the rates in two different regimes, e.g., $`h_x=0.04`$ and 0.1, and its trend continues in the whole range of the bias field. The origin of these behaviors is that, as described above, $`m_d`$ changes discontinuously around $`T_0`$ for the first-order crossover, while it does continuously for the second-order one. In other words, since the rate (1) is sensitive to the change of $`m`$, the abrupt jump in the first-order crossover induces sharp increase, e.g., $`\mathrm{\Gamma }(h_x=0.05,h_z=0.1)/\mathrm{\Gamma }(h_x=0.2,h_z=0.1)10^6`$ at $`T/(DS)0.12`$, and $`\mathrm{\Gamma }(h_x=0.04,h_z=0.4)/\mathrm{\Gamma }(h_x=0.1,h_z=0.4)10^5`$ at $`T/(DS)0.11`$.
For the uniaxial spin model considered in this paper, both type of crossover can be realized in the molecular magnet $`\mathrm{Mn}_{12}`$, and the situation can be controlled by the longitudinal and transverse field. For $`\mathrm{Mn}_{12}`$, $`D0.55`$K, and $`\mathrm{\Delta }U=DS^255`$K at $`h_x=h_z=0`$. The critical field in the $`x`$\- or $`z`$-direction is $`H_{xc}=H_{zc}=2DS/(g\mu _B)8.2`$T, and the longitudinal field for the resonance is $`H_0=D/(g\mu _B)0.41`$T. At these fields, the magnetic relaxation in $`\mathrm{Mn}_{12}`$ can occur on measurement time scales, and gives rise to the different behavior of the dominant level and the rate in magnetization depending on $`H_x`$ and $`H_z`$. Furthermore, in the unbiased case the first-order crossover can be observed in the field range $`0<H_x<0.93`$ Tesla and the crossover region occurs at the temperature range $`0.1\mathrm{K}<T<1\mathrm{K}`$. This field range decreases with increasing the bias field. Even though we presented the results of $`m_d`$ only for $`h_z=0.1`$ and 0.4 in Fig. 2, we have found that the levels which, at a certain temperature, change by more than 1 are $`m=7`$, 8 or 9 depending on the value of the resonant field.
In conclusion, we have studied the quantum-classical crossover of the escape-rate of a uniaxial spin model with an arbitrarily directed field. Employing the diagonalization method, we have obtained the dominant level for the thermally assisted tunneling and dependence of the escape rate on temperature. In comparison to the previously studied models, the first-order region is greatly suppressd, but still observable in molecular magnets. It is also found that the first-order crossover is fairly sharp while the second-order the crossover is smooth. This is found to be strongly related with the discontinuous (continuous) jump of $`m_d`$ in lower (higher) field $`h_x`$. These results have been applied to the high-spin molecule, $`\mathrm{Mn}_{12}`$. They are also relevant to the study of nanoparticles.
G.-H. K. is grateful to A. Garg for useful discussions. E. C. acknowledges support from the NSF Grant DMR-9978882. G.-H. K. was supported by grant No. 1999-1-114-002-5 from the Interdisciplinary Research Program of the KOSEF.
| $`m`$ | pert. | diag. |
| --- | --- | --- |
| -6 | $`0.156`$ | $`8.53\times 10^2`$ |
| -7 | $`1.01\times 10^3`$ | $`7.38\times 10^4`$ |
| -8 | $`2.72\times 10^6`$ | $`2.26\times 10^6`$ |
| -9 | $`3.14\times 10^9`$ | $`2.78\times 10^9`$ |
| -10 | $`1.40\times 10^{12}`$ | $`1.28\times 10^{12}`$ |
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# The Role of Spin-Dependent Interface Scattering in Generating Current-Induced Torques in Magnetic Multilayers
## I Introduction
Stacks of alternating ferromagnetic and nonmagnetic metal layers exhibit giant magnetoresistance (GMR), because their electrical resistance depends strongly on whether the moments of adjacent magnetic layers are parallel or antiparallel. This effect has allowed the development of new kinds of field-sensing and magnetic memory devices. The cause of the GMR effect is that conduction electrons are scattered more strongly by a magnetic layer when their spins lie antiparallel to the layer’s magnetic moment than when their spins are parallel to the moment. Devices with moments in adjacent magnetic layers aligned antiparallel thus have a larger overall resistance than when the moments are aligned parallel, giving rise to GMR. This paper discusses the converse effect: just as the orientations of magnetic moments can affect the flow of electrons, then by Newton’s third law, a polarized electron current scattering from a magnetic layer can have a reciprocal effect on the moment of the layer. As proposed by Berger and Slonczewski, an electric current passing perpendicularly through a magnetic multilayer may exert a torque on the moments of the magnetic layers. This effect which is known as “spin transfer”, may, at sufficiently high current densities, alter the magnetization state. It is a separate mechanism from the effects of current induced magnetic fields. Experimentally, spin-current-induced magnetic excitations such as spin-waves, and stable magnetic reversal, have been observed in multilayers, for current densities greater than $`10^7A/cm^2`$.
The spin-transfer effect offers the promise of new kinds of magnetic devices, and serves as a new means to excite and to probe the dynamics of magnetic moments at the nanometer scale. In order to controllably utilize these effects, however, it is necessary to achieve a better quantitative understanding of these current-induced torques. Slonczewski has presented a derivation of spin-transfer torques using a 1-D WKB approximation with spin-dependent potentials, but his calculations only take into account electrons which are either completely transmitted or completely reflected by the magnetic layers. For real materials the degree to which an electron is transmitted through a magnetic/nonmagnetic interface depends sensitively on the matching of the band structures across the interface. It is the goal of this paper to incorporate such band structure information together with the effect of multiple reflections between the ferromagnetic layers, into a more quantitative theory of the torques generated by spin-transfer. This could be done using the formalism of Brataas et al, which is based on kinetic equations for spin currents. Instead we choose to employ a modified Landauer-Büttiker formalism, in which we model the ferromagnetic layers as generalized spin-dependent scatterers. The calculations are carried out for a quasi-one dimensional geometry, for which we derive formulas for the torque generated on the magnetic layers when a current is applied to the system, for either ballistic or diffusive non-magnetic layers. The main difference between our approach and Ref. is that in our case, scattering in the normal layer is phase coherent, whereas Ref. assumes phase relaxation. However, in the case of a diffusive normal metal layer and for a large number of transverse modes, the two approaches would give the same answer.
The paper is organized as follows: in Section II, we present an intuitive picture (adapted from) of how spin-dependent scattering of a spin-polarized current produces a torque on a magnetic element. Section III is devoted to the introduction of the scattering matrix formalism for the spin-flux. This formalism is then used in Section IV to calculate the torque in a Ferromagnet-Normal-Ferromagnet (FNF) system where the normal part is disordered (diffusive). Section V contain a discussion of the results. Details of our calculation are presented in Appendix A. In appendix B, we derive the torque for an FNF system where transport in the normal layer is ballistic, rather than diffusive.
## II Physical Idea
In this section we will present a simple intuitive picture of the physics behind the spin-transfer effect. The connection between current-induced spin-transfer torques and the spin-dependent scattering that occurs when electrons pass through a magnetic/nonmagnetic interface can be illustrated most simply by considering the case of a spin-polarized current incident perpendicularly on a single thin ferromagnetic layer F, as shown in Figure 1. The layer lies in the $`yz`$ plane, with its magnetic moment uniformly pointed in the $`+z`$ direction, and we assume that the current is spin-polarized in the $`zx`$ plane at an angle $`\theta `$ to the layer moments. The incoming electrons can therefore be considered as a coherent linear superposition of basis states with spin in the +z direction (amplitude $`\mathrm{cos}(\theta /2)`$) and -z direction (amplitude $`\mathrm{sin}(\theta /2)`$). For this initial discussion we will assume that the layer is a perfect spin-filter, so that spins aligned with the layer moments are completely transmitted through the layer, while spins aligned antiparallel to the layer moment are completely reflected. For incident spins polarized at an angle $`\theta `$, the average outgoing current will have the relative weights $`\mathrm{cos}^2(\theta /2)`$ polarized in the $`+z`$ direction and transmitted to the right and $`\mathrm{sin}^2(\theta /2)`$ polarized in the $`z`$ direction and reflected to the left. Consequently, both of the outgoing electron spin fluxes (transmitted and reflected) lie along the $`z`$ axis, while the incoming (incident) electron flux has a component perpendicular to the magnetization, along the $`x`$ axis, with magnitude proportional to $`\mathrm{sin}\theta `$. This $`x`$-component of angular momentum must be absorbed by the layer in the process of filtering the spins. Because the spin-filtering is ultimately governed by the $`sd`$ exchange interaction between the conduction electrons and the magnetic moments of the layer, the angular momentum is imparted to the layer moments and produces a torque on them. This exchange torque, which is proportional to the electron current through the layer and to $`\mathrm{sin}\theta `$, is in the direction to align the moments with the polarization of the incident spin current.
The symmetry of this model precludes any generation of torque from the spin-filtering of a current of unpolarized electrons. To generate the effect, then, a second ferromagnetic layer is needed to first spin-polarize the current, see Fig. 2. In that case, spin angular momentum is transferred from one layer to the flowing electrons and then from the electrons to the second layer. However, the torques on the two layers are not equal and opposite, as spin angular momentum carried by the electrons can also flow away from the layers to infinity, see Fig. 2.
The presence of this second layer has the additional effect of allowing for multiple scattering of the electrons between the two layers, which gives rise to an explicit asymmetry of the torque with respect to current direction. This asymmetry is an important signature which can be used to distinguish spin-transfer-induced torques from the torques produced by current-generated magnetic fields. To see how the asymmetry arises, consider the ferromagnet–normal-metal–ferromagnet (FNF) junction shown in Figure 2. It consists of two ferromagnetic layers, F<sub>a</sub> and F<sub>b</sub>, with moments pointing in directions $`\widehat{m}_a`$ and $`\widehat{m}_b`$, separated by a normal metal spacer N. Normal metal leads on either side of the trilayer inject an initially unpolarized current into the system. When the current enters the sample from the left (Fig. 2a), electrons transmitted through F<sub>a</sub> will be polarized along $`\widehat{m}_a`$. As long as the normal metal spacer is smaller than the spin-diffusion length (100 nm for Cu), this current will remain spin-polarized when it impinges on F<sub>b</sub> and will exert a torque on the moment of F<sub>b</sub> in a direction so as to align $`\widehat{m}_b`$ with $`\widehat{m}_a`$. Repeating the argument for F<sub>b</sub>, we find that the spin of the electrons reflected from layer F<sub>b</sub> is aligned antiparallel to $`\widehat{m}_b`$, and, hence in turn, exerts a torque on the moment of F<sub>a</sub> trying to align $`\widehat{m}_a`$ antiparallel with $`\widehat{m}_b`$. (Subsequent multiple reflections of electrons between F<sub>a</sub> and F<sub>b</sub> can serve to reduce the magnitudes of the initial torques, but they do not eliminate or reverse them, as the electron flux is reduced upon each reflection.) When the current is injected from the right, the directions of the torques are reversed: Now the flow of electrons exerts a torque on F<sub>a</sub> trying to align its moment parallel with $`\widehat{m}_b`$, while it exerts a torque on F<sub>b</sub> so as to force the moment in layer F<sub>b</sub> antiparallel with $`\widehat{m}_a`$.
In the remainder of the paper, we assume that $`\widehat{m}_b`$ points in the $`+z`$ direction, while $`\widehat{m}_a`$ differs by a small angle $`\theta `$ in the $`xz`$ plane. (For thin films, demagnetizing forces will in general cause the $`yz`$ plane to be preferred, but this produces no change in our argument. We present the case that is easier to draw.) The overall effect of a left-going flow of electrons then, is to exert a torque $`\stackrel{}{\tau }_b`$ on F<sub>b</sub> in the $`+x`$ direction. If we reverse the current, so that electrons pass through F<sub>b</sub> first (Fig. 2b), the torque on F<sub>b</sub> is only exerted by the electrons after they have been reflected from F<sub>a</sub>. As seen before, the electrons reflected from F<sub>a</sub> have polarizations opposite to $`\widehat{m}_a`$, so that the torque on F<sub>b</sub> is in the $`x`$ direction.
In Refs. and , the layer F<sub>a</sub> was taken to be much thicker than F<sub>b</sub>, so that intralayer exchange and anisotropy forces will hold the orientation of $`\widehat{m}_a`$ fixed. In that case, one is only interested in the torque on F<sub>b</sub>, which serves to align $`\widehat{m}_b`$ either parallel or antiparallel with the fixed moment $`\widehat{m}_a`$ depending on the current direction. This asymmetric current response has been employed in both a point-contact geometry and in a thin-film pillar geometry to switch the moments in FNF trilayers from a parallel to an antiparallel configuration by a current pulse in one direction, and then from antiparallel to parallel by a reversed current. For weakly-interacting layers, either orientation can be stable in the absence of an applied current, so that the resistance versus current characteristic is hysteretic, and the devices can function as simple current-controlled memory elements.
## III Spin flux and torque in the scattering approach
Treating the ferromagnetic layers as perfect spin filters provides important qualitative insights into spin-transfer, but for a complete qualitative and quantitative picture, a more general approach is required. In this section, we introduce a scattering matrix description of the FNF junction which allows us to deal with non-ideal (magnetic and non-magnetic) layers. Our goal is to relate the torque $`\stackrel{}{\tau }_b`$ exerted on layer F<sub>b</sub> by an unpolarized incident electron beam to the scattering properties of the layers. Although we shall restrict our formulas to the FNF junction (see Fig. 3), our method is applicable for an arbitrary array of magnetic-nonmagnetic layers.
We first introduce the spin flux $`\stackrel{}{J}`$ in the $`x`$-direction (the direction of current flow):
$$\stackrel{}{J}(x)=\frac{\mathrm{}^2}{2m}\mathrm{Im}𝑑y𝑑z\left[\varphi ^{}(x)\stackrel{}{\sigma }\frac{}{x}\varphi (x)\right]$$
(1)
where $`\varphi (x)`$ is a spinor wavefunction and $`\stackrel{}{\sigma }`$ the vector of Pauli matrices,
$$\varphi (x)=\left(\begin{array}{c}\varphi _{}(x)\\ \varphi _{}(x)\end{array}\right),\stackrel{}{\sigma }=\left(\begin{array}{c}\sigma _x\\ \sigma _y\\ \sigma _z\end{array}\right).$$
Note that although Eq.(1) bears close formal analogy to the particle current, no local equation of conservation can be written for the spin flux, since in general, the Hamiltonian does not conserve spin. Specifically, the magnetic layers can act as sources and sinks of spin flux, so that the spin flux on different sides of a F layer can be different.
### A Definition of the scattering matrices
Fig. 3 shows the FNF junction where (fictitious) perfect leads (labeled $`0`$, $`1`$, $`2`$ and $`3`$) have been added in between the layers F and N and between the F layers and the electron reservoirs on either side of the sample. The introduction of these leads allows for a description of the system using scattering matrices. In the perfect leads, the transverse degree of freedom are quantized, giving $`N`$ propagating modes at the Fermi level, where $`NA/\lambda _F^2`$ $`A`$ being the cross section area of the junction and $`\lambda _F`$ the Fermi wave length. Expanding the electronic wave function in these modes, we can describe the system in terms of the projection $`\mathrm{\Psi }_{i,L/R}`$ of the wave function onto the left (right) going modes in region $`i`$. The $`\mathrm{\Psi }_{i,L/R}`$ are $`2N`$-component vectors, counting the $`N`$ transverse modes and spin. The amplitudes of the wave function in two neighboring ideal leads are connected through the scattering matrices $`S_b`$, $`S_a`$ and $`S_\mathrm{N}`$, that relate amplitudes of outgoing modes and incoming modes at the layer (see for example Ref. for a review of the scattering matrix approach),
$`\left(\begin{array}{c}\mathrm{\Psi }_{3L}\\ \mathrm{\Psi }_{2R}\end{array}\right)`$ $`=`$ $`S_b\left(\begin{array}{c}\mathrm{\Psi }_{3R}\\ \mathrm{\Psi }_{2L}\end{array}\right),`$ (7)
$`\left(\begin{array}{c}\mathrm{\Psi }_{1L}\\ \mathrm{\Psi }_{0R}\end{array}\right)`$ $`=`$ $`S_a\left(\begin{array}{c}\mathrm{\Psi }_{1R}\\ \mathrm{\Psi }_{0L}\end{array}\right),`$ (12)
$`\left(\begin{array}{c}\mathrm{\Psi }_{2L}\\ \mathrm{\Psi }_{1R}\end{array}\right)`$ $`=`$ $`S_\mathrm{N}\left(\begin{array}{c}\mathrm{\Psi }_{2R}\\ \mathrm{\Psi }_{1L}\end{array}\right).`$ (17)
The scattering matrices $`S_b`$, $`S_a`$ and $`S_\mathrm{N}`$ are $`4N\times 4N`$ unitary matrices. We decompose $`S_b`$ into $`2N\times 2N`$ reflection and transmission matrices,
$$S_b=\left(\begin{array}{cc}r_b& t_b^{}\\ t_b& r_b^{}\end{array}\right),$$
(18)
with similar decompositions of $`S_a`$ and $`S_\mathrm{N}`$. Normalization is done in such a way that each mode carries unit current. Due to the spin degree of freedom, the reflection and transmission matrices can be written in terms of four $`N\times N`$ blocks:
$$r_b=\left(\begin{array}{cc}r_b& r_b\\ r_b& r_b\end{array}\right),$$
(19)
where the subscripts $`,`$ refer to spin up and down in the $`z`$-axis basis.
The scattering matrix of the magnetic layers depends on the angle $`\theta `$ the moments may make with the z-axis. The matrix $`S_a(\theta )`$ is related to $`S_a(\theta =0)`$ through a rotation in spin space:
$`r_a(\theta )`$ $`=`$ $`R_\theta r_a(0)R_\theta ,r_a^{}(\theta )=R_\theta r_a^{}(0)R_\theta ,`$ (20)
$`t_a(\theta )`$ $`=`$ $`R_\theta t_a(0)R_\theta ,t_a^{}(\theta )=R_\theta t_a^{}(0)R_\theta `$ (21)
where
$$R_\theta =\left(\begin{array}{cc}\mathrm{cos}\frac{\theta }{2}& \mathrm{sin}\frac{\theta }{2}\\ \mathrm{sin}\frac{\theta }{2}& \mathrm{cos}\frac{\theta }{2}\end{array}\right)1_N.$$
(22)
The non-magnetic metallic layer will not affect the spin states, i.e., $`r_\mathrm{N}=r_\mathrm{N}=0`$ and $`r_\mathrm{N}=r_\mathrm{N}`$.
We need to keep track of the amplitudes within the system in order to calculate the net spin flux deposited into each magnetic layer. Therefore, we define $`2N\times 2N`$ matrices $`\mathrm{\Gamma }_i^{L/R}`$ and $`\mathrm{\Lambda }_i^{L/R}`$ ($`i=0,1,2,3`$) so that we may express all the $`\mathrm{\Psi }_{i,L/R}`$ as a function of the amplitudes incident from the two electrodes (regions 0 and 3):
$$\left(\begin{array}{c}\mathrm{\Psi }_{iL}\\ \mathrm{\Psi }_{iR}\end{array}\right)=\left(\begin{array}{cc}\mathrm{\Gamma }_{iL}& \mathrm{\Lambda }_{iL}\\ \mathrm{\Gamma }_{iR}& \mathrm{\Lambda }_{iR}\end{array}\right)\left(\begin{array}{c}\mathrm{\Psi }_{0L}\\ \mathrm{\Psi }_{3R}\end{array}\right)$$
(23)
with the convention that $`\mathrm{\Gamma }_{0L}=\mathrm{\Lambda }_{3R}=1`$ and $`\mathrm{\Gamma }_{3R}=\mathrm{\Lambda }_{0L}=0`$. In order to calculate the torque exercised on layer F<sub>b</sub> for a current entering from the left, we need the matrix, $`\mathrm{\Gamma }_{2L}`$. To simplify the notations in the rest of the paper, we write
$$\mathrm{\Omega }\mathrm{\Gamma }_{2L}.$$
(24)
The matrix $`\mathrm{\Omega }`$ relates $`\mathrm{\Psi }_{2L}`$ to the incoming amplitudes $`\mathrm{\Psi }_{0L}`$ coming from the right. To calculate it, we put $`\mathrm{\Psi }_{3R}=0`$, then, using Eq.(III A), we get the equations:
$`\mathrm{\Psi }_{1L}`$ $`=`$ $`t_a^{}\mathrm{\Psi }_{0L}+r_a\mathrm{\Psi }_{1R},`$ (25)
$`\mathrm{\Psi }_{1R}`$ $`=`$ $`t_n\mathrm{\Psi }_{2R}+r_n^{}\mathrm{\Psi }_{1L},`$ (26)
$`\mathrm{\Psi }_{2L}`$ $`=`$ $`r_n\mathrm{\Psi }_{2R}+t_n^{}\mathrm{\Psi }_{1L},`$ (27)
$`\mathrm{\Psi }_{2R}`$ $`=`$ $`r_b^{}\mathrm{\Psi }_{2L},`$ (28)
$`\mathrm{\Psi }_{3L}`$ $`=`$ $`t_b^{}\mathrm{\Psi }_{2L},`$ (29)
from which we obtain:
$$\mathrm{\Omega }=\frac{1}{1r_nr_b^{}}t_n^{}\frac{1}{1r_at_nr_b^{}\frac{1}{1r_nr_b^{}}t_n^{}r_ar_n^{}}t_a^{}.$$
(30)
### B Spin flux response
Let us now connect our system to two unpolarized electron reservoirs on its two sides, as shown in Fig. 4. In equilibrium, the modes in the reservoirs are filled up to the fermi level $`ϵ_F`$. We want to calculate the spin current that is generated when the chemical potential in the left (right) reservoir is slightly increased by $`\delta \mu _3`$ ($`\delta \mu _0`$). The spin current $`\stackrel{}{J_i}`$ is the difference of the left going and right going contributions. For each of the region $`i=0,1,2,3`$, we find from Eq.(1) and Eq.(23):
$$\frac{\stackrel{}{J_i}}{\mu _3}=\frac{1}{4\pi }\mathrm{Re}\left[\mathrm{Tr}\stackrel{}{\sigma }\mathrm{\Gamma }_{iR}\mathrm{\Gamma }_{iR}^{}\mathrm{Tr}\stackrel{}{\sigma }\mathrm{\Gamma }_{iL}\mathrm{\Gamma }_{iL}^{}\right],$$
(31)
and
$$\frac{\stackrel{}{J_i}}{\mu _0}=\frac{1}{4\pi }\mathrm{Re}\left[\mathrm{Tr}\stackrel{}{\sigma }\mathrm{\Lambda }_{iR}\mathrm{\Lambda }_{iR}^{}\mathrm{Tr}\stackrel{}{\sigma }\mathrm{\Lambda }_{iL}\mathrm{\Lambda }_{iL}^{}\right].$$
(32)
Derivation of Eq.(31) and Eq.(32) proceeds analogously to the derivation of the Landauer formula for the conductance.
### C Torque exercised on layer F<sub>b</sub>
If the spin flux on both sides of F<sub>b</sub> (region 2 and 3) is different, then angular momentum has been deposited in the layer F<sub>b</sub>. This creates a torque $`\stackrel{}{\tau _b}`$ on the moment of the ferromagnet,
$$\stackrel{}{\tau _b}=\stackrel{}{J}_3\stackrel{}{J}_2.$$
(33)
Setting $`\delta \mu _0=eV_0`$, we have:
$$\frac{\stackrel{}{\tau }_b}{V_0}=\frac{e}{4\pi }\mathrm{Re}\mathrm{Tr}_{2N}\left[\stackrel{}{\mathrm{\Sigma }}\mathrm{\Omega }\mathrm{\Omega }^{}\right],$$
(34)
with
$$\stackrel{}{\mathrm{\Sigma }}=\stackrel{}{\sigma }t_b^{^{}}\stackrel{}{\sigma }t_b^{}r_b^{}\stackrel{}{\sigma }r_b^{}.$$
(35)
This equation can be simplified further if the spin-transfer effect is due entirely to spin-filtering (as argued by Slonczewski) as opposed to spin-flip scattering of electrons from the magnetic layers. That is, if we assume that $`r_b=r_b=r_a(\theta =0)=r_a(\theta =0)=0`$, then:
$$\frac{\tau _b^x}{V_0}=\frac{e}{2\pi }\mathrm{Re}\mathrm{Tr}_N[(\mathrm{\Omega }_{}\mathrm{\Omega }_{}^{}+\mathrm{\Omega }_{}\mathrm{\Omega }_{}^{})$$
$$(1r_b^{}r_b^{}t_b^{}t_b^{})]$$
(36)
We will comment briefly in the conclusion of this paper about the physical implications of including the off-diagonal spin-flip terms, as well.
We note that, as there is no spin flux conservation in this system, $`\stackrel{}{J_i}/\mu _3`$ can be different from $`\stackrel{}{J_i}/\mu _0`$ and, hence, there can be a non zero spin flux even when the chemical potentials are identical in the two reservoirs. The existence of a zero-bias spin-flux and the resulting torques reflect the well-known itinerant-electron-mediated exchange interaction (a.k.a. the RKKY interaction) between two ferromagnetic films separated by a normal-metal spacer. This interaction can in fact be understood within a scattering framework. The zero-bias torque has to be added to the finite-bias contribution (given by Eq.(36)). Since the former is typically a factor $`N^1`$ smaller and vanishes upon ensemble averaging (see section IV and Ref.), we henceforth neglect the zero-bias contribution to the torque and restrict our attention to the bias induced torque, for which we have
$$\frac{\stackrel{}{\tau }_b}{V_0}=\frac{\stackrel{}{\tau }_b}{V_3}$$
up to a correction of order $`N^1`$.
## IV Averaging over the normal layer
Via Eq.(36), the torque on the moments of the ferromagnetic layers F<sub>a</sub> and F<sub>b</sub> not only depends on the scattering matrices $`S_a`$ and $`S_b`$ of these layers, but also of the scattering matrix $`S_\mathrm{N}`$ of the normal metal layer in between. If the normal layer is disordered, $`\stackrel{}{\tau }_a`$ and $`\stackrel{}{\tau }_b`$ depend on the location of the impurities; if N is ballistic the torque depend sensitively on the electronic phase shift accumulated in N. In general, sample to sample fluctuations of the torque will be a factor $`N^1`$ smaller than the average. Hence, if $`N`$ is large ($`N>10^3`$ in the experiments of Ref.), the torque is well characterized by its average. In this section, we average over $`S_\mathrm{N}`$ for the case where N is disordered. The case of ballistic N is addressed in appendix B. After averaging, the zero-bias spin transfer current, corresponding to the RKKY interaction described above, vanishes, and only the torque caused by the electron current remains. Because all effects of quantum interference in the N layer will disappear in the process of averaging (to leading order in the number of modes $`N`$), the results we derive are unchanged if the reflection and transmission matrices include processes in which the energy of the electron changes during scattering, in addition to the elastic proceses normally considered in scattering matrix calculations.
### A Averaged torque
The scattering matrix of the normal layer can be written using the standard polar decomposition:
$$S_n=\left(\begin{array}{cc}U& 0\\ 0& V^{}\end{array}\right)\left(\begin{array}{cc}\sqrt{1T}& i\sqrt{T}\\ i\sqrt{T}& \sqrt{1T}\end{array}\right)\left(\begin{array}{cc}U^{}& 0\\ 0& V\end{array}\right)$$
(37)
where $`U,V,U^{}`$ and $`V^{}`$ are $`2N\times 2N`$ unitary matrices and T is a diagonal matrix containing the eigenvalues of $`t_nt_n^{}`$. Since $`S_\mathrm{N}`$ is diagonal in spin space, we find that $`U`$, $`U^{}`$, $`V`$ and $`V^{}`$ are block diagonal.
$$U=\left(\begin{array}{cc}u& 0\\ 0& u\end{array}\right),U^{}=\left(\begin{array}{cc}u^{}& 0\\ 0& u^{}\end{array}\right),$$
(38)
and similar definitions for $`v`$ and $`v^{}`$. In the isotropic approximation, the $`N\times N`$ unitary matrices $`u,u^{},v`$ and $`v^{}`$ are uniformly distributed in the group $`U(N)`$. (The outer matrices in Eq.(37) thus mix the modes in a ergodic way while the central matrix contains the transmission properties of the layer, which determine the average conductance of N.)
We want to average Eq.(36) over both the unitary matrices and $`T`$. A diagrammatic technique for such averages has already been developed in Ref. and can be used to calculate $`\stackrel{}{\tau }_b/V_0`$ in leading order in $`1/N`$. It is a general property of such averages that the fluctuations are a factor of order $`N`$ smaller than the average. This justifies our statement above, that the ensemble averaged torque is sufficient to characterize the torque exerted on a single sample. Details of the calculation are presented in Appendix A.
The resulting expression for $`\stackrel{}{\tau }_b/V_0`$ can be written in a form very similar to the one for Eq.(36) if one uses a notation that involves $`4\times 4`$ matrices. To be specific, to each $`2N\times 2N`$ matrix $`A`$ appearing in Eq.(36) and Eq.(30), we assign a $`4\times 4`$ matrix $`\widehat{A}`$ as,
$$\widehat{A}=\frac{1}{N}\mathrm{Tr}_N\left[AA^{}\right],$$
(39)
where $`\mathrm{Tr}_N`$ means that the trace has been taken in each the $`N\times N`$ blocks, or in extenso:
$$\widehat{A}=\frac{1}{N}\mathrm{Tr}_N\left(\begin{array}{cccc}A_{}A_{}^{}& A_{}A_{}^{}& A_{}A_{}^{}& A_{}A_{}^{}\\ A_{}A_{}^{}& A_{}A_{}^{}& A_{}A_{}^{}& A_{}A_{}^{}\\ A_{}A_{}^{}& A_{}A_{}^{}& A_{}A_{}^{}& A_{}A_{}^{}\\ A_{}A_{}^{}& A_{}A_{}^{}& A_{}A_{}^{}& A_{}A_{}^{}\end{array}\right).$$
(40)
We also define $`\widehat{\stackrel{}{\mathrm{\Sigma }}}`$ by,
$$\widehat{\stackrel{}{\mathrm{\Sigma }}}=\mathrm{Tr}_N\left(\begin{array}{cccc}\stackrel{}{\mathrm{\Sigma }}_{}& \stackrel{}{\mathrm{\Sigma }}_{}& \stackrel{}{\mathrm{\Sigma }}_{}& \stackrel{}{\mathrm{\Sigma }}_{}\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ \stackrel{}{\mathrm{\Sigma }}_{}& \stackrel{}{\mathrm{\Sigma }}_{}& \stackrel{}{\mathrm{\Sigma }}_{}& \stackrel{}{\mathrm{\Sigma }}_{}\end{array}\right).$$
(41)
The average over the transmission eigenvalues $`T`$ follows if we note that the average of a function is the function of the average, to leading order in $`1/N`$. Thus the average over $`T`$ amounts to the replacement
$$\widehat{t}_n=\frac{g_\mathrm{N}}{N}𝟙_\mathrm{𝟜}\mathrm{and}\widehat{𝕣}_𝕟=\left(\mathrm{𝟙}\frac{𝕘_\mathrm{N}}{}\right)𝟙_\mathrm{𝟜},$$
(42)
where $`g_\mathrm{N}`$ is the conductance of the normal layer and $`𝟙_\mathrm{𝟜}`$ is the $`4\times 4`$ unit matrix. Using these “hat” matrices, the result has now the simple form:
$$\frac{\stackrel{}{\tau }_b}{V_0}=\frac{e}{4\pi }\mathrm{Re}\mathrm{Tr}_4\left[\widehat{\stackrel{}{\mathrm{\Sigma }}}\widehat{\mathrm{\Omega }}\right],$$
(43)
where (compare to Eq.(30)),
$$\widehat{\mathrm{\Omega }}=\frac{1}{1\widehat{r}_n\widehat{r}_b^{}}\widehat{t}_n^{}\frac{1}{1\widehat{r}_a\widehat{t}_n\widehat{r}_b^{}\frac{1}{1\widehat{r}_n\widehat{r}_b^{}}\widehat{t}_n^{}\widehat{r}_a\widehat{r}_n^{}}\widehat{t}_a^{}.$$
(44)
Equation (43) is the main result of this paper. In the absence of spin-flip scattering, it reduces to
$`{\displaystyle \frac{\tau _b^x}{V_0}}`$ $`=`$ $`{\displaystyle \frac{e}{2\pi }}\mathrm{Re}[(\widehat{\mathrm{\Omega }}_{3,1}+\widehat{\mathrm{\Omega }}_{3,4})`$ (45)
$`\times `$ $`\mathrm{Tr}_N(1r_b^{}r_b^{}t_b^{}t_b^{})].`$ (46)
The same formalism can be used to calculate the conductance $`g`$ of the system using the Landauer formula. One gets:
$$g=\frac{Ne^2}{h}\left[\widehat{t^{}}_{1,1}+\widehat{t^{}}_{1,4}+\widehat{t^{}}_{4,1}+\widehat{t^{}}_{4,4}\right],$$
(47)
$`t^{}`$ being the total transmission matrix:
$$t^{}=t_b^{}\mathrm{\Omega }.$$
(48)
We would like to note that, while our theory started from a fully phase coherent description of the FNF trilayer, including the full $`4N\times 4N`$ scattering matrices of the FN interfaces, the final result can be formulated in term of $`2\times 4`$ parameters, represented by the matrices $`\widehat{r}_a`$ and $`\widehat{r}_b^{}`$ ( $`2\times 16`$ parameters in case of spin-flip scattering). Such a reduction of the number of degrees of freedom was also found by Brataas et al., although their starting point is an hybrid ferromagnetic-normal metal circuit with incoherent nodes. This confirms the statement at the beginning of this section, that for a diffusive normal-metal spacer all effects of quantum interferences are washed out. The difference between our approach and the one of Ref. is important in the case of the ballistic normal layer, see Appendix B.
### B Symmetries
Before we proceed with a further analysis of Eq.(43), we identify the different symmetries of the torque. Due to the conservation of current, the total torque deposited on the full system is anti-symmetric with respect to current direction:
$$\frac{\stackrel{}{\tau }_b}{V_0}+\frac{\stackrel{}{\tau }_a}{V_0}=\left[\frac{\stackrel{}{\tau }_b}{V_3}+\frac{\stackrel{}{\tau }_a}{V_3}\right].$$
(49)
Eq.(49) holds before averaging. But, as pointed out in section III, equality for each of the torques $`\stackrel{}{\tau }_a`$ and $`\stackrel{}{\tau }_b`$ separately holds only after averaging,
$$\frac{\stackrel{}{\tau }_b}{V_0}=\frac{\stackrel{}{\tau }_b}{V_3};$$
(50)
sample to sample fluctuations of $`\stackrel{}{\tau }_b/V_0`$ and $`\stackrel{}{\tau }_b/V_3`$ of relative order $`1/N`$ are in general different. Thus, for $`N1`$, our calculation can be used to compute the linear response of the torque to a small bias voltage:
$$\stackrel{}{\tau }_b=\frac{\stackrel{}{\tau }_b}{V_0}(V_0V_3).$$
(51)
In our geometry, where F<sub>a</sub> and F<sub>b</sub> are in $`xz`$ plane, the only non-zero component of the torque is $`\tau _b^x`$. The torque vanishes when the moments are completely aligned or antialigned (all the matrices are diagonal in spin space and therefore no x-component of the spin can be found). Around these two limits, the torque is symmetric in respect to the angle ($`\theta \theta `$ and $`\pi \theta \pi +\theta `$). There is no symmetry between $`\theta `$ and $`\pi \theta `$. In addition, the two layers are not equivalent and exchanging the scattering matrices of F<sub>a</sub> and F<sub>b</sub> also changes the torque.
### C Discussion of some limiting cases
Eq.(43) can be simplified in some particular cases. Let us start with the simplest case of ideal spin filters, so that majority (minority) spins are totally transmitted (reflected) by either layer. Equation (43) then reduces to
$$\frac{\tau _b^x}{V_0}=\frac{e}{4\pi }\frac{g_\mathrm{N}\mathrm{sin}\theta }{3+\mathrm{cos}\theta }=\frac{h}{4\pi e}g\frac{\mathrm{tan}\theta /2}{2},$$
(52)
where $`g`$ is the average magnetoconductance, cf Eq.(47),
$$g=\frac{e^2}{h}g_\mathrm{N}\frac{4\mathrm{cos}^2\theta /2}{3+\mathrm{cos}\theta }.$$
(53)
Equation (52) reproduces a result of Slonczewski. As expected, for left-going electrons ($`V_0<0`$) the torque is positive, so it acts to align the moments of the two magnetic layers, see section II.
Let us now consider the case of weak $`sd`$ exchange coupling, i.e., when the scattering coefficients depend only weakly on spin. We continue to assume that no spin-flip scattering occurs in the ferromagnetic layers. We define $`g_a`$ and $`g_b`$ as the average conductance per spin of the two layers (in unit of $`e^2/h`$). Then, the conductance of F<sub>a</sub> alone is $`g_a+\delta g_a`$ and $`g_a\delta g_a`$ for respectively the majority and minority spins, which defines the spin scattering asymmetry $`\delta g_a`$. In that case, we get to lowest nontrivial order in $`\delta g_a`$ and $`\delta g_b`$:
$`g`$ $`=`$ $`{\displaystyle \frac{2e^2}{h}}({\displaystyle \frac{g_\mathrm{N}g_ag_b}{g_ag_b+g_\mathrm{N}(g_a+g_b2\frac{g_ag_b}{N})}}`$ (54)
$`+`$ $`{\displaystyle \frac{g_\mathrm{N}^2\delta g_a\delta g_b\mathrm{cos}\theta }{(\frac{g_ag_b}{N}+g_\mathrm{N}(g_a+g_b2\frac{g_ag_b}{N}))^3}})`$ (55)
and
$$\frac{\tau _b^x}{V_0}=\frac{e}{2\pi }\frac{g_\mathrm{N}^2\delta g_a\delta g_b^2\mathrm{sin}\theta }{2(1\frac{g_b}{N})(g_ag_b+g_\mathrm{N}(g_a+g_b2\frac{g_ag_b}{N}))^2}.$$
(56)
This last formula shows that:
(i) The torque is not symmetric with respect to interchanging the layer F<sub>a</sub> and F<sub>b</sub>, in contrast to the conductance. If one changes $`\delta g_a`$ to $`\delta g_a`$, the sign of the torque is reversed. However, $`\tau _b^x/V_0\delta g_b^2`$, so if one changes $`\delta g_b`$ to $`\delta g_b`$, the sign of the torque is unchanged. The sign of the torque on a ferromagnetic layer therefore depends on whether the other layer is a positive or negative polarizer, but not on the sign of filtering for the layer experiencing the torque. We have verified that this is true also in the general case. This point explains why the two layers can not be treated on an equal footing.
(ii) We see that $`g_\mathrm{N}`$ appears through its square. Indeed, in order for some spin to be deposited in the layer F<sub>b</sub>, some left going electrons have to be reflected by F<sub>b</sub> and exit the system from the right hand side. Therefore these electrons cross the normal layer at least twice and this leads to the factor $`g_\mathrm{N}^2`$. On the other hand the conductance is linear in $`g_\mathrm{N}`$. Therefore in order to maximize the torque deposited per current, one has to use the cleanest possible normal metal spacer. (This statement is true in this limit of weak filtering, but not in general, see section V.) Note that in the previous case (perfect spin-filtering) the torque is proportional to $`g_\mathrm{N}`$ instead of the expected $`g_\mathrm{N}^2`$. Indeed, in that case, once the electron has been reflected by the layer F<sub>b</sub>, it cannot go through F<sub>b</sub> which works as a perfect wall for it. Therefore current conservation implies that it goes out of the system through the right. For $`g_\mathrm{N}N`$, the torque is actually proportional to $`g_\mathrm{N}^2`$ for arbitrary spin asymmetry (except perfect filtering), and one gets:
$$\frac{\tau _b^x}{V_0}g_\mathrm{N}^2\mathrm{sin}\theta ,g_\mathrm{N}N,$$
(57)
the factor of proportionality being a complicated function of the transmission probabilities of the layers.
## V Application to current-driven switching of magnetic domains
In this section, we consider the general solution Eq.(43) for the spin-transfer torque. We first address strongly polarizing systems and then calculate torques for scattering parameters more appropriate for the transition metal trilayers that can be studied experimentally. As the trilayer devices are primarily current-driven, we calculate the torque per unit of current $`I`$,
$$\frac{\tau _b^x}{I}=\frac{1}{g}\frac{\tau _x}{V_0}.$$
The torque is measured in units of $`\frac{h}{2\pi e}`$.
Eq. (52) of the previous section gives the torque per unit current for the case that both layers F<sub>a</sub> and F<sub>b</sub> are perfect polarizers. The main feature of this system is that the $`\theta `$ dependence of the torque is not of a simple $`\mathrm{sin}\theta `$ form, and that the torque per unit current diverges at $`\theta =0`$. In Fig. 5, we look at what happens when one of the layers (F<sub>b</sub>) is a nearly perfect polarizer while the other one is not. Although the divergence at $`\theta =\pi `$ is regularized, $`\tau _b^x/I`$ remains sharply peaked near $`\theta =\pi `$. This is relevant for the critical current needed to switch the magnetization of F<sub>b</sub> from $`\theta =\pi `$ to $`\theta =0`$. Recall that the switching of the domains follows from a competition between the spin-transfer torque on the one hand and restoring forces from local fields, anisotropy, exchange coupling etc. (The competition between these forces has been considered phenomenologically in Ref. using a phenomenological Landau-Lifschitz-Gilbert equation.) The torques for $`\theta `$ close to $`0`$ and $`\pi `$ determine the critical currents to overturn a metastable parallel (antiparallel) alignment of the moment in F<sub>a</sub> and F<sub>b</sub>. Hence the critical current should be different at $`\theta =0`$ and $`\theta =\pi `$.
In Fig. 6, we consider the same system as in Fig. 5 (one perfectly polarizing F layer, one partially polarizing layer), but as a function of the conductance of the normal layer $`g_\mathrm{N}`$ for angles $`\theta `$ close to $`0`$. We find that switching the two layers has a drastic effect on the torque, even at a qualitative level. Interestingly, in the case where F<sub>a</sub> is the nearly perfect layer (dashed line in Fig. 6), a maximum of the torque is found for $`\frac{g_\mathrm{N}}{N}1`$, i.e., in that case, a dirty metal spacer would give a higher torque (per unit of current) than a clean one.
At this stage, it is interesting to compare our theory to that of Ref . In this work the WKB approximation was used, and the electrons at the FN interfaces are either totally transmitted or reflected. For non-perfect polarizers, only a fraction of the channels act as perfect filters while the others perfectly transmit both the minority and majority spins. However, this situation is different from having non perfect transmission probabilities $`T_{}`$, $`T_{}`$ per channel. In particular, having $`T_{}N`$ channels that do not filter and $`(1T_{})N`$ perfect filters is not equivalent to $`N`$ channels that all partly transmit the minority spins with probability $`T_{}`$. This situation is illustrated in Fig 7. The latter scenario is supported by ab-initio calculations. Moreover, for a disordered normal-metal spacer, multiple scattering from impurities mixes all channels and the notion of two type of channels become superfluous. In that case, the torque is described by Eq.(43) in all cases. The torque found in the second case can be significantly smaller than under the assumption of Ref .
We can also compare our model to the work of Berger. While the theories of Berger and Slonczewsi have much in common, Berger does invoke inelastic spin-flip scattering in a way that Slonczewski does not. (Slonczewski’s theory utilizes only spin-filtering, without spin-flip scattering.) This effect can in principle be treated in our model, by including the off-diagonal spin-flip reflection and transmission amplitudes that we have thus far neglected. We shall comment on some of the implications in the conclusion. We suspect that the differing treatments of this aspect of the physics may explain why Slonczewski and Berger predict slightly different forms for the current-induced torques.
In our theory, the scattering matrices of the ferromagnetic layers appear as free input parameters. However, it is in principle possible to calculate them from first principle calculations for specific materials. Such an approach has been taken in Ref. and the results can be used to give some estimates of torques that can be expected in realistic systems. In Fig. 8, we compare the Co-Cu-Co system considered in the experiment of Ref. with the Fe-Cr-Fe system. In the latter, the minority spins have a larger transmission probability than the majority ones, explaining the opposite sign of the torque.
## VI Conclusion
We have developed a theory for the spin-transfer-induced torques on the magnetic moments of a ferromagnet-normal-ferromagnet FNF trilayer system caused by a flowing current. Our theory deals with the effects of multiple scattering between the layers using the scattering matrices of the ferromagnet-normal metal interfaces as input parameters. We consider both the cases of a diffusive and ballistic normal metal spacer. Remarkably, in the diffusive case, the high-dimensional scattering matrices of the FN interfaces only appear through the reduced $`4\times 4`$ tensor products of Eq.(39) which greatly reduces the number of degrees of freedom of the theory (see also Ref. ). This reduction of the number of degrees of freedom allows us to make qualitative predictions about the role of the interface transparency, normal metal resistance etc., without detailed knowledge of the microscopic details of the system. However, for quantitative predictions, inclusion of the microscopic parameters in our theory, e.g. from ab initio calculations is still needed.
Having a complete theoretical description of the current-induced switching of magnetic domains in FN multilayers as a final goal, the theory here can be regarded as being an intermediate step. On the one hand, microscopic input is needed for the scattering matrices of the FN interfaces, as explained above. On the other hand, the output of our theory, the current-induced torques, needs to be combined with restoring (hysteretic) forces in a more phenomenological theory that describes the dynamics of the magnetic moments. Such a theory involves anisotropy forces and information about the mechanism by which the torque is exerted (spin wave excitation, local exchange field) – issues which are still subject of debate..
In this paper, we have focused on the effects of “spin filtering” as the mechanism for current-induced torque, i.e., the difference in the transmission and reflection probabilities for electrons with spins parallel and antiparallel to the moments of the ferromagnetic layers (the diagonal terms in the matrices for the reflection and transmission amplitudes, Eq.(19).) A different source of spin-dependent scattering, which we have not considered in detail, but which is included in our formalism, is that of spin-flip scattering – the off-diagonal terms in Eq.(19). Its effect can be twofold. In the normal spacer, it would decrease the effective polarization, and therefore the torque. However, in the ferromagnet, the rate of spin-flip scattering might be asymmetric with respect to minority and majority spins, and therefore spin-flip scattering may also be an additional source of torque. As the number of degrees of freedom involved is much larger than for spin filtering only, a realistic model for the scattering matrices in the ferromagnets would be a necessary starting point for a theory that would include the effect of spin-flip scattering. We leave such a theory for future work.
###### Acknowledgements.
We thank A. Brataas and G. Bauer for drawing our attention to Ref. .
## APPENDIX
### A Derivation of Eq.(43)
In this appendix, we describe the calculation of Eq.(43) step by step.
First, we substitute the expression (30) for $`\mathrm{\Omega }`$ into (34), and then formally expand the resulting equation in powers of the reflection matrices $`r_a`$, $`r_b`$, $`r_n`$ and $`r_n^{}`$. Using the polar decomposition Eq.(37) for the reflection and transmission matrices $`r_n,t_n,t_n^{}`$ and $`r_n^{}`$ of the normal layer, we get a sum of many terms, each of which is of a form where contributions from N are alternated with those of F<sub>a</sub> and F<sub>b</sub>. Writing spin indices explicitly (summation over repeated indices is implied), we can write those terms as,
$$\mathrm{Tr}_N\stackrel{}{\mathrm{\Sigma }}_{ij}\left(A_{jk}\alpha B_{kl}\beta \mathrm{}\eta C_{sm}\right)\left(F_{nm}^{}\omega ^{}\mathrm{}\delta ^{}E_{pn}^{}\gamma ^{}D_{ip}^{}\right),$$
(58)
where $`A`$, $`B`$, $`C`$, $`D`$, $`E`$, $`F`$ $`\{r_a,t_a,r_b,r_b^{}\mathrm{}\}`$ refer to the layer F<sub>a</sub> and F<sub>b</sub> while $`\alpha `$, $`\beta `$, $`\gamma `$, $`\delta `$, $`\eta `$, $`\omega \{ui\sqrt{T}v,u\sqrt{1T}u^{},\mathrm{}\}`$ refer to the normal layer.
We are now ready to do the average of Eq.(58) over the matrices $`u,u^{},v`$ and $`v^{}`$ using the diagrammatic technique of Ref. (In leading order in $`N`$ these integrals reduce to the application of wick theorem.) Doing so, each of the $`\alpha ,\beta ,\mathrm{}`$ has to be put in correspondence with one of the $`\gamma ^{},\delta ^{}`$, etc. To leading order in $`N`$, only the ladder diagram survives, in which $`\alpha =\gamma `$, $`\beta =\delta `$, $`\eta =\omega `$,… and hence, $`A=D`$, $`B=E`$, $`C=F`$,… Thus, after averaging, we get terms like:
$`\mathrm{Tr}_N\left[\stackrel{}{\mathrm{\Sigma }}_{ij}\right]{\displaystyle \frac{1}{N}}\mathrm{Tr}_N\left[A_{jk}A_{ip}^{}\right]a{\displaystyle \frac{1}{N}}\mathrm{Tr}_N\left[B_{kl}B_{pn}^{}\right]b`$ (59)
$`\mathrm{}c{\displaystyle \frac{1}{N}}\mathrm{Tr}_N\left[C_{sm}C_{nm}^{}\right],`$ (60)
where $`a,b,c,\mathrm{}`$ stands for either $`\frac{1}{N}\mathrm{Tr}T`$ or $`\frac{1}{N}\mathrm{Tr}(1T)`$.
To leading order in $`N`$, the average over $`T`$ can now be done by simply replacing $`a,b,c,\mathrm{}`$ by their average value $`g_\mathrm{N}/N`$ or $`1g_\mathrm{N}/N`$ where $`g_\mathrm{N}`$ is the average conductance (per spin) of the normal layer, in unit of $`e^2/h`$.
Finally, denoting $`\lambda =(i,j)`$ and $`\mu =(k,p)`$, let us now introduce $`4\times 4`$ matrices $`\widehat{A}`$, $`\widehat{B}`$, $`\widehat{C}`$,… that are defined as:
$$\widehat{A}_{\lambda \mu }=\frac{1}{N}\mathrm{Tr}_N\left[A_{ik}A_{jp}^{}\right],$$
(61)
and $`\widehat{\stackrel{}{\mathrm{\Sigma }}}`$ is defined as
$$\widehat{\stackrel{}{\mathrm{\Sigma }}}_{\lambda \mu }=\delta _{kp}\mathrm{Tr}_N\stackrel{}{\mathrm{\Sigma }}_{ji}.$$
(62)
In term of these new matrices, eq.(59) now reads as a simple matrix product:
$$\mathrm{Tr}_4\widehat{\stackrel{}{\mathrm{\Sigma }}}\widehat{A}\widehat{\alpha }\widehat{B}\widehat{\beta }\mathrm{}\widehat{\eta }\widehat{C},$$
(63)
with $`\widehat{\alpha }`$, $`\widehat{\beta }`$, $`\widehat{\eta }`$,…$`\{g_\mathrm{N}/N,1g_\mathrm{N}/N\}`$. Equation (63) is formally equal to the expansion of $`\mathrm{\Omega }`$ (see Eq.(58)) except that we are now dealing with “hat” matrices. Therefore, we can now resum all the terms of the expansion and get Eq. (43).
### B Ballistic normal layer: a pedestrian approach
If N is very clean, and the interfaces are very flat, it is reasonable to assume that the electrons propagate ballistically inside the normal layer. The different modes will not be mixed in that case, and the electron wavefunction only picks up a phase factor $`e^{ik_iL}`$ where $`L`$ is the width of N and $`k_i`$ the momentum of channel $`i`$. For a sufficiently thick normal layer (i.e. $`L\lambda _F`$), small fluctuations of $`k_i`$ lead to an arbitrary change in the phase factor, and it is justified to consider $`e^{ik_iL}`$ as a random phase and to average over it. This is different from the case of a disordered metal spacer, where the average involves unitary matrices $`u`$, $`u^{}`$… that mix the channels, cf Eq.(37). In the case where $`r_a`$, $`r_a`$,… are proportional to the identity matrix (i.e. the reflection amplitudes do not depend on the channel), the ballistic model reduces to the disordered model of Eq. (43) for $`g_\mathrm{N}=N`$.
The reflection matrices of N being zero, the matrix $`\mathrm{\Omega }`$ reads:
$$\mathrm{\Omega }=e^{ik_iL}\frac{1}{1e^{2ik_iL}r_ar_b^{}}t_a^{}.$$
(64)
Neglecting spin-flip scattering, denoting $`z=e^{2ik_iL}`$, and choosing $`r_{a11}=r_a`$, $`r_{a22}=r_a`$,… where $`r_a`$, $`r_a`$,… are diagonal matrices, one gets after some algebra:
$$\frac{\tau _b^x}{V_0}(z)=\frac{e\nu }{4\pi }\mathrm{Tr}\mathrm{Re}\frac{A(z)}{z|D(z)|^2}\mathrm{sin}\theta ,$$
(65)
where $`A(z)`$ and $`D(z)`$ stand for:
$`A(z)`$ $`=`$ $`\left(1t_b^{}t_b^{}r_b^{}r_b^{}\right)`$ (66)
$`(`$ $`|t_a^{}|^2(1zr_b^{}r_a)(zr_b^{}r_a^{})`$ (67)
$``$ $`|t_a^{}|^2(1zr_br_a)(zr_b^{}r_a^{})),`$ (68)
$`D(z)`$ $`=`$ $`1z[\mathrm{cos}^2{\displaystyle \frac{\theta }{2}}(r_ar_b^{}+r_ar_b^{})`$ (69)
$`+`$ $`\mathrm{sin}^2{\displaystyle \frac{\theta }{2}}(r_ar_b^{}+r_ar_b^{})]`$ (70)
$`+`$ $`z^2r_ar_ar_b^{}r_b^{}.`$ (71)
A similar formula can be written for the conductance $`g(z)`$:
$$g=\frac{e^2}{h}\mathrm{Tr}\frac{B(z)}{z|D(z)|^2},$$
(72)
with:
$`B(z)`$ $`=`$ $`|t_a^{}|^2|t_b^{}|^2\mathrm{cos}^2{\displaystyle \frac{\theta }{2}}(1zr_ar_b^{})(zr_a^{}r_b^{})`$ (73)
$`+`$ $`|t_a^{}|^2|t_b^{}|^2\mathrm{sin}^2{\displaystyle \frac{\theta }{2}}(1zr_ar_b^{})(zr_a^{}r_b^{})`$ (74)
$`+`$ $`|t_a^{}|^2|t_b^{}|^2\mathrm{sin}^2{\displaystyle \frac{\theta }{2}}(1zr_ar_b^{})(zr_a^{}r_b^{})`$ (75)
$`+`$ $`|t_a^{}|^2|t_b^{}|^2\mathrm{cos}^2{\displaystyle \frac{\theta }{2}}(1zr_ar_b^{})(zr_a^{}r_b^{}).`$ (76)
Taking the average over the phases now amounts to coutour integration for $`z`$:
$$f=\frac{1}{2\pi i}\frac{dz}{z}f(z),$$
(77)
where the integration is done along the unit circle. The result is then given by the sum of the poles that are inside the unit circle. The two poles of $`D(z)`$ are outside the unit circle, while the two poles $`z_1`$ and $`z_2`$ of $`z^2D\left(\frac{1}{z}\right)`$ are inside the circle. They are given by:
$`z_i`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{cos}^2{\displaystyle \frac{\theta }{2}}(r_ar_b^{}+r_ar_b^{})+\mathrm{sin}^2{\displaystyle \frac{\theta }{2}}(r_ar_b^{}+r_ar_b^{})`$ (78)
$`+`$ $`{\displaystyle \frac{1}{2}}(1)^i[\mathrm{cos}^4{\displaystyle \frac{\theta }{2}}(r_ar_b^{}r_ar_b^{})^2`$ (79)
$`+`$ $`2\mathrm{cos}^2{\displaystyle \frac{\theta }{2}}\mathrm{sin}^2{\displaystyle \frac{\theta }{2}}(r_b^{}r_b^{}(r_ar_a)^2+r_ar_a(r_b^2+r_a^2))`$ (80)
$`+`$ $`\mathrm{sin}^4{\displaystyle \frac{\theta }{2}}(r_ar_b^{}r_ar_b^{})^2]^{\frac{1}{2}}.`$ (81)
The averaged torque and conductance are then simply given by
$$\frac{\tau _b^x}{V_0}=\frac{e\nu }{4\pi }\frac{\mathrm{sin}\theta }{z_1z_2}\mathrm{Tr}\left(\frac{A(z_1)}{D(z_1)}\frac{A(z_2)}{D(z_2)}\right)$$
(82)
and
$$g=\frac{e^2}{h}\frac{1}{z_1z_2}\mathrm{Tr}\left(\frac{B(z_1)}{D(z_1)}\frac{B(z_2)}{D(z_2)}\right).$$
(83)
In the case where all the channels are not identical, these results can be generalized by introducing a $`k`$ dependence of the different transmission/reflection amplitudes.
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# Abstract
## Abstract
This paper discusses the implementation of diffeomorphism invariance in purely Hamiltonian formulations of General Relativity. We observe that, if a constrained Hamiltonian formulation derives from a manifestly covariant Lagrangian, the diffeomorphism invariance of the Lagrangian results in the following properties of the constrained Hamiltonian theory: the diffeomorphisms are generated by constraints on the phase space so that a) The algebra of the generators reflects the algebra of the diffeomorphism group. b) The Poisson brackets of the basic fields with the generators reflects the space-time transformation properties of these basic fields. This suggests that in a purely Hamiltonian approach the requirement of diffeomorphism invariance should be interpreted to include b) and not just a) as one might naively suppose. Giving up b) amounts to giving up objective histories, even at the classical level. This observation has implications for Loop Quantum Gravity which are spelled out in a companion paper. We also describe an analogy between canonical gravity and Relativistic particle dynamics to illustrate our main point.
## 1 Introduction
The diffeomorphism invariance of General Relativity presents both conceptual and technical problems for quantisation. At the conceptual level, it leads to deep questions about the nature of time, observables and the interpretation of quantum theory. At the technical level, diffeomorphism invariance leads to constraints on the classical phase space , which in a quantum theory, must be imposed on physical states. Solving these constraints has occupied much of the effort in the canonical approach to quantum gravity. Several constrained Hamiltonian formulations (CHFs) of General Relativity exist today, each with its own following. It remains to be seen which of these formulations will be the most advantageous in the approach to quantum General Relativity.
This paper seeks to clarify the meaning of diffeomorphism invariance in a classical, constrained Hamiltonian Theory. Given a constrained theory, how does one test for diffeomorphism invariance? The answer to this question involves a subtlety, on which we focus in this paper. There is a substantial literature on the constrained Hamiltonian formulation of diffeomorphism invariant theories. The point we wish to emphasise here is perhaps implicit in these earlier works, but we wish to make it explicit in order to use it in . Our strategy in addressing this question will be to start with CHF’s which we know are diffeomorphism invariant: those that are derived by a Legendre transformation from a manifestly covariant Lagrangian. We will then notice that the resulting constrained Hamiltonian formulation satisfies certain conditions as a consequence of the diffeomorphism invariance of the Lagrangian starting point. We will explicitly spell out these conditions and use these as a criterion for testing for diffeomorphism invariance even when a Lagrangian starting point is not available. For example many currently popular CHFs of General Relativity are derived by making canonical transformations on the phase space; they are entirely Hamiltonian in spirit and are often presented and discussed without any Lagrangian starting point. One would like to discuss the diffeomorphism invariance of such formulations in a purely Hamiltonian framework. The purpose of this paper is to clarify how this can be done.
The paper is organised as follows: In section II, we recapitulate some known results about the gauge symmetries of Lagrangian systems and show how these symmetries manifest themselves in a Hamiltonian framework. In section III we illustrate these general results using familiar examples like the ADM formalism, gravity in 2+1 dimensions and Ashtekar’s extended phase space construction (EPS). In section IV, we distinguish between strong and weak diffeomorphism invariance of a CHF and bring out an analogy with a much simpler situation: relativistic particle dynamics. Section V is a concluding discussion.
## 2 Symmetries of Singular Lagrangian systems
Consider a dynamical system with configuration manifold $`𝒬`$ on which local co-ordinates are $`q^r,r=1..n`$. The tangent bundle over $`𝒬`$ is $`T𝒬`$ and the Lagrangian $`L(q,\dot{q})`$ is a real valued function on $`T𝒬`$. The Lagrangian $`L`$ defines a map from $`T𝒬`$ to the cotangent bundle $`T^{}𝒬`$ defined locally by $`p_r=\frac{L}{\dot{q}^r}`$. In the cases of interest in this paper, the Lagrangian $`L`$ is singular, i.e, the Legendre map $`\mathrm{\Phi }:T𝒬T^{}𝒬`$ is not onto. Its image $`\mathrm{\Sigma }`$ is a proper subset of $`T^{}𝒬`$: $`\mathrm{\Phi }(T𝒬)=\mathrm{\Sigma }T^{}𝒬`$ and there are constraints on the phase space. Such situations are dealt with in Dirac’s theory of constrained systems . One iteratively demands preservation of the constraints and this leads, in general, to more constraints. The algorithm terminates when no new constraints emerge. The total set of constraints are divided into first and second class and we suppose that the second class constraints are eliminated by passage to the Dirac bracket. An elegant way to do this is to use the Bergmann-Komar starring procedure. One simply replaces all phase space functions by their starred counterparts. After the Dirac constraint analysis ends, one has a constrained Hamiltonian formulation which has the following ingredients: i) the basic variables (or fields, in a field theory) are $`(q^r,p_r)`$ which span the phase space obeying commutation relations <sup>1</sup><sup>1</sup>1 these may not be canonical if some second class constraints have been eliminated.. ii) a physical interpretation for $`q^r`$ and $`p_r`$ that derives from their definitions as functions of $`q`$ and $`\dot{q}`$. iii) a set of constraints which emerge from the constraint analysis. iv) A Hamiltonian function on the phase space, which generates the dynamics and preserves the constraints. The Hamiltonian is arbitrary to the extent of a primary first class constraint<sup>2</sup><sup>2</sup>2 We do not follow Dirac’s suggestion of “extending” the Hamiltonian by adding arbitrary combinations of the secondary first class constraints to it, since we wish to stick with the Lagrangian starting point. This equally means that we do not “extend” the symmetry vector field in a similar manner..
Let us recapitulate a few known results about the continuous symmetries of singular Lagrangian systems. Let $`S^r(q,\dot{q},t)`$ be a symmetry transformation. By this we mean that the change $`\delta _SL`$ in the Lagrangian under the changes $`\delta _Sq^r=ϵS^r(q,\dot{q},t)`$, $`\delta _S\dot{q}^r=ϵ\dot{S}^r`$ in $`(q^r,\dot{q}^r)`$ is given by a total divergence:
$$\delta _SL=ϵ\frac{dF(q,\dot{q},t)}{dt}.$$
(1)
(Note that in (1) we do not use the Euler-Lagrange equations, the accelerations are unrestricted.) From (1) it follows that on solutions to the equations of motion, the quantity $`G_{}(q,\dot{q},t):=\frac{L}{\dot{q}^r}S^rF`$ is conserved as a result of Nöther’s theorem. (1) also implies that $`G_{}(q,\dot{q},t)`$ is projectable under the Legendre map and therefore can be expressed as the pull back of a function on $`\mathrm{\Sigma }`$ : $`G_{}=\mathrm{\Phi }^{}G`$, . In general, the symmetry vector field $`X_S:=S^r\frac{}{q^r}+\dot{S}^r\frac{}{\dot{q}^r}`$ (which is defined on $`T𝒬`$ by using the equations of motion, or more briefly, the dynamics $`\mathrm{\Delta }`$) is not $`\mathrm{\Phi }`$ projectable<sup>3</sup><sup>3</sup>3 The non-projectability of a symmetry transformation has also been recently remarked in . While these papers too address the question of interplay between gauge symmetries and diffeomorphism invariance, their motivation is different from ours: they seek to find combinations of diffeos and gauge which are $`\mathrm{\Phi }`$ projectable. Our interest here is in pure diffeo’s , which in general are not projectable.. The vertical part of $`X_S`$ projects down to zero, and the horizontal part can be expressed in the form
$`\delta _Sq^r`$ $`=`$ $`ϵ({\displaystyle \frac{G}{p_r}}+u^\rho {\displaystyle \frac{\varphi _\rho }{p_r}})`$
$`\delta _Sp_r`$ $`=`$ $`ϵ({\displaystyle \frac{G}{q^r}}+u^\rho {\displaystyle \frac{\varphi _\rho }{q^r}})`$ (2)
where $`\varphi _\rho (q,p)`$ are the primary constraints. The functions $`u^\rho `$ are functions on $`T𝒬`$, which are not in general projectable under $`\mathrm{\Phi }`$. The non-projectability of $`X_S`$ has been isolated in the functions $`u^\rho `$, which depend not only on the phase space variables $`(q^r,p_r)`$, but also the “unsolved velocities” $`v^\rho `$. As the Dirac analysis proceeds, the dynamics, and with it the symmetry vector field (which depends on the dynamics), gets more sharply determined . From the basic identity (1) it follows that the symmetry is “compatible” with the dynamics throughout the constraint analysis: if the dynamics preserves constraints, so does the symmetry. The symmetry generator of $`X_S`$ is $`𝒢_S=G+u^\rho \varphi _\rho `$, which Kamimura refers to as a Generalised Canonical Quantity because the $`u(q,p,v)`$ are not strictly phase space functions. (They depend on the unsolved velocities $`v`$). Symmetries of the Lagrangian translate into the following properties of the constrained Hamiltonian formulation, which hold on shell, (i.e, modulo the equations of motion):
a) The Lie algebra of the symmetry group is reflected in the bracket relations of the symmetry generators $`𝒢`$.
b) The basic variables $`q^r`$ and $`p_r`$ are functions on $`T𝒬`$ and transform in a definite manner under the symmetry tranformation $`S`$. This transformation property is reflected in the bracket relations between these basic variables and $`𝒢_S`$.
$`\delta _Sq^r`$ $`=`$ $`ϵ\{q^r,𝒢_S\}`$
$`\delta _Sp_r`$ $`=`$ $`ϵ\{p_r,𝒢_S\},`$ (3)
where $`\{,\}`$ refers to the Dirac bracket resulting from elimination of second class constraints (if any).
In the rest of this paper we apply these general considerations to the case of interest to us. We consider constrained Hamiltonian formulations of General Relativity and the symmetry of interest is diffeomorphism invariance. In this case, as is well known, the generators $`𝒢_S`$ are a linear combination of constraints. The criteria listed above can be used to test for invariance even in a purely Hamiltonian framework i.e, even when a Lagrangian is absent. Below, we will slightly weaken them to allow for the possibility that they are satisfied modulo gauge<sup>4</sup><sup>4</sup>4 In this paper, we reserve the word ‘gauge’ to mean “internal” gauge. Diffeomorphisms will not be referred to as ‘gauge’ transformations. transformations<sup>5</sup><sup>5</sup>5 In this case the diffeomorphism group is twisted with the internal gauge group. This point is discussed further in the concluding section.
## 3 Diffeomorphism Invariant Formulations
We now examine some constrained Hamiltonian formulations of diffeomorphism invariant theories to see that they do indeed satisfy the conditions listed above. All of these formulations are derived from diffeomorphism invariant Lagrangians. Let $`(,g_{\mu \nu }),\mu =0,1,2,3`$ be a space-time manifold, topologically $`𝒮\times IR`$. To simplify matters, we will assume that $`𝒮`$ has no boundary so that we don’t need to keep track of spatial boundary terms. We are also interested only in infinitesimal diffeomorphisms and deal entirely with the Lie Algebra rather than the Lie Group of diffeomorphisms. These infinitesimal diffeomorphisms are generated by constraints. The constraint algebra ensures that a) is satisfied. The property a) is discussed extensively in the canonical gravity literature as “path independence” of evolution and we do not dwell on it any further. We wish to concentrate on the condition b), which is perhaps implicitly assumed to be true in the above references. From the general discussion of the last section, we expect that the CHF’s will satisfy the conditions a) and b) above. We write down condition b) explicitly in a few concrete cases and note that it is satisfied.
ADM formulation: The ADM formulation consists of the following ingredients: the basic variables are $`(q_{ab},\stackrel{~}{\pi }^{ab})`$, which are canonically conjugate. $`q_{ab}`$ is the pullback of the space-time metric to a spatial slice $`𝒮`$ and $`\stackrel{~}{\pi }^{ab}`$ is its conjugate momentum. The basic variables ($`q_{ab},\stackrel{~}{\pi }^{ab}`$) are subject to the Hamiltonian constraint
$$\stackrel{~}{}=\frac{1}{\sqrt{q}}(\stackrel{~}{\pi }^{ab}\stackrel{~}{\pi }_{ab}\frac{1}{2}\stackrel{~}{\pi }^2)\sqrt{q}{}_{}{}^{3}R0$$
and the spatial diffeomorphism constraint
$$\stackrel{~}{}^b=D_a\stackrel{~}{\pi }^{ab}0$$
where $`D`$ is the covariant derivative compatible with the three-metric $`q_{ab}`$.
The condition (b) holds in the ADM formalism, as one would expect from the general analysis of the last section. The basic variables of the theory are $`(q_{ab},\stackrel{~}{\pi }^{ab})`$ and they have definite space-time meaning: $`q_{ab}`$ is the pull-back of the space-time metric to a spatial slice $`𝒮`$. By using Hamilton’s equations of motion we see that $`\stackrel{~}{\pi }^{ab}`$ is algebraically related to the extrinsic curvature of $`𝒮`$. Since the basic fields have a clear space-time meaning, they have a definite transformation property under space-time diffeomorphisms. For example, under an infinitesimal diffeomorphism generated by a vector field $`\xi ^a`$ tangent to $`𝒮`$ ($`a=1,2,3`$ is a spatial index), we expect
$$\delta q_{ab}(\mathrm{space}\mathrm{time})=(D_a\xi _b+D_b\xi _a).$$
If $`\xi `$ is normal to $`𝒮`$, $`\xi ^\mu =N\widehat{n}^\mu `$ we expect
$$\delta q_{ab}(\mathrm{space}\mathrm{time})=_\xi q_{ab}=NK_{ab},$$
where $`K_{ab}`$ is the extrinsic curvature of $`𝒮`$.
One can also compute the change in the basic variables by taking their Poisson brackets with the diffeomorphism generator $`C(\xi )`$:
$`\delta _\xi q_{ab}(\mathrm{canonical})`$ $`=`$ $`\{q_{ab},C(\xi )\},`$
$`\delta _\xi \stackrel{~}{\pi }^{ab}(\mathrm{canonical})`$ $`=`$ $`\{\stackrel{~}{\pi }^{ab},C(\xi )\}.`$ (4)
The condition (b) is satisfied in the ADM formalism since as follows from Hamilton’s equations
$`\delta _\xi q_{ab}(\mathrm{space}\mathrm{time})`$ $`=`$ $`\delta _\xi q_{ab}(\mathrm{canonical})`$
$`\delta _\xi \stackrel{~}{\pi }^{ab}(\mathrm{space}\mathrm{time})`$ $`=`$ $`\delta _\xi \stackrel{~}{\pi }^{ab}(\mathrm{canonical}).`$
2+1 Palatini gravity: The next example we consider is gravity in 2+1 dimensions in its Palatini formulation. The basic fields are $`e_\mu ^I`$ and $`A_\mu ^{IJ}`$; $`\mu =0,1,2`$ is a tangent space index and $`I=0,1,2`$ is an internal Minkowski index. $`e_\mu ^I`$ is a triad and $`A_\mu ^{IJ}`$ an $`SO(2,1)`$ connection. The action is given by
$$I=\frac{1}{2}e_IF^I,$$
where $`F=dA+AA`$ in the notation of differential forms. A standard constraint analysis leads to the following Hamiltonian formulation: The basic variables are the canonically conjugate pair $`(\stackrel{~}{e}^I{}_{}{}^{a}:=\stackrel{~}{\eta }^{ab}e_b^I,A_a^I)`$, where $`a`$ is a spatial index. The constraints of the theory are
$`F^I`$ $`=`$ $`0`$
$`G^I`$ $`=`$ $`𝒟e^I=0,`$
where it is understood that these two-forms are pulled back to a spatial slice $`𝒮`$. Diffeomorphism are generated by combinations of constraints. If $`\xi ^\mu `$ is a vector field on $``$,
$$C(\xi )=_𝒮(\xi ^\mu e_\mu ^IF_I+\xi ^\mu A_\mu ^IG_I)$$
generates a pure diffeomorphism on the basic variables. It is easily checked, using the $`ISO(2,1)`$ algebra satisfied by the constraints that the condition (b) above is satisfied on shell (using the equations of motion).
Extended phase space construction (EPS): As a last example, we consider the extended phase space of Ashtekar. This CHF was originally arrived at by Ashtekar by extending the ADM phase space to incorporate triads. This example is instructive because it can also be derived from a manifestly covariant Lagrangian by fixing the “time gauge”. This example will illustrate how internal gauge fixing interacts with diffeomorphism invariance. As we will see, because of the gauge fixing a) and b) are not satisfied as they stand but they are satisfied modulo $`SO(3)`$ gauge.
Let us start with the following action principle. The basic fields are $`e_\mu ^I,A_\mu ^{IJ}`$, where $`e_\mu ^I`$ is a tetrad field and $`A_\mu ^{IJ}`$ an $`SO(3,1)`$ connection field. The action is
$$I=\frac{1}{2}e^Ie^JF^{KL}ϵ_{IJKL},$$
(5)
where we use differential form notation and $`F=dA+AA`$. A straight forward Legendre transformation results in the following CHF. The basic conjugate variables are $`(A_a,\stackrel{~}{\alpha }^a)`$ where
$$\stackrel{~}{\alpha }_{IJ}^a=\stackrel{~}{\eta }^{abc}e_{bI}e_{cJ}.$$
(6)
These variables are subject to the constraints
$$G_{IJ}=𝒟_a\stackrel{~}{\alpha }_{IJ}^a0$$
(7)
$$V_a=Tr\stackrel{~}{\alpha }^bF_{ab}0$$
(8)
$$S=Tr\stackrel{~}{\alpha }^a\stackrel{~}{\alpha }^bF_{ab}0$$
(9)
$$\varphi ^{ab}:=ϵ^{IJKL}\stackrel{~}{\alpha }_{IJ}^a\stackrel{~}{\alpha }_{KL}^b0$$
(10)
$$\chi ^{ab}:=ϵ^{IJKL}\stackrel{~}{\alpha }_I^{cM}\stackrel{~}{\alpha }_{MJ}^{(a}(𝒟_c\stackrel{~}{\alpha }^{b)})_{KL}0.$$
(11)
Of these the last two (10, 11) are second class. Let us suppose these second class constraints to be formally eliminated by passing to the Dirac bracket. No gauge fixing has been done so far and it follows from the general theory summarised in the last section that the Hamiltonian formulation above satisfies a) as well as b).
(10) implies that $`\stackrel{~}{\alpha }_{IJ}`$ is of the form $`\stackrel{~}{E}^a{}_{[I}{}^{}n_{J]}^{}`$ for some internal vector $`n_J`$. Let us now impose the “time” gauge, i.e., pick $`n_I`$ to have the standard form $`\stackrel{}{n}^I=(1,0,0,0)`$. This corresponds to choosing $`e^0`$ to be normal to the spatial slice $`𝒮`$. One is, of course, at liberty to make this gauge choice. In order to enforce this gauge choice, we need to impose a constraint
$$\chi ^I=n^I\stackrel{}{n}^I0.$$
(12)
This constraint breaks the $`SO(3,1)`$ gauge generated by the Gauss law constraint (7) down to $`SO(3)`$. The “Boost part”
$$B_I=G_{IJ}\stackrel{}{n}^J$$
(13)
of (7) does not commute with (12) and in fact $`(B_I,\chi ^I)`$ form a second class set. If one eliminates this second class set one arrives at EPS. Writing $`i,j`$ instead of $`I,J`$ for indices orthogonal to $`\stackrel{}{n}^I`$, we find that basic variables of EPS are $`(\stackrel{~}{E}^a{}_{i}{}^{},K_a{}_{}{}^{i})`$ which are canonically conjugate and have the space-time interpretation of densitised triad and extrinsic curvature respectively. The constraints of the theory are:
$`ϵ_{ijk}K_a^j\stackrel{~}{E}^{ak}`$ $``$ $`0`$
$`D_a[\stackrel{~}{E}_k^aK_b^k\delta _b^a\stackrel{~}{E}_k^cK_c^k]`$ $``$ $`0`$
$`\sqrt{q}R+{\displaystyle \frac{2}{\sqrt{q}}}\stackrel{~}{E}_i^{[a}\stackrel{~}{E}_j^{b]}K_a^iK_b^j`$ $``$ $`0,`$
where $`D_a`$ is the covariant derivative associated with $`q_{ab}`$ and $`R`$, its scalar curvature.
Are conditions a) and b) satisfied in the gauge fixed theory? Diffeomorphisms that displace $`𝒮`$ normal to itself will in general, spoil the “time gauge”. In order to restore the “time gauge” (and this is the BK starring procedure of passing to Dirac brackets) one has to add some definite linear combination of $`B_I`$ to the diffeomorphism generator. As a result, (since the commutator of two boosts is a rotation) the diffeomorphism algebra closes only up to $`SO(3)`$ gauge rotations. In the same way, (b) is only satisfied up to $`SO(3)`$ gauge rotations. We describe this theory as satisfying a) and b) (mod $`SO(3)`$ gauge). The lesson to be learned from this example is that if one derives a Hamiltonian formulation from a Lagrangian and fixes gauge in the derivation, the resulting Hamiltonian formulation is diffeomorphism invariant(modulo gauge).
## 4 Strong and Weak Diffeomorphism Invariance
It is clear from these examples that diffeomorphism invariance in the Hamiltonian framework means more than getting the constraint algebra right. It is also necessary that under the action of the diffeomorphism generators, the basic variables must transform as expected from their space-time interpretation. We will refer to a CHF which satisfies the first condition (a) as weakly diffeomorphism invariant. A theory that also satisfies (b) is called strongly diffeomorphism invariant. It is clear that before we can test a CHF for diffeomorphism invariance, the space-time meaning of the basic variables has to be declared, since condition (b) explicitly needs this knowledge.
To better understand the meaning of Strong Diffeomorphism invariance, it is useful to consider a simpler but analogous situation: classical relativistic particle dynamics . Direct interactions between $`N`$ relativistic particles in Minkowski space can be described by mathematical models which are constrained Hamiltonian formulations. The models are defined as follows: the basic variables are $`(x_a{}_{}{}^{\mu },p_{a\mu }),(a=1..N,\mu =0,1,2,3)`$, where $`a`$ is particle index (for the duration of this section) and $`\mu `$ a Minkowski space-time index. One can define the system by imposing $`2N`$ second class constraints. The constraints are needed to reduce the phase space degrees of freedom from $`8N`$ to $`6N`$, which is the right number for $`N`$ particles. The symmetry of interest here is the Poincare group. We will say that a model is Poincare invariant if the following conditions hold:
a)There exist $`10`$ functions (one for each of the Poincare generators) on the phase space whose Dirac brackets reflect the Lie Algebra of the Poincare group.
b) The Dirac brackets between the basic variables $`(x^\mu {}_{a}{}^{},p_{a\mu })`$ and the Poincare generators reflect the space-time transformation properties of the basic variables.
As was first pointed out by Pryce , Poincare invariance means both a) and b) and not just a). Bakamjian and Thomas were able to construct interacting models, but at the cost of giving up condition b). In these models , particle world lines would depend on the Lorentz frame of the observer. (To clarify this point, it is not just the same world line viewed from different Lorentz frames, but different world lines.) This amounts to giving up the objectivity of world lines, or particle histories, which is unacceptable, since classically, particle world lines can be experimentally measured <sup>6</sup><sup>6</sup>6 There are models (see model 1 of this paper) in which these conditions are satisfied modulo reparametrisation gauge. This is quite acceptable since, it does not compromise the objectivity of particle World lines. All that happens is that the World line is reparametrised under a Poincare transformation..
To clarify the meaning of the conditions a) and b) above, we discuss two models which are in the literature . Both models describe two interacting particles and are defined by imposing $`4`$ constraints on the sixteen dimensional phase space spanned by $`(x_a{}_{}{}^{\mu },p_a{}_{\mu }{}^{})`$, the position and momentum four vectors of the particles.
model 1 The constraints that define this model are
$`K_1`$ $`=`$ $`p_1^2+m_1^2+V((x_1x_2)^2)`$ (14)
$`K_2`$ $`=`$ $`p_2^2+m_2^2+V((x_1x_2)^2)`$ (15)
$`\chi _1`$ $`=`$ $`P.(x_1x_2)`$ (16)
$`\chi _2`$ $`=`$ $`P.x_1\tau ,`$ (17)
where $`P=p_1+p_2`$ is the total momentum four-vector of the two particles, $`V((x_1x_2)^2)`$ is a potential function which depends only on the invariant interval $`(x_1x_2)^2`$ between the two particle position four vectors and $`\tau `$ is an evolution parameter, which plays the role of “time”. The constraints reduce the dimension of the phase space to 12 and do this in a Poincare invariant manner: both a) and b) above are satisfied. To see this, note that $`P_\mu `$ and $`M_{\mu \nu }=\mathrm{\Sigma }_{a=1,2}(x_a{}_{\mu }{}^{}p_{a}^{}{}_{\nu }{}^{}x_a{}_{\nu }{}^{}p_{a}^{}{}_{\mu }{}^{})`$, the ten generators of the Poincare group commute with $`K_1,K_2`$ and $`\chi _1`$ (with all but one of the constraints). It follows from this (and the definition of the Dirac bracket) that a) above is satisfied. It also follows that b) is satisfied modulo reparametrisation, since the Dirac bracket of the basic variables $`(x_a{}_{}{}^{\mu },p_a{}_{\mu }{}^{})`$ with the Poincare generators agrees with the Poisson bracket (which in turn agrees with the four-vector space-time transformation property of $`(x_a{}_{}{}^{\mu },p_a{}_{\mu }{}^{})`$) apart from a term representing the reparametrisation of the world line. Writing $`G`$ for any one of the ten generators of the Poincare group,
$$\{x_a{}_{}{}^{\mu },G\}^{}=\{x_a{}_{}{}^{\mu },G\}+\frac{dx_a^\mu }{d\tau }\delta _a\tau $$
(19)
for some $`\delta _a\tau `$.
model 2 The second model is defined by the constraints
$`K_1`$ $`=`$ $`p_1^2+m_1^2+V(x_1,x_2,p_1,p_2)`$ (20)
$`K_2`$ $`=`$ $`p_2^2+m_2^2+V(x_1,x_2,p_1,p_2)`$ (21)
$`\chi _1`$ $`=`$ $`(x_1x_2)^0`$ (22)
$`\chi _2`$ $`=`$ $`(x_1)^0\tau ,`$ (23)
$`K_1`$ and $`K_2`$ are required to commute with each other and with all the Poincare generators. From the definition of the Dirac bracket it follows that this model satisfies a). However, it does not satisfy b). It must therefore be rejected as a description of two relativistic particles.
The analogy between the models described above and canonical gravity is as follows: The symmetry group of interest in the first case is the Poincare group and in the second case the diffeomorphism group. The classical histories of the first system describe the world lines of N particles and in the second case a space–time. In both cases, the problem is one of realising a symmetry group in a purely Hamiltonian framework. Our main point here is that it is possible to get the symmetry algebra right but still violate the symmetry by giving up b). Model 2 is an example of this.
Returning to our problem in canonical gravity, a CHF which is only weakly diffeomorphism invariant suffers from the following feature: Given a space-time history (a solution of the field equations), one can slice it up in many ways in a 3+1 formalism. Conversely, given initial data and particular slicing one can evolve the initial data and produce a “history” by “stacking” the spatial slices in temporal order. In theories where b) is given up the “history” which is produced depends on the slicing. This means that the history has no objective reality. One would of course like measureable quantities to have an objective meaning (independent of slicing). One should therefore be aware of which fields in a theory are objectively real. For example, in the EPS, the fields $`q_{ab}`$ and $`K_{ab}`$ (which are $`SO(3)`$ gauge invariant) do have objective reality. But the basic fields in the formulation $`(\stackrel{~}{E}^a{}_{i}{}^{},K_a{}_{}{}^{i})`$ do not. They are only defined modulo $`SO(3)`$ gauge.
## 5 Conclusion
We have shown that for a constrained Hamiltonian formulation of gravity to be diffeomorphism invariant there must be diffeomorphism generators on the phase space so that a) the generators reflect the algebra of the diffeomorphism group in their brackets and b) the space-time interpretation of the basic fields is reflected in their brackets with the diffeomorphism generators. These conditions are automatically satisfied by CHFs which derive from a covariant Lagrangian. In the absence of a Lagrangian these conditions can be used to test for diffeomorphism invariance. Condition a) has been emphasised in the literature, but it appears that condition b) is usually left implicit. In this paper we point out what goes wrong if one gives up condtion b): one loses the objectivity of history.
Notice that a covariant Lagrangian automatically gives us space-time interpretations for all the phase space variables appearing in the Hamiltonian formulation. In a purely Hamiltonian approach one has to not only prescribe the basic variables, their brackets, and the constraints, but also give a space-time interpretation for the basic variables. Unless this is done, it is not possible to physically interpret the Hamiltonian system. If the PB of the diffeomorphism generator with a phase space variable does not reflect its space-time interpretation, one loses a space-time interpretation for that variable even at the classical level.
The Diffeomorphism invariance of the theory can only be decided after the space-time interpretation of the basic variables has been declared (since condition b) explicitly needs this knowledge). Indeed, unless the space-time interpretation of the basic variables is declared, the CHF is not even fully defined. A CHF may be diffeomorphism invariant with one interpretation and not invariant with another space-time interpretation of the basic variables. An example of this phenomenon is discussed in . Barbero’s Hamiltonian formulation of General Relativity is strongly diffeomorphism invariant with another space–time interpretation of the basic variables, but not with the space–time “gauge field interpretation” that one might prefer. In contrast, Ashtekar’s Hamiltonian formulation is SDI with both interpretations for the connection variable: that deriving from the canonical transformation as well as for the space–time gauge field interpretation.
In the EPS model, the Lie Algebra of the Diffeomorphism group is not a subalgebra of the constraint generators, but appears as a quotient. We arrived at conditions a) and b) by assuming that the symmetry group of interest was a subgroup of the total Lagrangian symmetry group. The known Lagrangian formulations of General Relativity all have the property that the diffeomorphism group is a subgroup. One can slightly relax this assumption and allow for “twisted products”, where the diffeomorphism group only appears as a quotient. The situation then is very similar to the EPS formulation, where the diffeomorphism Lie Algebra only closes modulo gauge.
One may object that one should not demand that the basic variables be objectively defined in space-time, since they are not “observables” in the Dirac sense. This objection is easily met: it is easy to construct “observables” from the basic variables by using a device explained in . Although $`q_{ab}`$ is not an “observable”, the distance between invariantly specified events is an “observable”. E.g, one can locate an event as the intersection of two particle world lines or (in the absence of matter) as a point where four scalars constructed from the gravitational field vanish. If a CHF is strongly diffeomorphism invariant in the sense of this paper, such “observables” do have an objective meaning. Otherwise, the answer predicted by the CHF could depend on slicing. A CHF which violates strong diffeomorphism invariance classically should be rejected as an unsuitable starting point for building a quantum theory.
Acknowledgement: It is a pleasure to thank Richard Epp, B.R. Iyer, Sukanya Sinha and Madhavan Varadarajan for extended discussions.
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# Flavour Stability of the Chiral Vacuum and Scalar Meson Dynamics
## 1 Introduction
Chiral perturbation theory for $`N_F=3`$ light flavours at next-to-leading order was shown by Gasser and Leutwyler to involve ten independent coupling constants $`L_i^r(\mu )`$ (referred to, sometimes, as LEC’s). In ref. all ten couplings were determined from low-energy data except two, $`L_4`$ and $`L_6`$. According to the OZI rule or, alternatively, large $`N_c`$ considerations these couplings are expected to be suppressed relative to the other ones. Physically, $`L_4`$ and $`L_6`$ control how observables of the pion, like $`M_\pi `$ or chiral $`SU(2)`$ order parameters $`(F_\pi )_{SU(2)}`$, $`\overline{u}u_{SU(2)}`$ vary if the mass of the strange quark mass $`m_s`$ varies. If OZI suppression holds these quantities are expected to be essentially insensitive to variations of $`m_s`$ and the value of the chiral order parameters $`F_\pi `$, $`\overline{u}u`$ in the $`SU(2)`$ chiral limit or in the $`SU(3)`$ chiral limit (i.e. $`m_s=0`$) ought to be very nearly the same. In nature, of course, the values of the light quark masses are fixed but they can be made to vary in unquenched lattice simulations of QCD.
There are reasons to suspect that the couplings $`L_4`$ and $`L_6`$ are, in fact, far from being suppressed. One reason is suggested by two recent unquenched lattice simulations which have investigated the phase structure of QCD-like theories with $`N_F`$ equal-mass flavours when the value of $`N_F`$ is varied. It was found in ref. that for $`N_F>6`$ a phase with no confinement and no spontaneous chiral symmetry breaking prevails, at any value of the QCD coupling constant. In ref. a very strong decrease of the chiral order parameters was observed upon varying $`N_F`$ from $`N_F=2`$ to $`N_F=4`$. The values of the coupling constants $`L_i`$ encode informations on the physics of the massive states in QCD. In particular, $`L_4`$, $`L_6`$ are linked to the scalar resonances. It is notorious that the OZI or large $`N_c`$ rules seem to fail in this sector. In fact, is not clear at present exactly which of the scalar resonances form the lowest lying nonet (see e.g. the reviews in the pdg). In other words, the fact that $`L_4`$, $`L_6`$ may be unsuppressed, the failure of the OZI rule in the scalar sector are connected and these features may be related in an interesting way to the phase structure of QCD, in particular to the fact that a phase transition could occur for a value of $`N_F`$ not exceedingly larger than $`N_F=3`$. An interpretation of these features in terms of a paramagnetic effect of the quark loops is discussed in ref.
These topics have started to be investigated in a previous paper. The couplings $`L_4,L_6`$ were related to the scalar resonances via the scalar form-factors of the pion and the Kaon. These form-factors can be reconstructed directly from experimental data on S-wave scattering, modulo a few plausible hypothesis. The coupling $`L_4`$ is related to the derivative of the strange form-factor of the pion at the origin and, for $`L_6`$, a chiral sum rule can be derived in the terms of the correlator $`\mathrm{\Pi }_6`$ of the strange scalar current $`\overline{s}s`$ and the non-strange current $`\overline{u}u+\overline{d}d`$. In the present paper, we investigate the $`O(m_s)`$ corrections to the results of ref., which could be sizable. For this purpose, one needs to use CHPT at order $`p^6`$. Renormalizability of the theory at this order was recently proved, and the set of independent chiral lagrangian terms classified. In particular, we will use the expansions of $`M_\pi `$ and $`F_\pi `$ which were derived in refs. and and we have computed the $`O(p^6)`$ expansion of the correlator $`\mathrm{\Pi }_6`$. The result of this effort, at first sight, will turn out to be mitigated: as new coupling constants appear in the formulas, there is a loss in predictivity over the leading order calculation. As discussed in sec.2, conclusions may nevertheless be drawn provided one makes some assumptions on the convergence of the chiral expansion.
In order to gain further insight, we have also investigated a different, more model dependent approach in sec.3. Following the idea of ref. one starts from a lagrangian for the scalar resonances and the LEC’s are generated by integrating out the resonances. If we attempt to determine some $`O(p^6)`$ LEC’s in this way, from the most general lagrangian, we find that there are many undetermined resonance parameters which get involved, like three resonance couplings, couplings of one resonance to two scalar sources etc… In sec.3, we discuss the more specific dynamics generated by the linear sigma-model. This model is still attracting interest in connection with the phenomenon of disordered chiral condensate which could be formed in heavy ion collisions. In the symmetry breaking sector, we will consider a lagrangian more general than previously done, still compatible with the criterion of softness (restricting ourselves to first order symmetry breaking). The model can accommodate exactly, in principle, a given set of scalar nonet masses and one obtains the $`O(p^4)`$ and $`O(p^6)`$ LEC’s as definite predictions. The question at this point is which scalar mesons are to be included in the nonet? In particular, does one have to include a light $`\sigma `$ and a light $`\kappa `$ meson? It has long been known that an S-matrix pole can be identified in low-energy $`\pi \pi `$ scattering with a very large imaginary part. A similar structure was argued to exist also in $`\pi K`$ scattering, but it is unclear whether such poles should be interpreted as physical scalar resonances. For our purposes, we will leave this question open and consider two different possibilities for the scalar nonet. As we will see, in the framework of the sigma-model, these will correspond to rather different behaviour of the chiral vacuum.
## 2 $`L_4`$, $`L_6`$ beyond $`O(p^4)`$
Our estimates of $`L_4`$ and $`L_6`$ are based, essentially, on a method to extract the scalar form-factors of the pion and the kaon using experimental $`\pi \pi K\overline{K}`$ scattering data, proposed in ref.. The form-factors get determined up to a normalization factor, for which one uses CHPT. In order to determine $`L_4`$ one equals the experimental determination of $`G_\pi ^{}(0)`$, the derivative of the strange scalar form-factor of the pion and the chiral expansion of this quantity. In ref. this matching was performed at leading order in the chiral expansion, i.e. $`G_\pi ^{}(0)`$ was expanded up to $`O(p^4)`$ and for the normalization condition we used CHPT at order $`p^2`$. In a similar way, for $`L_6`$ we use the scalar form factors to express the spectral function of the correlator $`\mathrm{\Pi }_6(s)`$ (see (12) below) and $`L_6`$ is obtained by matching $`\mathrm{\Pi }_6(0)`$ evaluated from experimental data and its chiral expansion. Below, we discuss the chiral corrections to these results, which involve two parts a) one must use the chiral expansions of $`G_\pi ^{}(0)`$ and $`\mathrm{\Pi }_6(0)`$ up to $`O(p^6)`$ and b) we must use CHPT at $`O(p^4)`$ instead of $`O(p^2)`$ in the normalization conditions.
Let us designate by $`F_\pi (s)`$ and $`G_\pi (s)`$ the non-strange and strange scalar form factors of the pion, and by $`F_K(s)`$ and $`G_K(s)`$ the analogous ones for the Kaon. We will use the following normalizations,
$`F_\pi ={\displaystyle \frac{1}{B_0}}\sqrt{{\displaystyle \frac{3}{2}}}0|\overline{u}u+\overline{d}d|\pi ^0\pi ^0`$ $`G_\pi ={\displaystyle \frac{1}{B_0}}\sqrt{{\displaystyle \frac{3}{2}}}0|\overline{s}s|\pi ^0\pi ^0`$ (1)
$`F_K={\displaystyle \frac{1}{B_0}}\sqrt{2}0|\overline{u}u+\overline{d}d|K^+K^{}`$ $`G_K={\displaystyle \frac{1}{B_0}}\sqrt{2}0|\overline{s}s|K^+K^{}.`$
The method of ref. consists in solving numerically a set of Muskhelishvili-Omnès coupled-channel equations using experimentally determined $`\pi \pi K\overline{K}`$ S-wave scattering T-matrix elements as input(in practice we made use of the parametrisations of ref. and ref.). Strictly speaking, these equations hold under the assumption of exact two-channel unitarity up to $`s=\mathrm{}`$ while in practice, two-channel unitarity is a good approximation up to the $`\eta \eta `$ threshold. The resulting form factors are expected to be reliable in a finite energy range, which we will assume to extend up to 1 GeV. One also chooses appropriate boundary conditions for the T-matrix at $`s=\mathrm{}`$ which insure existence of a solution with a minimal number of free parameters: a solution vector $`u_1(s),u_2(s)`$ will be determined in the entire energy range once the values at one point, say at $`s=0`$, are given. Extension to more channels and the stability of this scheme were discussed in ref.. In practice, one first constructs numerically two independent solutions of the equation set: $`u_i(s)`$ and $`v_i(s)`$, $`i=1,2`$, normalized at the origin such that $`u_1(0)=1`$, $`u_2(0)=0`$ and $`v_1(0)=0`$, $`v_2(0)=1`$. The form factors are then given as
$$\left(\genfrac{}{}{0pt}{}{F_\pi (s)}{F_K(s)}\right)=F_\pi (0)\left(\genfrac{}{}{0pt}{}{u_1(s)}{u_2(s)}\right)+F_K(0)\left(\genfrac{}{}{0pt}{}{v_1(s)}{v_2(s)}\right)$$
(2)
and similarly for the strange form factors. In refs. the normalization at the origin was taken from CHPT at $`O(p^2)`$. Here, we wish to investigate the chiral corrections to these results, so we need to go to the next chiral order.
### 2.1 Evolution of $`F_\pi `$
In the $`SU(2)`$ chiral limit, the strange form factor of the pion vanishes at the origin $`G_\pi (0)=0`$, the LEC $`L_4`$ is related to the first derivative at the origin of $`G_\pi `$. It is convenient to consider the proportional quantity
$$d_F^{exp}=\sqrt{\frac{2}{3}}m_sB_0G_\pi ^{}(0)\sqrt{\frac{2}{3}}m_sB_0G_K(0)v_1^{}(0).$$
(3)
Here, $`v_1^{}(0)`$ can be determined from experiment by the procedure outlined above, and one finds,
$$v_1^{}(0)0.27\mathrm{GeV}^2,$$
(4)
using the T-matrix parametrisation of ref. (using that of ref. one would find $`v_1^{}(0)0.31`$ GeV<sup>-2</sup>). We next need to express $`m_sB_0`$ and $`G_K(0)`$ in eq.(3). At leading chiral order, one has $`m_sB_0=m_K^2m_\pi ^2/2`$ and $`G_K(0)=\sqrt{2}`$. Including $`O(m_s)`$ corrections to this result, one obtains $`d_F^{exp}`$ in the form,
$$d_F^{exp}=\frac{2}{\sqrt{3}}\left[m_K^2m_\pi ^2/2+\frac{m_K^4}{F_\pi ^2}\left(8S_816S_6+\frac{1}{36\pi ^2}(L_\eta +1)\right)\right]v_1^{}(0),$$
(5)
with
$$S_8=2L_8^r(\mu )+L_5^r(\mu ),S_6=2L_6^r(\mu )+L_4^r(\mu ),L_\eta =\mathrm{log}\frac{m_\eta ^2}{\mu ^2}.$$
(6)
Corrections of order $`m_\pi ^2m_K^2`$ or $`m_\pi ^4`$ have been neglected here and CHPT at $`O(p^4)`$ has been used to obtain $`G_K(0)`$ and to express $`m_sB_0`$ in terms of the physical Kaon mass. We observe that low-energy constants appear now but it its consistent to use their values obtained at $`O(p^4)`$, i.e. for $`L_5`$, $`L_8`$ the values given in ref. and for $`L_4`$, $`L_6`$ those obtained in ref.. We can then match $`d_F^{exp}`$ with its expression as a chiral expansion. The latter is obtained from the relation
$$d_F=\frac{m_s}{F_\pi }\frac{dF_\pi }{dm_s}$$
(7)
which holds in the chiral $`SU(2)`$ limit, $`m_u=m_d=0`$, and one can use the chiral expansion of $`F_\pi `$ which has been determined up to $`O(p^6)`$ . One obtains
$$d_F=\frac{m_sB_0}{F_0^2}d^{(4)}+\frac{m_s^2B_0^2}{F_0^4}d^{(6)}+O(m_s^3),$$
(8)
with
$$d^{(4)}=8L_4^r\frac{1}{32\pi ^2}(L_K+1),L_K=\mathrm{log}\frac{m_sB_0}{\mu ^2},$$
(9)
and
$`d^{(6)}=`$ $`64C_{16}^r128(L_4^r)^2{\displaystyle \frac{17}{3072\pi ^4}}L_K^2{\displaystyle \frac{5}{1152\pi ^4}}L_KL_\eta `$
$`+{\displaystyle \frac{L_K}{\pi ^2}}\left({\displaystyle \frac{73}{9216\pi ^2}}+4L_1^r+L_2^r+{\displaystyle \frac{5}{4}}L_3^r+{\displaystyle \frac{1}{2}}L_5^r2L_6^rL_8^r\right)`$
$`+{\displaystyle \frac{L_\eta }{\pi ^2}}\left({\displaystyle \frac{1}{768\pi ^2}}+{\displaystyle \frac{16}{9}}L_1^r+{\displaystyle \frac{4}{9}}L_2^r+{\displaystyle \frac{4}{9}}L_3^r{\displaystyle \frac{8}{9}}L_4^r\right)`$
$`+{\displaystyle \frac{1}{\pi ^2}}\left({\displaystyle \frac{26}{9}}L_1^r+{\displaystyle \frac{35}{54}}L_3^r+{\displaystyle \frac{5}{9}}L_4^r+{\displaystyle \frac{3}{4}}L_5^r3L_6^r{\displaystyle \frac{3}{2}}L_8^r{\displaystyle \frac{0.0022}{\pi ^2}}\right).`$
In this formula, the finite contribution from the so-called sunset diagrams has been evaluated numerically. Equating (3) and (8) at the leading, linear order in $`m_s`$, gives the $`O(p^4)`$ determination of $`L_4`$. From the numerical results of refs. one obtains
$$L_4^r(m_\eta )\mathrm{0.6\hspace{0.17em}10}^3[\mathrm{order}p^4].$$
If we include the corrections quadratic in $`m_s`$ now, a new $`O(p^6)`$ coupling constant appears, labelled $`C_{16}`$ in ref., so strictly speaking, we have one equation and two unknowns. A similar situation will prevail in the determination of $`L_6`$ to be discussed below. We expect, however, the chiral expansion to be meaningful and the $`O(p^6)`$ part of the expansion to be smaller than the $`O(p^4)`$ one. If we assume a given ratio for these two parts we can determine $`L_4`$ and $`C_{16}`$ separately. Once these two couplings are known we can, furthermore, using the CHPT expressions of ref., determine how $`F_\pi `$ varies in going from an $`SU(2)`$ chiral limit to an $`SU(3)`$ one. Some results are collected in table 1, where we have varied the $`O(p^6)`$ over $`O(p^4)`$ ratio between 10% and 100%.
This exercise seems to indicate that the determination of $`L_4`$ is rather stable and rather close to its determination at leading order, and provides an estimate for the size of the $`O(p^6)`$ coupling constant $`C_{16}^r`$. A priori, however, one cannot completely exclude a different solution, with $`L_4(m_\eta )=0`$, for instance. In this case, all the OZI violation would be concentrated in the $`O(p^6)`$ parameter, which seems somewhat unplausible. Taking into account the experimental uncertainties in the T-matrix, the uncertainty from the energy region above 1 GeV (which were estimated in ref.) and the discussion above of the $`O(m_s)`$ corrections, we find that the error on $`L_4`$ should be of the order of 30%, i.e.,
$$L_4(m_\eta )=(0.6\pm 0.2)\mathrm{\hspace{0.17em}10}^3.$$
(11)
### 2.2 Evolution of $`\overline{u}u`$
Let us now discuss the coupling constant $`L_6`$. Using the same scalar form factors, it was remarked in ref. that one can derive an estimate of $`L_6`$. The idea is to consider the two-point correlation function,
$$\mathrm{\Pi }_6(p^2)=\frac{i}{B_0^2}d^4x\mathrm{e}^{ipx}0|T(\overline{u}u(x)+\overline{d}d(x))\overline{s}s(0)|0_{conn.}.$$
(12)
In an energy domain extending roughly up to the $`\eta \eta `$ threshold, the spectral function is given in terms of the scalar form factors of the pion and the Kaon,
$`Im\mathrm{\Pi }_6(s)=`$ $`\sqrt{{\displaystyle \frac{s4m_\pi ^2}{s}}}F_\pi (s)G_\pi ^{}(s)\theta (s4m_\pi ^2)`$ (13)
$`+\sqrt{{\displaystyle \frac{s4m_K^2}{s}}}F_K(s)G_K^{}(s)\theta (s4m_K^2).`$
Essentially, a usual assumption is made there, that the $`4\pi `$ or $`6\pi `$ contributions are negligibly small in this energy region (see, e.g.). In the $`SU(2)`$ chiral limit, $`m_u=m_d=0`$, a super-convergence relation holds,
$$_0^{\mathrm{}}Im\mathrm{\Pi }_6(s)𝑑s=0.$$
(14)
Similar chiral sum rules are known to be saturated to a fairly good approximation in terms of a few low-lying resonances (see, e.g.). In the energy region below 1 GeV, inserting the form factors constructed as discussed above, the contribution to the integral is found to be positive, and dominated by the $`f_0(980)`$ resonance. A plausible assumption, then, is that the higher energy contribution to the sum rule (14) will be saturated, at least approximately, by the next prominent scalar resonance, the $`f_0(1500)`$, which we expect to make a negative contribution. We observe that the signs of the contributions of the $`f_0(980)`$ and the $`f_0(1500)`$ conforms with the recent assignments of Minkowski and Ochs as (essentially) $`SU_F(3)`$ singlet and octet respectively. Finally, one can calculate the value of the correlation function at zero energy, $`\mathrm{\Pi }_6(0)`$, making use of the scalar form factors below 1 GeV, and the super-convergence relation (14) as,
$$\mathrm{\Pi }_6(0)\frac{1}{\pi }\left(_0^1\frac{Im\mathrm{\Pi }_6(s)}{s}𝑑s\frac{1}{s_0}_0^1Im\mathrm{\Pi }_6(s)𝑑s\right),$$
(15)
with $`s_01.5^2`$ GeV<sup>2</sup>. This gives us an “experimental” value of $`\mathrm{\Pi }_6(0)`$, the difference with ref. is that we now wish to take into account chiral corrections in the normalization of the form factors. Below the $`K\overline{K}`$ threshold, we can express the spectral function in terms of the basis solutions to the Muskhelishvili-Omnès equations $`u_i(s)`$, $`v_i(s)`$ as,
$`Im\mathrm{\Pi }_6(s)=\sqrt{{\displaystyle \frac{s4m_\pi ^2}{s}}}[`$ $`F_\pi (0)G_K(0)u_1(s)v_1^{}(s)`$ (16)
$`+F_K(0)G_K(0)v_1(s)v_1^{}(s)].`$
Including $`O(m_s)`$ chiral corrections in the normalizing factors, we now have
$`F_\pi (0)G_K(0)=2\sqrt{3}\left(1+{\displaystyle \frac{m_K^2}{F_\pi ^2}}\left[16S_848S_6+{\displaystyle \frac{1}{24\pi ^2}}L_\eta +{\displaystyle \frac{1}{36\pi ^2}}\right]\right)`$ (17)
$`F_K(0)G_K(0)=2\left(1+{\displaystyle \frac{m_K^2}{F_\pi ^2}}\left[32S_880S_6+{\displaystyle \frac{7}{72\pi ^2}}L_\eta +{\displaystyle \frac{1}{24\pi ^2}}\right]\right),`$
where $`S_8=2L_8^r+L_5^r`$, $`S_6=2L_6^r+L_4^r`$ and we may use in this part the values of the LEC’s as determined at $`O(p^4)`$.
We then use the experimental result for $`\mathrm{\Pi }_6(0)`$ in conjunction with its chiral expansion. The computation of $`\mathrm{\Pi }_6(s)`$ to $`O(p^6)`$ is exposed in the appendix. Concerning $`\mathrm{\Pi }_6(0)`$, the chiral expansion goes as follows
$$\mathrm{\Pi }_6(0)=\mathrm{\Pi }^{(4)}+\frac{m_sB_0}{F_0^2}\mathrm{\Pi }^{(6)},$$
(18)
with
$$\mathrm{\Pi }^{(4)}=64L_6^r\frac{1}{16\pi ^2}\left(\frac{22}{9}(L_K+1)+\frac{4}{9}L_{43}\right),$$
(19)
and
$$L_K=\mathrm{log}\frac{m_sB_0}{\mu ^2},L_{43}=\mathrm{log}\frac{4}{3}.$$
(20)
This expression generates the $`O(p^4)`$ determination of $`L_6`$,
$$L_6^r(m_\eta )\mathrm{0.5\hspace{0.17em}10}^3[\mathrm{order}p^4].$$
(21)
The expression of the $`O(p^6)`$ part in (18) is,
$`\mathrm{\Pi }^{(6)}=`$ $`256(C_{20}^r+3C_{21}^r){\displaystyle \frac{1}{72\pi ^4}}L_K^2{\displaystyle \frac{1}{96\pi ^4}}L_KL_{43}+{\displaystyle \frac{1}{288\pi ^4}}L_{43}^2`$
$`+{\displaystyle \frac{L_K}{\pi ^2}}\left({\displaystyle \frac{64}{3}}S_6+{\displaystyle \frac{70}{9}}S_8+{\displaystyle \frac{16}{9}}S_7{\displaystyle \frac{55}{2592\pi ^2}}\right)`$
$`+{\displaystyle \frac{L_{43}}{\pi ^2}}\left({\displaystyle \frac{16}{3}}S_6+{\displaystyle \frac{16}{9}}S_8+{\displaystyle \frac{16}{9}}S_7{\displaystyle \frac{47}{5184\pi ^2}}\right)`$
$`+{\displaystyle \frac{1}{\pi ^2}}\left({\displaystyle \frac{140}{9}}S_6+{\displaystyle \frac{175}{27}}S_8+{\displaystyle \frac{8}{27}}S_7{\displaystyle \frac{19}{5184\pi ^2}}\right).`$
with
$$S_7=3L_7+L_8^r.$$
(23)
In these formulas, we have set $`m_u=m_d=0`$. These expressions exhibit the explicit $`m_s`$ dependence (which will prove usefull below as we plan to integrate over $`m_s`$). For numerical application, we may use $`m_sB_0=m_K^2m_\pi ^2/2`$ in $`\mathrm{\Pi }^{(6)}`$ while in $`\mathrm{\Pi }^{(4)}`$ we have to use the $`O(p^4)`$ expansions of $`m_K^2`$ and $`m_\pi ^2`$. Equating the chiral expansion of $`\mathrm{\Pi }_6(0)`$ with its sum rule evaluation gives us a linear relation involving the $`O(p^4)`$ LEC $`L_6`$ and the $`O(p^6)`$ combination $`C_{20}+3C_{21}`$. As before, we must make a hypothesis concerning the convergence of the chiral expansion if we want to separately estimate these two contributions.
Another interesting information contained in the correlation function $`\mathrm{\Pi }_6`$ concerns the variation of the quark condensate in the chiral $`SU(2)`$ limit as a function of the strange quark mass:
$$\frac{d\overline{u}u}{dm_s}=\frac{B_0^2}{2}\mathrm{\Pi }_6(0).$$
(24)
Using the chiral expansion of $`\mathrm{\Pi }_6(0)`$ we can integrate this equation in the variable $`m_s`$ from its physical value down to $`m_s=0`$. In this way we obtain an estimate for how the condensate varies from a chiral $`SU(2)`$ limit to a chiral $`SU(3)`$ limit. This variation is obtained as an expansion in powers of the physical strange quark mass (or, alternatively, in powers of the Kaon mass),
$$\overline{u}u_{SU(2)}=\overline{u}u_{SU(3)}\left(1+\frac{m_sB_0}{F_0^2}R^{(4)}+\frac{m_s^2B_0^2}{F_0^4}R^{(6)}\right),$$
(25)
with
$$R^{(4)}=32L_6^r\frac{11}{144\pi ^2}L_K\frac{1}{72\pi ^2}L_{43},$$
(26)
and
$`R^{(6)}=`$ $`64(C_{20}^r+3C_{21}^r){\displaystyle \frac{1}{288\pi ^4}}L_K^2{\displaystyle \frac{1}{384\pi ^4}}L_KL_{43}+{\displaystyle \frac{1}{1152\pi ^4}}L_{43}^2`$
$`+{\displaystyle \frac{L_K}{\pi ^2}}\left({\displaystyle \frac{16}{3}}S_6+{\displaystyle \frac{4}{9}}S_7+{\displaystyle \frac{35}{18}}S_8{\displaystyle \frac{19}{10368\pi ^2}}\right)`$
$`+{\displaystyle \frac{L_{43}}{\pi ^2}}\left({\displaystyle \frac{4}{3}}S_6+{\displaystyle \frac{4}{9}}S_7+{\displaystyle \frac{4}{9}}S_8{\displaystyle \frac{5}{5184\pi ^2}}\right)`$
$`+{\displaystyle \frac{1}{\pi ^2}}\left({\displaystyle \frac{11}{9}}S_6{\displaystyle \frac{4}{27}}S_7+{\displaystyle \frac{35}{54}}S_8\right).`$
Using the relation
$$2\overline{u}u_{SU(2)}=\frac{d}{dm}\left(m_\pi ^2F_\pi ^2\right)_{m=0},m=\frac{1}{2}(m_u+m_d),$$
(28)
one sees that the expansion of the condensate can be rederived from the expansion of the product $`m_\pi ^2F_\pi ^2`$ in powers of $`m_s`$. We have verified, using the $`O(p^6)`$ expansion of $`m_\pi `$, $`F_\pi `$ that the formulas (25), (26), (2.2) are exactly reproduced.
Let us now discuss some numerical results. At first, it is instructive to compare the value obtained for $`\mathrm{\Pi }_6(0)`$ using the normalization of the form factors at $`O(p^2)`$ and the normalization at $`O(p^4)`$,
$$\mathrm{\Pi }_6(0)0.022[O(p^2)\mathrm{norm}.],\mathrm{\Pi }_6(0)0.043[O(p^4)\mathrm{norm}.].$$
(29)
We see that including $`O(m_s)`$ corrections in the normalizations brings a rather substantial change in the result by approximately a factor of two. This can be traced to the large numbers appearing in front of the combination $`S_6=2L_6+L_4`$ in the normalization factors (17). Evidently, the difference $`2L_6+L_4`$ is very sensitive to a small variation of $`L_6`$ or $`L_4`$. Taking $`L_6`$ slightly smaller than the central value (21) would decrease significantly the modification in $`\mathrm{\Pi }_6(0)`$ without essentially altering the other results. Upon considering the chiral expansion of $`\mathrm{\Pi }_6(0)`$ now, this result suggests that the $`O(p^6)`$ contribution could be substantial and it is likely to be positive. We have collected some results in table 2 assuming that the ratio of the $`O(p^6)`$ to the $`O(p^4)`$ contributions is positive and ranges from 10% to 100%. In this range, we find that the rate of convergence for the ratio of quark condensates seems reasonable.
Both $`F_\pi `$ and $`\overline{u}u`$ show a tendency towards chiral restoration in going from chiral $`SU(2)`$ to chiral $`SU(3)`$. This tendency seems much stronger for $`\overline{u}u`$. This is in agreement with the arguments of ref. based on the spectrum of the Dirac operator. One must however bear in mind that its dimensionality is different from $`F_\pi `$: if we had considered $`\overline{u}u^{1/3}`$, we would have found a smaller variation. This behaviour could perhaps suggest the possibility, at larger $`N_F`$, of a phase with vanishing quark condensate, yet with chiral symmetry still spontaneously broken. Kogan et al. have discussed how this could result from a discrete $`Z_2`$ axial subgroup remaining unbroken, but argue against this possibility in QCD. Concerning the error in this evaluation of $`L_6`$, finally, it is expected to be somewhat larger than the error on $`L_4`$ because of a greater sensitivity to the energy region above 1 GeV. The discussion above also suggests that one should have $`L_6(m_\eta )<L_4(m_\eta )`$ otherwise the $`O(m_s)`$ corrections could become out of control. Keeping this mind, we find that $`L_6`$ should lie in the following range,
$$L_6(m_\eta )=(0.5\pm 0.3)\mathrm{\hspace{0.17em}10}^3.$$
(30)
## 3 The LEC’s in a linear sigma-model
### 3.1 General scalar meson lagrangian
Here, we adopt a simple resonance saturation point of view for estimating the low-energy coupling constants, which was discussed in detail in ref.. It consists in making a tree level calculation starting from a lagrangian for the resonances. We will consider the scalar resonances here. It is most convenient to start from a representation of the resonances in which they transform under a non-linear representation of the chiral group. A detailed discussion can be found in ref., we adopt essentially the same notations here. If one is interested in the $`O(p^4)`$ LEC’s then one needs consider only those terms which are quadratic in the resonance fields,
$$_{SS}=\frac{1}{2}_\mu S_0^\mu S_0+\frac{1}{2}_\mu S^\mu S\frac{1}{2}M_0^2S_0^2\frac{1}{2}M_8^2S^2,$$
(31)
or containing one resonance field and one scalar source, or one scalar field and two chiral fields $`u_\mu `$, i.e.
$$_{S\chi }=c_dSu_\mu u^\mu +\stackrel{~}{c}_dS_0u_\mu u_\mu +c_mS\chi _++\stackrel{~}{c}_mS_0\chi _+,$$
(32)
where $`S_0`$ is the singlet scalar, with chiral limit mass $`M_0`$, $`S`$ is a traceless matrix encoding the scalar octet, which has a common mass $`M_8`$ in the chiral limit. We use the same conventions otherwise as ref.. This lagrangian yields the resonance saturation estimates for the following low-energy constants,
$`L_4={\displaystyle \frac{c_dc_m}{3M_8^2}}+{\displaystyle \frac{\stackrel{~}{c}_d\stackrel{~}{c}_m}{M_0^2}},L_5={\displaystyle \frac{c_dc_m}{M_8^2}},`$
$`L_6={\displaystyle \frac{c_m^2}{6M_8^2}}+{\displaystyle \frac{\stackrel{~}{c}_m^2}{2M_0^2}},L_8={\displaystyle \frac{c_m^2}{2M_8^2}}.`$ (33)
In addition, the scalar resonances make contributions to $`L_1`$, $`L_3`$,
$$L_1=\frac{c_d^2}{6M_8^2}+\frac{\stackrel{~}{c}_d^2}{2M_0^2},L_3=\frac{c_d^2}{2M_8^2},$$
(34)
which receive other important contributions from the vector and axial-vector resonance sector. These tree-level estimates yield coupling constants which are scale independent<sup>1</sup><sup>1</sup>1If one starts from the renormalizable sigma-model it is possible, in principle, to compute the one-loop effective action and then generate the $`O(p^4)`$ LEC’s with the correct scale dependence. This was done for the SU(2) sigma-model in ref. : one usually assumes that they should represent meaningful estimates of the $`L_i^r(\mu )`$ for values of $`\mu `$ for which chiral logarithms are numerically small, that is of the order of $`\mu =0.5`$ to $`\mu =1`$ GeV. In ref. the OZI rule was assumed to apply, implying $`L_4=L_6=0`$. The formulas (3.1) provide one relation between the values of $`L_5`$, $`L_8`$ and the experimental value of the decay width $`a_0\pi \eta `$. Values of the couplings $`c_d`$ and $`c_m`$ were extracted,
$$|c_d|32\mathrm{MeV}|c_m|42\mathrm{MeV}.$$
(35)
Below, we will investigate the more detailed predictions that one can make if one assumes the validity of a renormalizable, sigma-model type interaction lagrangian. Before we discuss this, let us consider the generalisation of the considerations above to LEC’s of chiral order $`O(p^6)`$. We will illustrate the complications which arise at this order by considering only the lagrangian terms which are cubic in the scalar source $`\chi _+`$. These terms involve the three coupling constants $`C_{19}`$, $`C_{20}`$, $`C_{21}`$,
$$^{(6)}=\mathrm{}+C_{19}\chi _+^3+C_{20}\chi _+^2\chi _++C_{21}\chi _+^3+\mathrm{}$$
(36)
Firstly, one needs to consider interaction terms which are cubic in the resonance fields,
$$_{SSS}=aS_0^3+bS_0S^2+cS^3,$$
(37)
then terms which are quadratic in the resonance fields and linear in the scalar source,
$`_{SS\chi }=e_mS^2\chi _++f_mS_0S\chi _+`$
$`+g_mS^2\chi _++h_mS_0^2\chi _+,`$ (38)
and, finally, terms linear in the scalar field and quadratic in the scalar source,
$`_{S\chi \chi }=a^{}S_0\chi _+^2+b^{}S_0\chi _+^2+c^{}S\chi _+\chi _+`$
$`+d^{}S\chi _+^2.`$ (39)
This already brings in a large number of couplings and it would not be possible to make definite predictions without additional assumptions. The terms in eqs.(37),(3.1) and (3.1) generate chiral $`O(p^6)`$ coupling constants as well as chiral symmetry breaking for the scalar meson masses. It is convenient, both for generating the $`O(p^4)`$ and $`O(p^6)`$ LEC’s and for expressing the scalar meson masses to make the following field redefinitions,
$`S_0S_0+{\displaystyle \frac{a}{M_0^2}}S_0^2+{\displaystyle \frac{b}{M_0^2}}S^2+{\displaystyle \frac{\stackrel{~}{c}_m}{M_0^2}}\chi _+`$
$`SS+{\displaystyle \frac{c}{M_8^2}}\left(S^2{\displaystyle \frac{1}{3}}S^2I_d\right)+{\displaystyle \frac{c_m}{M_8^2}}\left(\chi _+{\displaystyle \frac{1}{3}}\chi _+I_d\right)`$ (40)
The redefined fields $`S_0`$ and $`S`$ have, by construction, no trilinear couplings and no minimal one-resonance to one-source coupling. The couplings of two resonances to one scalar source, eq.(3.1), get redefined to
$`g_m=g_m+{\displaystyle \frac{3c_mc}{M_8^2}},h_m=h_m+{\displaystyle \frac{3\stackrel{~}{c}_ma}{M_0^2}},`$
$`e_m=e_m+{\displaystyle \frac{\stackrel{~}{c}_mb}{M_0^2}}{\displaystyle \frac{c_mc}{M_8^2}},f_m=f_m+{\displaystyle \frac{2c_mb}{M_8^2}}.`$ (41)
The scalar meson masses, to linear order in the quark masses, can be straightforwardly expressed in terms of these redefined parameters. The masses of the isospin $`I=1`$ and $`I=1/2`$ mesons read,
$`M_{a_0}^2=M_8^24e_m(2m_K^2+m_\pi ^2)4g_mm_\pi ^2`$
$`M_{\kappa _0}^2=M_{a_0}^24g_m(m_K^2m_\pi ^2),`$ (42)
while the isospin zero meson masses diagonalise the singlet-octet mass matrix which gets parametrised as follows
$`_{11}=M_0^24h_m(2m_K^2+m_\pi ^2)`$
$`_{12}=4f_m\sqrt{{\displaystyle \frac{2}{3}}}(m_K^2m_\pi ^2)`$
$`_{22}=M_{a_0}^2{\displaystyle \frac{16}{3}}g_m(m_K^2m_\pi ^2).`$ (43)
In deriving these expressions for the masses, we have ignored possible terms in the scalar lagrangian of the form $`_\mu S_0^\mu S_0\chi _+`$, $`_\mu S^\mu S\chi _+`$, anticipating on the fact that such terms will not appear in the model to be discussed below.
Concerning the LEC’s $`C_{19}`$, $`C_{20}`$, $`C_{21}`$, the following expressions are obtained,
$`F_0^2C_{19}={\displaystyle \frac{c_m^2}{M_8^4}}(g_m+{\displaystyle \frac{cc_m}{M_8^2}})+d^{}{\displaystyle \frac{c_m}{M_8^2}}`$
$`F_0^2C_{20}={\displaystyle \frac{c_m^2}{M_8^4}}({\displaystyle \frac{2}{3}}g_m+e_m{\displaystyle \frac{cc_m}{M_8^2}})+{\displaystyle \frac{c_m\stackrel{~}{c}_m}{M_0^2M_8^2}}(f_m+{\displaystyle \frac{bc_m}{M_8^2}})`$
$`+(c^{}{\displaystyle \frac{1}{3}}d^{}){\displaystyle \frac{c_m}{M_8^2}}+b^{}{\displaystyle \frac{\stackrel{~}{c}_m}{M_0^2}}`$
$`F_0^2(C_{20}+3C_{21})={\displaystyle \frac{c_m^2}{3M_8^4}}(g_m+{\displaystyle \frac{cc_m}{M_8^2}})+{\displaystyle \frac{3\stackrel{~}{c}_m^2}{M_0^4}}(h_m+{\displaystyle \frac{a\stackrel{~}{c}_m}{M_0^2}})`$
$`+(3a^{}+b^{}){\displaystyle \frac{\stackrel{~}{c}_m}{M_0^2}}d^{}{\displaystyle \frac{c_m}{3M_8^2}}.`$ (44)
### 3.2 Parameters of the linear sigma-model
In the $`SU(3)`$ linear sigma-model, one first encodes the scalar nonet and the pseudo-scalar nonet into a complex $`3\times 3`$ matrix $`\mathrm{\Sigma }`$, which, under a chiral transformation transforms as,
$$\mathrm{\Sigma }g_R\mathrm{\Sigma }g_L^{},$$
(45)
and the lagrangian is assumed to be renormalizable. In the chiral limit, the most general renormalizable lagrangian, invariant under the chiral group (except for the $`U_A(1)`$ subgroup) contains four parameters,
$`={\displaystyle \frac{1}{2}}_\mu \mathrm{\Sigma }^\mu \mathrm{\Sigma }^{}{\displaystyle \frac{1}{2}}\mu ^2\mathrm{\Sigma }\mathrm{\Sigma }^{}\lambda (\mathrm{\Sigma }\mathrm{\Sigma }^{})^2`$
$`\lambda ^{}\mathrm{\Sigma }\mathrm{\Sigma }^{}^2+\beta (\mathrm{det}(\mathrm{\Sigma })+\mathrm{det}(\mathrm{\Sigma }^{})).`$ (46)
This lagrangian was reconsidered recently by several authors and we use the same notations as in ref.. Suitable choices of the parameters $`\mu ^2,\lambda ,\lambda ^{},\beta `$ ensure that spontaneous chiral symmetry breaking occurs, i.e. that the potential is minimised for a non-zero value of the vacuum expectation value of the $`\mathrm{\Sigma }`$ matrix,
$$\overline{\mathrm{\Sigma }}=vI_d,$$
(47)
(which we take to be diagonal, assuming that no spontaneous breaking of flavour symmetry occurs). Values of $`v`$ which correspond to extremums of the potential are solutions of the cubic equation,
$$v\left((4\lambda +12\lambda ^{})v^22\beta v+\mu ^2\right)=0,$$
(48)
which has real solutions $`v0`$ provided the parameters satisfy
$$\beta ^2\mu ^2(4\lambda +12\lambda ^{}).$$
(49)
Next, expanding around the minimum we find that the pseudoscalars are massless, except for the singlet, $`\eta _0`$ which mass is proportional to $`\beta `$,
$$M_{\eta _0}^2=6\beta v.$$
(50)
We can also read off the expressions for the singlet and octet scalar meson masses in the chiral limit, $`M_0`$ and $`M_8`$,
$$M_8^2=8\lambda v^2+4\beta v,M_0^2=8(\lambda +3\lambda ^{})v^22\beta v.$$
(51)
Stability of the vacuum requires that these squared masses be positive. Using eqs.(48)(50) and (51) we can trade the original four parameters of the lagrangian (3.2) for the more physically relevant ones $`v`$, $`M_{\eta _0}`$, $`M_0`$ and $`M_8`$.
In QCD, chiral symmetry is broken by the light quark masses. We will make the simplifying assumption, at the level of the sigma-model, to consider symmetry breaking to linear order in the quark masses but we accept all the terms which are soft (i.e. of dimensionality strictly smaller than four). Under these assumptions the most general symmetry breaking sector has, again, four parameters,
$`_{SB}=\gamma _0\mathrm{\Sigma }\chi ^{}+\gamma _1\mathrm{\Sigma }^1\chi \mathrm{det}(\mathrm{\Sigma })+\gamma _2\mathrm{\Sigma }\mathrm{\Sigma }^{}\mathrm{\Sigma }\chi ^{}`$
$`+\gamma _3\mathrm{\Sigma }\mathrm{\Sigma }^{}\mathrm{\Sigma }\chi ^{}+h.c..`$ (52)
Usually, only the first term is considered. This increased phenomenological flexibility will allow us to obtain better fits of the scalar nonet masses. While more general, this lagrangian will prove nevertheless to be reasonably constraining. We note that only the first two terms are renormalizable in the strict sense that the counterterms are exactly of the same form. The last two terms generate counterterms which are of higher order in the scalar source $`\chi `$. Consistent with the assumption of not including further terms quadratic in the quark mass matrix at this level, we will compute the scalar meson masses and the vacuum expectation values at linear order as well. We have checked this approximation by comparing with exactly calculated masses.
The connection with the representation of the scalar fields used in sec.3.1 is performed by making a change of variables
$$\mathrm{\Sigma }=\mathrm{exp}\left(i\sqrt{\frac{2}{3}}\frac{\eta _0}{F_0}\right)uHu$$
(53)
with
$$u=\mathrm{exp}\left(i\underset{1}{\overset{8}{}}\frac{\lambda _a\pi _a}{2F_0}\right),H=(v+\frac{S_0}{\sqrt{3}})I_d+S.$$
(54)
This change of variable makes sense only if $`F_00`$, i.e. when chiral symmetry is spontaneously broken. The kinetic energy part of the linear sigma-model lagrangian gets transformed into
$`_{kin}=`$ $`{\displaystyle \frac{1}{2}}_\mu S_0^\mu S_0+{\displaystyle \frac{1}{2}}_\mu S^\mu S+{\displaystyle \frac{1}{8}}\{u_\mu ,H\}^2+{\displaystyle \frac{1}{3F_0^2}}_\mu \eta _0^\mu \eta _0H^2`$ (55)
$`\sqrt{{\displaystyle \frac{2}{3}}}{\displaystyle \frac{_\mu \eta _0}{F_0}}H^2u^\mu `$
with
$$_\mu S=_\mu S+[\mathrm{\Gamma }_\mu ,S],\mathrm{\Gamma }_\mu =\frac{1}{2}(u^{}_\mu u+u_\mu u^{}),u_\mu =i(u^{}_\mu uu_\mu u^{}).$$
(56)
Replacing $`H`$ by its expression above (54) one identifies $`v`$ with the pion decay constant in the chiral limit $`F_0`$ and one obtains a prediction for the two couplings $`c_d`$ and $`\stackrel{~}{c}_d`$,
$$v^2=\frac{F_0^2}{2},c_d=v,\stackrel{~}{c}_d=\frac{v}{\sqrt{3}}.$$
(57)
The symmetry breaking lagrangian, expressed in terms of the new fields reads
$`_{SB}=`$ $`(\gamma _0H+\gamma _2H^3+\gamma _3H^2H)(\chi _+\mathrm{cos}\sqrt{{\displaystyle \frac{2}{3}}}{\displaystyle \frac{\eta _0}{F_0}}i\chi _{}\mathrm{sin}\sqrt{{\displaystyle \frac{2}{3}}}{\displaystyle \frac{\eta _0}{F_0}})`$ (58)
$`+\gamma _1\mathrm{det}HH^1(\chi _+\mathrm{cos}\sqrt{{\displaystyle \frac{8}{3}}}{\displaystyle \frac{\eta _0}{F_0}}+i\chi _{}\mathrm{sin}\sqrt{{\displaystyle \frac{8}{3}}}{\displaystyle \frac{\eta _0}{F_0}})`$
In this sector, we obtain one relation among the four parameters $`\gamma _i`$ from the requirement that the coefficient of the chiral lagrangian term $`\chi _+`$ be correctly normalized to $`F_0^2/4`$ and we can express two other parameters in terms of the meson-source couplings $`c_m`$ and $`\stackrel{~}{c}_m`$. This gives the relations
$`\gamma _0={\displaystyle \frac{3}{4}}v+{\displaystyle \frac{1}{6}}c_m{\displaystyle \frac{2}{\sqrt{3}}}\stackrel{~}{c}_m+\gamma _3v^2`$
$`\gamma _1v={\displaystyle \frac{1}{3}}(c_m\sqrt{3}\stackrel{~}{c}_m)2\gamma _3v^2`$
$`\gamma _2v^2={\displaystyle \frac{1}{4}}v+{\displaystyle \frac{1}{6}}c_m+{\displaystyle \frac{1}{\sqrt{3}}}\stackrel{~}{c}_m2\gamma _3v^2.`$ (59)
In this sector, we can take $`c_m`$, $`\stackrel{~}{c}_m`$ and $`\gamma _3`$ as arbitrary parameters.
In this model, we can now express the many parameters that appeared in the general scalar lagrangian discussed above. Firstly, for the terms coupling one scalar field to two scalar sources we have simply $`a^{}=b^{}=c^{}=d^{}=0`$. For the trilinear couplings of resonances, one obtains the relations
$`a={\displaystyle \frac{1}{2\sqrt{3}v}}\left(M_0^2+{\displaystyle \frac{1}{9}}M_{\eta _0}^2\right)`$
$`b={\displaystyle \frac{1}{2\sqrt{3}v}}\left(M_0^2+2M_8^2{\displaystyle \frac{2}{3}}M_{\eta _0}^2\right)`$
$`c={\displaystyle \frac{1}{2v}}\left(M_8^2{\displaystyle \frac{8}{9}}M_{\eta _0}^2\right).`$ (60)
Finally, the parameters $`g_m,h_m,e_m,f_m`$ which control linear symmetry breaking of the scalar meson masses (see eqs.(3.1), (3.1)) obey the following relations,
$`g_m={\displaystyle \frac{4}{3v}}\left(\sqrt{3}\stackrel{~}{c}_mc_m+c_m{\displaystyle \frac{M_{\eta _0}^2}{M_8^2}}\right){\displaystyle \frac{3}{4}}8\gamma _3v`$
$`\sqrt{3}f_m={\displaystyle \frac{1}{v}}\left(\sqrt{3}\stackrel{~}{c}_mc_m{\displaystyle \frac{M_0^2}{M_8^2}}\right){\displaystyle \frac{9}{8}}+{\displaystyle \frac{1}{2}}g_m`$
$`e_m={\displaystyle \frac{M_{\eta _0}^23M_8^2}{9M_0^2v}}\left(\sqrt{3}\stackrel{~}{c}_mc_m{\displaystyle \frac{M_0^2}{M_8^2}}\right){\displaystyle \frac{3}{16}}{\displaystyle \frac{1}{4}}g_m`$
$`h_m={\displaystyle \frac{M_{\eta _0}^2}{18M_0^2v}}\left(\sqrt{3}\stackrel{~}{c}_mc_m{\displaystyle \frac{M_0^2}{M_8^2}}\right){\displaystyle \frac{9}{32}}{\displaystyle \frac{1}{24}}g_m.`$ (61)
Making use of the relations (3.1) we can express the three $`O(p^6)`$ LEC’s $`C_{19}`$, $`C_{20}`$, $`C_{21}`$ discussed above in terms of the parameters of the linear sigma-model. We are also interested in the prediction for the OZI suppressed $`O(p^6)`$ LEC $`C_{16}`$ which participates in the flavour variation of the order parameter $`F_\pi `$ as discussed above. After a small calculation, the following expression is obtained,
$`F_0^2C_{16}=`$ $`{\displaystyle \frac{1}{54}}(2c_m+10\sqrt{3}\stackrel{~}{c}_m9v)\left({\displaystyle \frac{1}{M_0^2}}{\displaystyle \frac{1}{M_8^2}}\right)\left({\displaystyle \frac{2c_m}{M_8^2}}+{\displaystyle \frac{\sqrt{3}\stackrel{~}{c}_m}{M_0^2}}\right)`$
$`+{\displaystyle \frac{\gamma _3v^2}{9}}\left({\displaystyle \frac{2\sqrt{3}\stackrel{~}{c}_m}{M_0^2}}\left({\displaystyle \frac{1}{M_0^2}}+{\displaystyle \frac{2}{M_8^2}}\right){\displaystyle \frac{4c_m}{M_8^2}}\left({\displaystyle \frac{2}{M_0^2}}{\displaystyle \frac{5}{M_8^2}}\right)\right)`$
$`+{\displaystyle \frac{M_{\eta _0}^2}{54M_0^2M_8^2}}\left({\displaystyle \frac{2c_m}{M_8^2}}+{\displaystyle \frac{\sqrt{3}\stackrel{~}{c}_m}{M_0^2}}\right)\left(2c_m(12{\displaystyle \frac{M_0^2}{M_8^2}})\sqrt{3}c_m{\displaystyle \frac{M_8^2}{M_0^2}}\right).`$
### 3.3 Phenomenological applications
Altogether, the sigma-model as considered here has six independent parameters (apart from $`F_0`$): three chiral limit masses $`M_0`$, $`M_8`$ and $`M_{\eta _0}`$ and three couplings to the scalar source $`c_m`$, $`\stackrel{~}{c}_m`$ and $`\gamma _3`$. The parameter $`M_{\eta _0}`$, i.e. the mass of the $`\eta ^{}`$ in the chiral limit can be estimated to be,
$$M_{\eta _0}0.8\mathrm{GeV}.$$
(63)
Let us now discuss the determination of the remaining five parameters. A priori, one expects to be able to reproduce exactly the four independent masses ($`M_{a_0}`$, $`M_{\kappa _0}`$, $`M_{\sigma _0}`$, $`M_{f_0}`$) in the scalar nonet and, as an additional constraint a natural choice would be to enforce the super-convergence relation (14) which implies here $`c_m=\sqrt{3}\stackrel{~}{c}_m`$. In practice, however, because of non-linearities, this choice strongly restricts the possible range of masses for the $`\kappa _0`$. This is an indication that this model should be considered as a toy model rather than a really physically meaningful one. Instead, we did not try to impose the super-convergence relation and use a constraint from the pseudo-scalar sector, the ratio $`F_K/F_\pi `$ which allows for a fairly broad range of values for $`M_{\kappa _0}`$. Because of the specific parametrisation of the mass matrix, eq.(3.1) its entries get evaluated as follows in terms of physical masses,
$$_{22}=\frac{4M_{\kappa _0}^2M_{a_0}^2}{3},_{11}=M_{\sigma _0}^2+M_{f_0}^2_{22},$$
(64)
and, using the determinant,
$$_{12}=\pm \sqrt{(_{22}M_{\sigma _0}^2)(_{22}M_{f_0}^2)}.$$
(65)
We observe that there are two possible sign assignments. Reality of $`_{12}`$ requires that the following inequality be satisfied
$$\frac{M_{a_0}^2+3M_{\sigma _0}^2}{4}M_{\kappa _0}^2\frac{M_{a_0}^2+3M_{f_0}^2}{4},$$
(66)
As shown in , this inequality holds in a model independent way if symmetry breaking is assumed to be linear. From the expression of the mass matrix (3.1) it is then easy to see that the scalar singlet chiral mass $`M_0^2`$ satisfies a quadratic equation
$$M_0^4BM_0^2+CM_{\eta _0}^2=0,$$
(67)
with
$`B=_{11}{\displaystyle \frac{9}{8}}(2m_K^2+m_\pi ^2)\left(1{\displaystyle \frac{1}{27}}{\displaystyle \frac{M_{\kappa _0}^2M_{a_0}^2}{m_K^2m_\pi ^2}}\right)`$
$`C={\displaystyle \frac{1}{4}}(2m_K^2+m_\pi ^2)\left(1+{\displaystyle \frac{1}{9}}{\displaystyle \frac{M_{\kappa _0}^2M_{a_0}^2}{m_K^2m_\pi ^2}}+{\displaystyle \frac{\sqrt{2}}{3}}{\displaystyle \frac{_{12}}{m_K^2m_\pi ^2}}\right).`$ (68)
Real solutions exist provided
$$B^24CM_{\eta _0}^20,$$
(69)
and it must further be checked that at least one solution is positive. This puts further constraints on the allowed values of the scalar meson masses in this model. Once $`M_0^2`$ is determined, the other parameters are easily evaluated. For instance $`M_8^2`$ is given by,
$`M_8^2={\displaystyle \frac{1}{1+\frac{6C}{M_0^2}}}[M_{a_0}^2+{\displaystyle \frac{2CM_{\eta _0}^2}{M_0^2}}{\displaystyle \frac{3}{4}}(2m_K^2+m_\pi ^2)`$
$`+{\displaystyle \frac{1}{4}}(2m_K^23m_\pi ^2){\displaystyle \frac{M_{\kappa _0}^2M_{a_0}^2}{m_K^2m_\pi ^2}}].`$ (70)
Next, using as experimental input $`M_{\kappa _0}^2M_{a_0}^2`$ and $`_{12}`$ together with eqs.(3.1), (3.1) and (3.2)) gives two linear equations for the three quantities $`c_m/v`$, $`\stackrel{~}{c}_m/v`$ and $`\gamma _3v`$. We will use as an additional constraint, the ratio $`F_K/F_\pi `$, which determines $`c_m/v`$ from the relation,
$$\frac{c_m}{v}=\frac{M_8^2}{2F_\pi }\frac{F_KF_\pi }{m_K^2m_\pi ^2}.$$
(71)
Finally, there remains to determine $`v=F_0/\sqrt{2}`$ . To linear order in the quark masses, $`v`$ is given by the following expression
$$v=\frac{F_\pi }{\sqrt{2}}\left[1\frac{2}{3}m_\pi ^2\left(\frac{2c_m/v}{M_8^2}+\frac{\sqrt{3}\stackrel{~}{c}_m/v}{M_0^2}\right)\frac{4}{3}m_K^2\left(\frac{c_m/v}{M_8^2}+\frac{\sqrt{3}\stackrel{~}{c}_m/v}{M_0^2}\right)\right].$$
(72)
This expression can be recovered in two different ways: one can either use CHPT together with the expressions (3.1) for $`L_4`$ and $`L_5`$ or write down the equations for the vacuum expectation values in the presence of quark masses and solve these equations to linear order in the quark masses.
#### 3.3.1 Light $`\sigma `$ meson
The existence of a very broad scalar resonance in $`\pi \pi `$ scattering with $`M\mathrm{\Gamma }0.50.6`$ GeV is by now well established since the work of Basdevant Froggatt and Petersen, in which the whole set of constraints from unitarity, analyticity, and crossing symmetry has been implemented (see ref. for a recent survey and the particle data book for a complete list of references). What is unclear is whether this state should be interpreted as a physical light scalar resonance. Recently , Black et al. have proposed arguments based on perturbative unitarity favouring the existence of a light $`\sigma `$ and also, of a light $`\kappa `$ meson. According to the inequality (66), if the sigma meson is light then one must have $`M_{\kappa _0}M_{a_0}`$ unless our assumption of linearity in the quark masses fails. Recently, Törnqvist (similar fits were also discussed in ref.) has attempted to accommodate in the linear sigma-model a scalar meson multiplet with a light sigma and a heavy $`\kappa `$. Such fits fail to obtain the $`f_0(980)`$ at the correct mass and therefore would not correctly reproduce the $`\pi \pi `$ phase-shifts around 1 GeV, even if one-loop corrections are included.
Let us now assume that a light sigma meson exists, e.g. $`M_{\sigma _0}0.6`$ GeV, and discuss the consequences. We take otherwise $`M_{a_0}=0.983`$ GeV, and $`M_{f_0}=0.980`$ GeV from experiment and $`M_{\kappa _0}=0.9`$ GeV as proposed in ref.. From this input one easily calculates the parameters $`A`$ and $`B`$ in the quadratic equation (67),
$$A=0.024\mathrm{GeV}^2B=0.042\mathrm{GeV}^2.$$
(73)
There results that no real solution to the equation for $`M_0^2`$ exists unless $`M_{\eta _0}0.06`$ GeV, which appears as an absurdly small value. Reality and positivity of $`M_0^2`$ (and $`M_8^2`$) are necessary conditions for the existence of a stable minimum of the potential, with spontaneous chiral symmetry breaking. One must thus conclude that with such input scalar meson masses $`N_F=3`$ chiral symmetry is not spontaneously broken in the linear sigma-model. The same conclusion holds if one picks up larger values for $`M_{\kappa _0}`$: using the set of parameters determined in ref. for $`M_\sigma =0.4`$ and $`M_\sigma =0.6`$ GeV one verifies that the equation for the vacuum expectation value $`v`$ eq.(48) has no real solution $`v0`$. A better situation prevails if we accept smaller values for $`M_{\kappa _0}`$. If we take, for instance, $`M_{\kappa _0}=0.750`$ GeV, which is the smallest value allowed by the inequality(66), then, eq.(67) has real solutions provided $`M_{\eta _0}0.62`$ GeV, which is still reasonable. A problem remains, however, that being close to a situation where chiral symmetry is unbroken for $`N_F=3`$, while it is in general broken for $`N_F=2`$, the expansion in the strange quark mass is likely to be non-converging. Indeed, if we try to compute $`F_0`$ with, say, $`M_{\eta _0}=0.50`$ GeV, using our linear expansion formula (72), we find $`F_028`$ MeV, which is very small compared to $`F_\pi `$, suggesting that the expansion in $`m_s`$ is not perturbative in this case.
#### 3.3.2 Heavy $`\sigma `$ meson
Let us now consider the scenario of a heavy sigma meson, i.e. we identify the sigma meson with the well established resonance $`f_0(980)`$ and we identify the heavier $`I=0`$ member of the nonet with the $`f_0(1500)`$ (see ). With these assignments, we find that a solution of eq.(67) with real $`M_0^2`$ exists in the range
$$M_{\kappa _0}1.37\mathrm{GeV},$$
(74)
taking $`M_{\eta _0}=0.8`$ GeV and the negative sign for $`_{12}`$ (see (65)). This means that we cannot exactly reproduce the experimental mass of the $`K_0^{}(1430)`$ but we can get reasonably close. We will eventually allow $`M_{\kappa _0}`$ to vary somewhat away from the experimental result. Numerical values of the parameters $`M_0`$, $`M_8`$, $`c_m`$, $`\stackrel{~}{c}_m`$, $`\gamma _3`$ as well as $`v`$ determined for several values of $`M_{\kappa _0}`$ are collected in table (3).
We note that the situation with a minimal symmetry breaking lagrangian (i.e. $`\gamma _1=\gamma _2=\gamma _3=0`$ ) is close to the case when $`M_{\kappa _0}=1.20`$ in the table. In this case, one has $`c_m=\sqrt{3}\stackrel{~}{c}_m=v/2`$. The value of $`F_0=\sqrt{2}v`$ remains reasonable close to $`F_\pi `$ except when $`M_{\kappa _0}`$ gets very near to its upper bound. The coupling $`c_m`$ comes out smaller than that determined in ref. (see (35)) while the coupling $`c_d=v`$ comes out larger. The coupling $`c_d`$ controls the decay width $`a_0\pi \eta `$, which tends to be too large in this model (accepting that the experimental value is $`\mathrm{\Gamma }60`$ MeV, which is subject to some debate), the result in this respect improves as $`M_{\kappa _0}`$ gets larger.
Once these couplings are known it is easy to determine the LEC’s with the formulas given above. Numerical values for the LEC’s dominated by the scalar mesons and not suppressed in the large $`N_c`$ limit are shown in table (4). We have also included $`L_7`$, which is assumed to be saturated by the $`\eta _0`$, computed from the expression,
$$L_7=\frac{\stackrel{~}{d}_m^2}{2M_{\eta _0}^2},$$
(75)
where the coupling $`\stackrel{~}{d}_m`$ of the $`\eta _0`$ field to the pseudo-scalar current has the following expression in the present version of the linear sigma-model,
$$\stackrel{~}{d}_m=\frac{F_0}{2\sqrt{6}}\left(1+\frac{2}{v}(c_m\sqrt{3}\stackrel{~}{c}_m)+12\gamma _3v\right).$$
(76)
The fact that $`L_5`$ does not remain constant when $`M_{\kappa _0}`$ increases is a reflection of the effect of non-linearities in $`m_s`$. The values of the LEC’s may be compared with the ones determined at $`O(p^4)`$ from experimental data
$$10^3L_5^r(m_\eta )=2.2\pm 0.5,10^3L_7=0.4\pm 0.15,10^3L_8^r(m_\eta )=1.1\pm 0.3.$$
(77)
While the prediction for $`L_7`$ is acceptable, that for $`L_8`$ is too small. This is related to the small size of the meson to source coupling $`c_m`$ predicted by the model. Still, the order of magnitude is correct except, perhaps, for the very last line in the table. The LEC’s were recently re-determined in an $`O(p^6)`$ analysis of the experimental data and this brings the values somewhat down in magnitude: $`10^3L_5^r(m_\eta )=1.45\pm 0.12,10^3L_7=0.31\pm 0.15,10^3L_8^r(m_\eta )=0.68\pm 0.18`$. A resonance-saturated estimate for the $`O(p^6)`$ coupling $`C_{19}`$ was provided in ref.. The value that we obtain is significantly smaller, by a factor $`34`$. One difference in the estimates is that we include the effect of cubic interaction terms, but the numerical influence of these terms turns out to be unimportant. The main reason for our smaller result is, again, the fact that the coupling $`c_m`$ is smaller here.
Let us now return to the OZI suppressed coupling constants $`L_4`$ and $`L_6`$ and their $`O(p^6)`$ counterparts $`C_{16}`$ and $`C_{20}+3C_{21}`$. The predictions from the linear sigma-model are collected in table (5).
The values and signs of $`L_4`$, $`L_6`$ are seen to depend very much on the mass of $`\kappa _0`$. They both approximately vanish when $`M_{\kappa _0}1.24`$ GeV. This is the point of minimal OZI violation. As one pushes the $`\kappa _0`$ to higher masses, as required by experiment, then the pattern of OZI violation is not unlike the one found on the basis of sum rules. Both $`L_4`$ and $`L_6`$ become positive, and the orders of magnitudes seem to be in agreement with the discussion based on sum rules. There is even a relatively good agreement as far as $`L_4`$ and $`C_{16}`$ are concerned.
## 4 Conclusions
In this paper, we have pursued an investigation on the evolution of order parameters of the QCD chiral vacuum as one increases the number of massless flavours, as can be inferred from experimental data. An information on this evolution is contained in the $`O(p^4)`$ coupling constants $`L_4`$ and $`L_6`$ which control how $`F_\pi `$ and $`\overline{u}u`$ respectively vary from the $`SU(2)`$ chiral limit with $`m_s200`$ MeV to the $`SU(3)`$ chiral limit. Estimates for the couplings $`L_4`$ and $`L_6`$ are based on the scalar form factors of the pion and the Kaon which are reconstructed from experimental S-wave scattering data modulo some assumptions. We have investigated the $`O(m_s)`$ corrections to our previous results. These have two origins: a) the normalization of the form factors at the origin must include the $`O(p^4)`$ contributions and b) the chiral expansions (for $`F_\pi `$ and $`\mathrm{\Pi }_6(0)`$ ) must be pursued up to order $`p^6`$.
These corrections seem not to affect in a significant way the results obtained previously. In particular there is no sign that these corrections go in the sense of decreasing the values of $`L_4`$ and $`L_6`$ to make them compatible with the naive large $`N_c`$ expectation. This conclusion, however, only holds provided one makes the assumption that the chiral expansion is reasonable, in other terms, that the $`O(p^6)`$ contribution (say in the expansion of $`\mathrm{\Pi }_6(0)`$) is not larger than the $`O(p^4)`$ contribution. Modulo this assumption, one also obtains from this analysis an estimate of the $`O(p^6)`$ coupling constant $`C_{16}`$ and the combination $`C_{20}+3C_{21}`$. Influence of the values of $`L_4`$, $`L_6`$ on the convergence rate of the chiral expansion for various observables was studied very recently in . Because of such assumptions which enter into the calculation, the error bars on our coupling constant estimates must be considered as educated guesses. Our results are compatible with a tendency towards chiral restoration upon increasing $`N_F`$ from $`N_F=2`$ to $`N_F=3`$ which is significant, particularly so for the quark condensate (in agreement with the arguments of ), for which the decrease by a factor of approximately two found in ref. is confirmed. This behaviour of the quark condensate is in qualitative agreement with the one obtained from the instanton liquid model. It would be also interesting to compare this result with unquenched lattice simulations with $`N_F=3`$. Unfortunately, such simulations are not available yet, but will exist in the near future.
We also discussed some predictions from the linear sigma-model, considered as a simple toy model. One application is to gauge the possible importance of self-couplings of resonances (which are present in the model) in resonance saturation estimates of the $`O(p^6)`$ LEC’s. Furthermore, the model provides a connection between simply the spectrum of the scalars and the evolution of chiral order parameters with $`N_F`$. In order to improve the ability of the model to reproduce experimental masses, the symmetry breaking sector was made more general than usually done, but still linear in the quark mass matrix. We have considered two possibilities for the light scalar nonet:
a)The nonet contains a light $`\sigma `$ and a light $`\kappa `$. In this case, we find that chiral symmetry tends to be restored already at $`N_F=3`$ or, at least, for a value of $`N_F`$ somewhat too small, making the chiral expansion in $`m_s`$ unreliable.
b) The nonet is composed of the resonances $`a_0(980)`$, $`f_0(980)`$, $`K_0^{}(1430)`$ and $`f_0(1500)`$ (which are well established). In this case, we find that we need only a small amount of nonlinear effects to exactly reproduce the $`K_0^{}(1430)`$ mass, and $`SU(3)`$ chiral symmetry is realised in the Goldstone mode. Concerning $`L_4`$ and $`L_6`$, we find that they become different from zero as one increases the mass of the $`K_0^{}(1430)`$ towards its experimental value. The signs and orders of magnitude, then, are compatible with the preceding analysis. The same is true of the related $`O(p^6)`$ constants.
###### Acknowledgments.
The author thanks B. Ananthanarayan for drawing his attention to ref. and Jan Stern for useful remarks. This work is supported in part by funds provided by the U.S. Department of Energy (D.O.E.) under cooperative research agreement #DF-FC02-94ER40818 and by the EURODAPHNE network EEC contract #TMR-ERBFMRX-CT980169.
## Appendix A The correlation function $`\mathrm{\Pi }_6(s)`$ to $`O(p^6)`$
At chiral order $`p^4`$, a simple calculation gives the correlation function $`\mathrm{\Pi }_6(s)`$,
$$\mathrm{\Pi }_6(s)=2\overline{J}_K(s)+\frac{4}{9}\overline{J}_\eta (s)+64L_6^r\frac{1}{16\pi ^2}\left[2L_K+\frac{4}{9}L_\eta +\frac{22}{9}\right],$$
(78)
where $`\overline{J}_P(s)`$ is the one-loop function defined as to vanish for $`s=0`$,
$$\overline{J}_P(s)=\frac{1}{16\pi ^2}\left(\sigma _P\mathrm{log}\frac{\sigma _P1}{\sigma _P+1}+2\right),\sigma _P=\sqrt{1\frac{4M_P^2}{s}}.$$
(79)
In all the formulas $`M_P^2`$ stands not for the physical pseudo-scalar meson masses but for their lowest order chiral expansion,
$$M_\pi ^22mB_0,M_K^2(m+m_s)B_0,M_\eta ^2\frac{(2m+4m_s)B_0}{3}.$$
(80)
The following simplified notation for logarithms was introduced
$$L_P=\mathrm{log}\frac{M_P^2}{\mu ^2}.$$
(81)
At order $`p^6`$, one must compute the one-loop and two-loops diagrams shown in fig.1 and add the tree-level contributions from the $`O(p^6)`$ chiral lagrangian. The calculation is lengthy but not excessively difficult because of the absence of sunset-type diagrams in the present case. The result may be written as a sum of six terms:
$$\mathrm{\Delta }\mathrm{\Pi }_6(s)=\mathrm{\Delta }\mathrm{\Pi }_6^a(s)+\mathrm{\Delta }\mathrm{\Pi }_6^b(s)+\mathrm{\Delta }\mathrm{\Pi }_6^{resc}(s)+A\frac{s}{F_0^2}+B+C.$$
(82)
The first term encodes the $`O(p^4)`$ corrections to the pseudo-scalar meson masses in the one-loop functions,
$$\mathrm{\Delta }\mathrm{\Pi }_6^a(s)=s\left(2\frac{\mathrm{\Delta }M_K^2}{M_K^2}\frac{d\overline{J}_K(s)}{ds}+\frac{4}{9}\frac{\mathrm{\Delta }M_\eta ^2}{M_\eta ^2}\frac{d\overline{J}_\eta (s)}{ds}\right).$$
(83)
These mass corrections are given by the following expressions,
$`\mathrm{\Delta }M_K^2=M_K^2[{\displaystyle \frac{mB_0}{F_0^2}}(32S_68S_8+{\displaystyle \frac{1}{72\pi ^2}}L_\eta )`$
$`+{\displaystyle \frac{m_sB_0}{F_0^2}}(16S_68S_8+{\displaystyle \frac{1}{36\pi ^2}}L_\eta )]`$ (84)
and
$`\mathrm{\Delta }M_\eta ^2=M_\eta ^2[{\displaystyle \frac{mB_0}{F_0^2}}(32S_6{\displaystyle \frac{16}{3}}S_8{\displaystyle \frac{5}{48\pi ^2}}L_\pi +{\displaystyle \frac{1}{9\pi ^2}}L_K{\displaystyle \frac{1}{144\pi ^2}}L_\eta )`$
$`+{\displaystyle \frac{m_sB_0}{F_0^2}}(16S_6{\displaystyle \frac{32}{3}}S_8+{\displaystyle \frac{1}{24\pi ^2}}L_\pi +{\displaystyle \frac{1}{18\pi ^2}}L_K{\displaystyle \frac{1}{18\pi ^2}}L_\eta )]`$
$`+{\displaystyle \frac{(m_sm)^2B_0^2}{F_0^2}}\left[{\displaystyle \frac{128}{9}}S_7+{\displaystyle \frac{1}{108\pi ^2}}L_K{\displaystyle \frac{1}{18\pi ^2}}L_\pi \right].`$ (85)
In these expressions, $`S_6`$, $`S_7`$ and $`S_8`$ denote the following combinations of $`O(p^4)`$ coupling constants,
$$S_6=2L_6^r+L_4^r,S_7=3L_7^r+L_8^r,S_8=2L_8^r+L_5^r.$$
(86)
The second corrective term $`\mathrm{\Delta }\mathrm{\Pi }_6^b(s)`$ is generated by $`O(p^4)`$ corrections to the scalar form factors and has the following form,
$`\mathrm{\Delta }\mathrm{\Pi }_6^b(s)=3\overline{J}_\pi (s)\left[{\displaystyle \frac{d\mathrm{\Delta }M_\pi ^2}{dm_sB_0}}+{\displaystyle \frac{s}{F_0^2}}\left(8L_4^r{\displaystyle \frac{1}{32\pi ^2}}(L_K+1)\right)\right]`$
$`+\overline{J}_K(s)[2{\displaystyle \frac{d\mathrm{\Delta }M_K^2}{dmB_0}}+2{\displaystyle \frac{d\mathrm{\Delta }M_K^2}{dm_sB_0}}+{\displaystyle \frac{s}{F_0^2}}(16L_5^r+48L_4^r`$
$`{\displaystyle \frac{1}{32\pi ^2}}(3L_\pi +6L_K+3L_\eta +12))]+{\displaystyle \frac{1}{3}}\overline{J}_\eta (s)[2{\displaystyle \frac{d\mathrm{\Delta }M_\eta ^2}{dmB_0}}+{\displaystyle \frac{d\mathrm{\Delta }M_\eta ^2}{dm_sB_0}}`$
$`+{\displaystyle \frac{s}{F_0^2}}({\displaystyle \frac{32}{3}}L_5^r+40L_4^r{\displaystyle \frac{9}{32\pi ^2}}(L_K+1))].`$ (87)
The derivatives of $`\mathrm{\Delta }M_K^2`$ and $`\mathrm{\Delta }M_\eta ^2`$ are easily computed from the formulas above and for $`\mathrm{\Delta }M_\pi ^2`$ one has,
$$\frac{d\mathrm{\Delta }M_\pi ^2}{m_sB_0}=mB_0\left(32S_6\frac{1}{36\pi ^2}(L_\eta +1)\right).$$
(88)
The third corrective term $`\mathrm{\Delta }\mathrm{\Pi }_6^{resc}(s)`$ is the rescattering contribution,
$`\mathrm{\Delta }\mathrm{\Pi }_6^{resc}(s)={\displaystyle \frac{3s}{2F_0^2}}\overline{J}_K(s)\left(\overline{J}_\pi (s)+\overline{J}_K(s)+\overline{J}_\eta (s)\right)`$ (89)
$`+{\displaystyle \frac{4}{3}}\overline{J}_\eta (s)\left({\displaystyle \frac{(m+8m_s)B_0}{27F_0^2}}\overline{J}_\eta (s){\displaystyle \frac{(m+m_s)B_0}{F_0^2}}\overline{J}_K(s)+{\displaystyle \frac{mB_0}{F_0^2}}\overline{J}_\pi (s)\right).`$
The last corrective term is a polynomial, linear in $`s`$. This term picks up contributions from four coupling constants of the $`O(p^6)`$ chiral lagrangian, $`C_{39}^r`$, $`C_{20}^r`$, $`C_{21}^r`$ and $`C_{94}^r`$ in the notation of . The part proportional to $`s`$ is as follows
$`A=64C_{39}^r+{\displaystyle \frac{1}{\pi ^2}}[{\displaystyle \frac{3}{512\pi ^2}}L_K(L_\pi +L_K+L_\eta )+L_\pi ({\displaystyle \frac{3}{512\pi ^2}}{\displaystyle \frac{3}{2}}L_4^r)`$
$`+L_K\left({\displaystyle \frac{3}{128\pi ^2}}3L_4^rL_5^r\right)+L_\eta \left({\displaystyle \frac{3}{512\pi ^2}}{\displaystyle \frac{5}{6}}L_4^r{\displaystyle \frac{2}{9}}L_5^r\right)`$
$`+{\displaystyle \frac{9}{512\pi ^2}}{\displaystyle \frac{16}{3}}L_4^r{\displaystyle \frac{11}{9}}L_5^r]`$ (90)
The constant terms, finally are
$`B={\displaystyle \frac{mB_0}{F_0^2}}[32(8C_{20}^r+48C_{21}^r+C_{94}){\displaystyle \frac{1}{864\pi ^4}}L_\eta (9L_\pi +12L_KL_\eta )`$
$`+{\displaystyle \frac{L_\pi }{\pi ^2}}\left(12S_6+{\displaystyle \frac{1}{96\pi ^2}}\right)+{\displaystyle \frac{L_K}{\pi ^2}}\left(20S_6+6S_8{\displaystyle \frac{1}{96\pi ^2}}\right)`$
$`+{\displaystyle \frac{L_\eta }{\pi ^2}}\left(4S_6+{\displaystyle \frac{8}{9}}S_8{\displaystyle \frac{11}{5184\pi ^2}}\right)`$
$`+{\displaystyle \frac{1}{\pi ^2}}({\displaystyle \frac{206}{9}}S_6{\displaystyle \frac{16}{27}}S_7+{\displaystyle \frac{124}{27}}S_8+{\displaystyle \frac{1}{5184\pi ^2}})]`$
$`+{\displaystyle \frac{m_sB_0}{F_0^2}}[256(C_{20}^r+3C_{21}^r){\displaystyle \frac{1}{288\pi ^4}}L_\eta (L_\eta +5L_K)`$
$`+{\displaystyle \frac{L_K}{\pi ^2}}\left(16S_6+6S_8{\displaystyle \frac{1}{96\pi ^2}}\right)+{\displaystyle \frac{L_\eta }{9\pi ^2}}\left(48S_6+16S_7+16S_8{\displaystyle \frac{37}{576\pi ^2}}\right)`$
$`+{\displaystyle \frac{1}{\pi ^2}}({\displaystyle \frac{118}{9}}S_6+{\displaystyle \frac{16}{27}}S_7+{\displaystyle \frac{140}{27}}S_8{\displaystyle \frac{19}{5184\pi ^2}})],`$ (91)
and
$$C=\frac{1}{16\pi ^2}\left(2\frac{\mathrm{\Delta }M_K^2}{M_K^2}+\frac{4}{9}\frac{\mathrm{\Delta }M_\eta ^2}{M_\eta ^2}\right)$$
(92)
The terms proportional to the mass corrections $`\mathrm{\Delta }M_K^2`$ and $`\mathrm{\Delta }M_\eta ^2`$ (see eq.(83) and (92)) may eventually be absorbed into the $`O(p^4)`$ expression for $`\mathrm{\Pi }_6(s)`$, which amounts to replace the lowest order expression for the pseudo-scalar masses there by their expression up to $`O(p^4)`$. Formula (2.2) in the text is obtained from $`B+C`$, setting $`m=0`$.
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# Wave propagation in linear electrodynamics
## I Introduction
Electromagnetic wave represents perhaps the most important classical device with the help of which one can carry out physical measurements and transmit information. Intrinsic properties and motion of material media, as well as the geometrical structure of spacetime, can affect the propagation of electromagnetic waves. In the most general setting , electromagnetic phenomena are described by the pair of two-forms $`H,F`$ (called the electromagnetic excitation and the field strength, respectively) which satisfy the Maxwell equations $`dH=J,dF=0`$, together with the constitutive law $`H=H(F)`$. The latter relation contains a crucial information about the underlying physical continuum (i.e., material medium and/or spacetime). Mathematically, this constitutive law arises either from a suitable phenomenological theory of a medium or from the electromagnetic field Lagrangian.
In general, the constitutive law establishes a nonlinear (or even nonlocal) relation between the electromagnetic excitation and the field strength. The function (or functional) $`H(F)`$ may depend on the polarization and magnetization properties of matter, and/or on the spacetime geometry, i.e. metric, curvature, torsion and nonmetricity. Previously, propagation of electromagnetic waves was analysed for a variety of constitutive laws: for nonlinear models in Minkowski and Riemannian spacetimes , for electrodynamics in a Riemann-Cartan manifold , and also for certain nonminimal and higher derivative gravity models . Numerous authors discussed the electromagnetic waves in the Einstein-Maxwell theory. The main aim of this paper is to investigate the wave propagation in the Maxwell electrodynamics with the most general linear constitutive law. We derive the generalized Fresnel equation which determines the wave normals directly from the constitutive coefficients. This result is of interest, e.g., for various applications in crystaloptics and related domains.
Another motivation for the present work came from the study of a deep relationship between the duality operators defined on two-forms and the conformal classes of spacetime metrics in four dimensions. Within the classical Maxwell electrodynamics, Toupin, Schönberg and others have noticed that the constitutive coefficients define a duality operator, provided a certain reciprocity or closure condition is fulfilled, and gave first demonstrations of the existence of the corresponding conformal metric structure. Later these observations were rediscovered and developed in mathematics and in the gravity theory . Recently the complete explicit solution of the closure relation has been given , and it was conjectured that the reciprocity condition is a necessary and sufficient condition for the standard null-cone structure for the light propagation (see also independent arguments in ). Here we give a partial answer to this question.
## II Electrodynamics with linear constitutive law
Let us consider the Maxwell equations in vacuum,
$$dH=0,dF=0,$$
(1)
i.e. we assume that the electric current three-form $`J`$ vanishes in the spacetime region under consideration. Given the local coordinates $`x^i`$, $`i=0,1,2,3`$, we can decompose the exterior forms as
$$H=\frac{1}{2}H_{ij}dx^idx^j,F=\frac{1}{2}F_{ij}dx^idx^j.$$
(2)
Following , we write the linear constitutive law in terms of the electromagnetic excitation and field strength tensors as
$$H_{ij}=\frac{1}{4}ϵ_{ijkl}\chi ^{klmn}F_{mn},i,j,\mathrm{}=0,1,2,3.$$
(3)
Here $`ϵ_{ijkl}`$ is the Levi-Civita symbol and $`\chi ^{ijkl}(x)`$ an even tensor density of weight $`+1`$ (called the constitutive tensor density) which can be decomposed according to
$$\chi ^{ijkl}=f(x)\stackrel{\mathrm{o}}{\chi }{}_{}{}^{ijkl}+\alpha (x)ϵ^{ijkl},\mathrm{with}\stackrel{\mathrm{o}}{\chi }{}_{}{}^{[ijkl]}0.$$
(4)
Here $`f(x)`$ is a dimensionfull scalar function such that $`\stackrel{\mathrm{o}}{\chi }^{ijkl}`$ is dimensionless. The pseudo-scalar constitutive function $`\alpha (x)`$ can be identified (on the kinematic level) as an Abelian axion field, whereas $`f(x)`$ can be interpreted as a dilaton scalar field. Note that $`\stackrel{\mathrm{o}}{\chi }^{ijkl}`$ has the same algebraic symmetries and therefore the same number of 20 independent components as a Riemannian curvature tensor:
$$\stackrel{\mathrm{o}}{\chi }{}_{}{}^{ijkl}=\stackrel{\mathrm{o}}{\chi }{}_{}{}^{jikl}=\stackrel{\mathrm{o}}{\chi }{}_{}{}^{ijlk}=\stackrel{\mathrm{o}}{\chi }{}_{}{}^{klij},\stackrel{\mathrm{o}}{\chi }{}_{}{}^{[ijkl]}=0.$$
(5)
This follows from the existence and the structure of the Lagrangian for the linear electrodynamics $`V_{\mathrm{lin}}=\frac{1}{2}HF`$, see . It is convenient to adopt a more compact (essentially bivector) notation by defining the 3-(co)vector quantities
$$𝒟^a:=\left(\begin{array}{c}H_{23}\\ H_{31}\\ H_{12}\end{array}\right),_a:=\left(\begin{array}{c}H_{01}\\ H_{02}\\ H_{03}\end{array}\right),\mathrm{and}B^a:=\left(\begin{array}{c}F_{23}\\ F_{31}\\ F_{12}\end{array}\right),E_a:=\left(\begin{array}{c}F_{10}\\ F_{20}\\ F_{30}\end{array}\right),$$
(6)
for the electric and magnetic excitations, and for the magnetic and electric field strengths, respectively. The Latin indices label now $`a,b,c,\mathrm{}=1,2,3`$. The constitutive tensor is then naturally parametrized by a triplet of $`3\times 3`$ matrices, $`\stackrel{\mathrm{o}}{\chi }{}_{}{}^{ijkl}=\{𝒜^{ab},_{ab},𝒞^a{}_{b}{}^{}\}`$, so that the constitutive law (4) is finally recasted into
$$\left(\begin{array}{c}_a\\ 𝒟^a\end{array}\right)=f(x)\left(\begin{array}{cc}𝒞^b_a& _{ab}\\ 𝒜^{ab}& 𝒞^a_b\end{array}\right)\left(\begin{array}{c}E_b\\ B^b\end{array}\right)+\alpha (x)\left(\begin{array}{c}E_a\\ B^a\end{array}\right).$$
(7)
Here the $`3\times 3`$ matrices satisfy $`𝒜^{ab}=𝒜^{ba}`$, $`_{ab}=_{ba}`$, and $`𝒞^a{}_{a}{}^{}=0`$, thereby providing the algebraic properties (5).
## III Wave propagation: Fresnel equation
In the theory of partial differential equations, the propagation of waves is described by Hadamard discontinuities of solutions across a characteristic (wave front) hypersurface $`S`$ . One can locally define $`S`$ by the equation $`\mathrm{\Phi }(x^i)=const`$. The Hadamard discontinuity of any function $`(x)`$ across the hypersurface $`S`$ is determined as the difference $`\left[\right](x):=(x_+)(x_{})`$, where $`x_\pm :=\underset{\epsilon 0}{lim}(x\pm \epsilon )`$ are points on the opposite sides of $`Sx`$. An ordinary electromagnetic wave is a solution of the Maxwell equations (1) for which the derivatives of $`H`$ and $`F`$ have regular discontinuities across the wave front hypersurface $`S`$.
In terms of the (co)vector components, we have on the characteristic hypersurface $`S`$:
$`[𝒟^a]=0,[_i𝒟^a]=d^aq_i,[_a]=0,[_i_a]=h_aq_i,`$ (8)
$`[B^a]=0,[_iB^a]=b^aq_i,[E_a]=0,[_iE_a]=e_aq_i,`$ (9)
where $`d^a,h_a,b^a,e_a`$ describe discontinuities of the corresponding quantities across $`S`$, and the wave-covector normal to the front is given by
$$q_i:=_i\mathrm{\Phi }.$$
(10)
Equations (8)-(9) represent the Hadamard geometrical compatibility conditions. Substituting (2) into (1), and using (6) and (8)-(9), we find
$`q_0d^aϵ^{abc}q_bh_c=0,q_0b^a+ϵ^{abc}q_be_c=0,`$ (11)
$`q_ad^a=0,q_ab^a=0,`$ (12)
where $`ϵ^{abc}`$ is the 3-dimensional Levi-Civita symbol. In this system only the 6 equations (11) are independent. Assuming that $`q_00`$, one finds that the equations (12) are trivially satisfied if one substitutes (11) into them. \[Note that the characteristics with $`q_0=0`$ do not have intrinsic meaning for the evolution equations, since they obviously depend on the arbitrary choice of coordinates\].
Differentiating (7) and using the compatibility conditions (8)-(9), we find additionally 6 algebraic equations:
$$\left(\begin{array}{c}h_a\\ d^a\end{array}\right)=f(x)\left(\begin{array}{cc}𝒞^b_a& _{ab}\\ 𝒜^{ab}& 𝒞^a_b\end{array}\right)\left(\begin{array}{c}e_b\\ b^b\end{array}\right)+\alpha (x)\left(\begin{array}{c}e_a\\ b^a\end{array}\right).$$
(13)
Note that the constitutive coefficients and their first derivatives are assumed to be continuous across $`S`$.
We can now substitute $`d^a`$ and $`h_a`$ from (13) into the first equation (11), which gives
$$f(x)q_0\left(𝒜^{ab}e_b+𝒞^a{}_{b}{}^{}b_{}^{b}\right)+\alpha (x)q_0b^a=f(x)ϵ^{abc}q_b\left(𝒞^d{}_{c}{}^{}e_{d}^{}+_{cd}b^d\right)\alpha (x)ϵ^{abc}q_be_c.$$
(14)
The terms proportional to the axion field $`\alpha (x)`$ drop out completely due to (11), and then one can also remove the common dilaton factor $`f(x)`$ on both sides of the equation. \[We assume $`f(x)0`$, since otherwise there is no hyperbolic evolution system\]. It remains finally to substitute $`b^a`$ in terms of $`e_b`$ from the second equation (11), and after some rearrangements one finds:
$$\left(q_0^2𝒜^{ab}+q_0q_d\left[𝒞^a{}_{c}{}^{}ϵ_{}^{cdb}+𝒞^b{}_{c}{}^{}ϵ_{}^{cda}\right]+q_eq_fϵ^{aec}ϵ^{bfd}_{cd}\right)e_b=0.$$
(15)
This homogeneous algebraic equation has a nontrivial solution when
$$𝒲:=det\left|q_0^2𝒜^{ab}+q_0q_d\left[𝒞^a{}_{c}{}^{}ϵ_{}^{cdb}+𝒞^b{}_{c}{}^{}ϵ_{}^{cda}\right]+q_eq_fϵ^{aec}ϵ^{bfd}_{cd}\right|=0.$$
(16)
This is a Fresnel equation which is central in the wave propagation analysis. It determines the geometry of the wave normals in terms of the constitutive coefficients $`𝒜,,𝒞`$. A direct calculation yields the general result:
$$𝒲=q_0^2\left(q_0^4M+q_0^3q_aM^a+q_0^2q_aq_bM^{ab}+q_0q_aq_bq_cM^{abc}+q_aq_bq_cq_dM^{abcd}\right)=0,$$
(17)
where we have denoted
$`M:=det𝒜,M^a:=2ϵ_{bcd}𝒜^{ab}𝒞^c{}_{e}{}^{}𝒜_{}^{ed},`$ (18)
$`M^{ab}:=_{cd}(𝒜^{ab}𝒜^{cd}𝒜^{ac}𝒜^{bd})𝒜^{cd}𝒞^a{}_{c}{}^{}𝒞_{}^{b}{}_{d}{}^{}+4𝒜^{ac}𝒞^b{}_{d}{}^{}𝒞_{}^{d}{}_{c}{}^{}2𝒜^{ab}𝒞^c{}_{d}{}^{}𝒞_{}^{d}{}_{c}{}^{},`$ (19)
$`M^{abc}:=2ϵ^{cde}[_{df}(𝒜^{ab}𝒞^f{}_{e}{}^{}𝒜^{af}𝒞^b{}_{e}{}^{})+𝒞^a{}_{e}{}^{}𝒞_{}^{b}{}_{f}{}^{}𝒞_{}^{f}{}_{d}{}^{}]`$ (20)
$`M^{abcd}:=ϵ^{cef}ϵ^{dgh}_{fh}[{\scriptscriptstyle \frac{1}{2}}𝒜^{ab}_{eg}𝒞^a{}_{e}{}^{}𝒞_{}^{b}{}_{g}{}^{}].`$ (21)
Note that only the completely symmetric parts $`M^{(a_1\mathrm{}a_p)}`$, $`p=2,3,4`$, contribute to the Fresnel equation. Since $`q_00`$, one can delete the first factor in (17), and thus we finally find that the wave covector $`q_i`$ lies, in general, on a 4th order surface. This is different from the light cone (i.e., 2nd order) structure which arises only in a particular case. In the next section we demonstrate that the latter corresponds to the closure condition. Earlier, the relation between the fourth- and the second-order wave geometry was studied by Tamm for a special case of the linear constitutive law.
## IV The closure relation as a sufficient condition
The linear constitutive law defines a duality operator when the constitutive coefficients satisfy the ‘reciprocity’ or ‘closure’ relation :
$$\frac{1}{4}ϵ_{ijmn}ϵ_{pqrs}\stackrel{\mathrm{o}}{\chi }{}_{}{}^{mnpq}\stackrel{\mathrm{o}}{\chi }{}_{}{}^{rskl}=\delta _{ij}^{kl},$$
(22)
or in terms of the $`3\times 3`$ matrices:
$$𝒜^{ac}_{cb}+𝒞^a{}_{c}{}^{}𝒞_{}^{c}{}_{b}{}^{}=\delta _b^a,𝒞^{(a}{}_{c}{}^{}𝒜_{}^{b)c}=0,𝒞^c{}_{(a}{}^{}_{b)c}^{}=0.$$
(23)
The general solution of the closure condition (22)-(23) reads :
$`𝒜^{ab}`$ $`=`$ $`{\displaystyle \frac{1}{det}}(k^2^{ab}k^ak^b)^{ab},`$ (24)
$`𝒞^a_b`$ $`=`$ $`^{ad}ϵ_{dbc}k^c={\displaystyle \frac{1}{det}}ϵ^{adc}_{db}k_c.`$ (25)
Here $`k^a`$ is an arbitrary 3-vector, $`k_b:=_{ab}k^a`$, $`k^2:=_{ab}k^ak^b`$, and $`^{ab}`$ denotes the inverse matrix to $`_{ab}`$.
Starting from (24)-(25), the direct calculation yields:
$`M={\displaystyle \frac{1}{det}}\left(1{\displaystyle \frac{k^2}{det}}\right)^2,`$ (26)
$`M^a={\displaystyle \frac{1}{det}}\mathrm{\hspace{0.17em}4}k^a\left(1{\displaystyle \frac{k^2}{det}}\right),`$ (27)
$`M^{ab}={\displaystyle \frac{1}{det}}\mathrm{\hspace{0.17em}4}k^ak^b+2^{ab}\left(1{\displaystyle \frac{k^2}{det}}\right),`$ (28)
$`M^{abc}=\mathrm{\hspace{0.17em}4}^{b(a}k^{c)},`$ (29)
$`M^{(abcd)}=(det)^{(ab}^{cd)}.`$ (30)
Substituting all this into the general Fresnel equation (17), we find
$`𝒲`$ $`=`$ $`\sigma q_0^2\left[{\displaystyle \frac{q_0^2}{\sqrt{|det}|}}\left(1{\displaystyle \frac{k^2}{det}}\right){\displaystyle \frac{2q_0(q_ak^a)}{\sqrt{|det}|}}\sqrt{|det|}(q_aq_b^{ab})\right]^2`$ (31)
$`=`$ $`\sigma q_0^2\left(q_iq_jg^{ij}\right)^2.`$ (32)
Here $`\sigma =sign(det)`$, and $`g^{ij}`$ is the (inverse) 4-dimensional metric which arises from the duality operator and the closure relation :
$`g^{00}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{|det}|}}\left(1{\displaystyle \frac{k^2}{det}}\right),`$ (33)
$`g^{0a}`$ $`=`$ $`{\displaystyle \frac{k^a}{\sqrt{|det|}}},`$ (34)
$`g^{ab}`$ $`=`$ $`\sqrt{|det|}^{ab}.`$ (35)
This metric $`g_{ij}`$ (defined up to a conformal factor) always has the Lorentzian signature, although it is not necessarily interpretable as a spacetime metric (this is a so called optical metric, in general; see, e.g., ). As shown in , the constitutive tensor density (4) can be rewritten in terms of this metric as
$$\chi ^{ijkl}=f(x)\sqrt{g}\left(g^{ik}g^{jl}g^{jk}g^{il}\right)+\alpha (x)ϵ^{ijkl}.$$
(36)
Thus we indeed recover the null cone $`q_iq^i=q_iq_jg^{ij}=0`$ structure for the propagation of electromagnetic waves from our general analysis: provided the constitutive matrices satisfy the closure relation (22)-(23), the quartic surface (17) degenerates to the null cone for the induced metric $`g_{ij}`$.
It is worthwhile to note that the Fresnel equation (17) can be rewritten in an explicitly covariant form
$$G^{ijkl}q_iq_jq_kq_l=0,i,j,\mathrm{}=0,1,2,3,$$
(37)
where the fourth order totally symmetric tensor density $`G^{ijkl}`$ is constructed as the cubic polynomial of the components of the constitutive tensor:
$$G^{ijkl}:=\frac{1}{4!}\stackrel{\mathrm{o}}{\chi }{}_{}{}^{mnp(i}\stackrel{\mathrm{o}}{\chi }{}_{}{}^{j|qr|k}\stackrel{\mathrm{o}}{\chi }{}_{}{}^{l)stu}ϵ_{mnrs}^{}ϵ_{pqtu}.$$
(38)
\[Here the total symmetrization is extended only over the four indices $`i,j,k,l`$ with all the summation indices excluded\]. Tamm has introduced a similar ‘fourth-order metric’ for the particular case of the linear constitutive law.
## V The closure relation as a necessary condition
It was conjectured that the closure relation is not only sufficient, but also a necessary condition for the reduction of the quartic geometry (17) to the null cone. The complete proof of this conjecture requires a rather lengthy algebra and will be considered elsewhere. Here we demonstrate the validity of the necessary condition in a particular case when the matrix $`𝒞=0`$.
Putting $`𝒞^a{}_{b}{}^{}=0`$, we find from (18)-(21) that $`M^a=0`$ and $`M^{abc}=0`$, whereas
$`M^{ab}`$ $`=`$ $`_{cd}(𝒜^{ab}𝒜^{cd}𝒜^{ac}𝒜^{bd}),`$ (39)
$`M^{(abcd)}`$ $`=`$ $`(det)𝒜^{(ab}^{cd)}.`$ (40)
Consequently, (17) reduces to
$$𝒲=q_0^2\left(det𝒜q_0^4+q_0^2\gamma +det\alpha \beta \right),$$
(41)
where $`\alpha :=A^{ab}q_aq_b`$, $`\beta :=B^{ab}q_aq_b`$, and $`\gamma :=M^{ab}q_aq_b`$. Assuming that the last equation describes a null cone, one concludes that the roots for $`q_0^2`$ should coincide and thus necessarily
$$\gamma ^2=4det𝒜det\alpha \beta .$$
(42)
Let us write $`(det𝒜det)=s|det𝒜det|`$, with $`s=sign(det𝒜det)`$. Then (42) yields
$$2\sqrt{|det𝒜det|}\frac{\alpha }{\gamma }=s\lambda ,2\sqrt{|det𝒜det|}\frac{\beta }{\gamma }=\frac{1}{\lambda },$$
(43)
where $`\lambda `$ is an arbitrary scalar factor. Recalling the definitions of $`\alpha ,\beta ,\gamma `$, we then find
$$𝒜^{ab}=s\lambda ^2^{ab}.$$
(44)
Consequently, $`M=det𝒜=s\lambda ^6/det`$ and $`M^{ab}=2\lambda ^4^{ab}`$, and therefore one verifies that
$$𝒲=\frac{s\lambda ^2q_0^2}{det}\left(\lambda ^2q_0^2+sq_aq_b^{ab}det\right)^2.$$
(45)
We immediately see that for $`s=1`$ the quadratic form in (45) can have either the $`(+)`$ signature, or $`(+++)`$. Similarly, for $`s=1`$ the signature is either $`(++++)`$, or $`(++)`$. Therefore, the Fresnel equation describes a correct light cone (hyperbolic) structure only in the case $`s=1`$.
Finally, one can verify that the above solutions satisfy
$$\frac{1}{4}ϵ_{ijmn}ϵ_{pqrs}\stackrel{\mathrm{o}}{\chi }{}_{}{}^{mnpq}\stackrel{\mathrm{o}}{\chi }{}_{}{}^{rskl}=s\lambda ^2\delta _{ij}^{kl},$$
(46)
which for $`s=1`$ reproduces the closure relation (22) after a trivial rescaling of the constitutive tensor density (and subsequently absorbing the factor $`\lambda `$ into the ‘dilaton’ field $`f`$).
## VI Conclusions
In this paper we have derived, extending the earlier results (see e.g., ), the Fresnel equation governing the propagation of electromagnetic waves for the most general linear constitutive law. The wave covector lies, in general, on a fourth order surface. Such generic fourth order structure is not affected by the axion- and dilaton-like parts of the constitutive tensor. Note however that the linear constitutive law $`H=\alpha (x)F`$ does not lead to hyperbolic evolution equations, and hence necessarily $`f(x)0`$.
We have proved that the closure relation (22) is a sufficient condition for the reduction of the fourth order surface to the familiar second order light cone structure. The corresponding family of conformally related metrics $`g`$ coincides with that derived in , see also . This result may be considered as an alternative (as compared to Urbantke’s scheme ) derivation of the Lorentzian metric $`g`$ from a duality operator. In terms of the Lagrangian, the closure relation is equivalent to the statement that $`V_{\mathrm{lin}}=\frac{1}{2}\left[f(x)F{}_{}{}^{}F+\alpha (x)FF\right]`$, where the Hodge operator is defined by the metric $`g`$.
For the special case $`𝒞^a{}_{b}{}^{}=0`$ we have proved that the requirement of reduction of the fourth order Fresnel structure to a second order one implies a relation between the constitutive coefficients which is slightly weaker than the closure relation (22), in that it allows for an arbitrary scalar factor. The latter though can be removed by the redefinition of the dilaton field $`f(x)`$. Also the signature of the resulting quadratic form is not fixed, so that one has to impose hyperbolicity as a separate condition.
It is worthwhile to note that the results obtained can be directly applied to the refinement and generalization of the previous analyses of the observational tests of the equivalence principle. See, for instance, , where some particular cases of the Fresnel equation have been studied in this context.
Acknowledgments
We are grateful to Friedrich W. Hehl for useful comments and discussion of the results obtained. TF thanks the Institute for Theoretical Physics, University of Cologne, for the warm hospitality. GFR would like to thank the German Academic Exchange Service (DAAD) for a graduate fellowship (Kennziffer A/98/00829).
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# Acknowledgments
## Acknowledgments
The author would like to thank Y. Nomura, T. Watari and T. Yanagida for valuable discussions.
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# 1 Introduction
## 1 Introduction
In the past years many aspects of domain wall (DW) solutions of 5-dimensional supergravity have been discussed and one of the most interesting application concerns a supergravity description of the infra-red (IR) physics of 4-dimensional field theories. In particular it is possible to obtain a supergravity description of the renormalization group (RG) flow towards non-trivial IR fixed points, which can be conformal or non-conformal. Most progress has been made for non-conformal IR fixed points , where the scalars flow towards a singularity in the superpotential. Many superpotentials coming from string theory compactifications have poles, which are IR attractive and imply a singularity in the supergravity solution . In special cases the singularity indicates the appearance of a Coulomb branch or can be resolved as discussed in . On the other hand less is known about conformal IR fixed points. Especially for supersymmetric RG-flows the construction of the corresponding supergravity solutions prove difficult and the only known example that we are aware of has been discussed in . Moreover, potentials with IR-attractive fixed points are essential for a string- or M-theory embedding of the Randall-Sundrum (RS) scenario . Only if such potentials exist a thin-wall approximation is justified and we can approximate the scalars as constants given by their fixed-point values.
In this letter we discuss domain wall solutions, where the superpotential in the simplest case depend on a single scalar $`\theta `$ and has the form
$$W=2(a+b\mathrm{cos}\theta ).$$
(1)
Potentials of this type are generated naturally by an instanton/monopole condensation and have an old history in gauge theory in the discussion of confinement and for domain walls ; for a recent discussion see also .
As we show in the next section, this superpotential can be obtained from a model where the scalar fields parameterize the coset $`SU(2,1)/U(2)`$ and a linear combination of both Abelian isometries is gauged. Depending on the constant parameters $`a`$ and $`b`$ it yields ultra-violet (UV) as well as IR attractive fixed points and the extrema of $`W`$ give a (negative) cosmological constant yielding an anti deSitter spacetime. In the special case of $`a=\pm b`$, the cosmological constant vanishes at one extrema and we obtain flat spacetime.
In section three we derive explicit domain wall solutions interpolating between the (different) extrema of $`W`$. Depending the choice for $`a`$ and $`b`$, these solutions are not only interesting from the RG-flow point of view, but provide also a realization of the RS scenario; the no-go statements do not apply for this model.
Embedding this model into $`N`$=2, $`D`$=5 gauged supergravity , the scalars parameterizing the $`SU(2,1)/U(2)`$ coset build up the universal hypermultiplet. As we argue in the last section, from the M-theory perspective the potential can be understood by a superposition of non-trivial G-fluxes and an M5-brane instanton gas.
## 2 Abelian gauging of the $`SU(2,1)/U(2)`$ coset
In 5-d supergravity obtained by string or M-theory compactification, the scalars parameterize a coset space and the only potential consistent with supersymmetry comes from gauging of isometries of this coset. Since we are interested in flat domain walls, we can omit the gauge fields and the bosonic part of the action reads
$$S=d^5x\sqrt{g}\left[\frac{1}{2}RV(\mathrm{\Phi })\frac{1}{2}g_{MN}_\mu \mathrm{\Phi }^M^\mu \mathrm{\Phi }^N\right]$$
(2)
where the scalars $`\mathrm{\Phi }^M`$ are real and parameterize a space with the metric $`g_{MN}`$. Following general arguments , the potential can be expressed in terms of a superpotential $`W`$ as
$$V=6\left(\frac{3}{4}g^{MN}_MW_NWW^2\right).$$
(3)
Before we can discuss BPS domain wall solutions, we have to derive the superpotential. We consider the coset $`SU(2,1)/U(2)`$ and the metric can be derived from the Kähler potential (in more general models will be discussed)
$$K=\mathrm{log}(1|z_1|^2|z_2|^2),|z_1|^2+|z_2|^2<1.$$
(4)
There is another commonly used parameterization of this coset, which makes the quaternionic nature of this manifold manifest and corresponds to the Kähler potential $`K=\mathrm{log}(\frac{S+\overline{S}}{2}C\overline{C})`$ , where the complex scalars $`S,C`$ are known to enter the universal hypermultiplet of $`N=2`$ supergravity. In the ungauged case both parameterizations are equivalent, but after gauging the resulting superpotentials differ significantly<sup>2</sup><sup>2</sup>2A similar effect is also known for $`N=4`$, $`D=4`$ supergravity, where the $`SO(4)`$ and $`SU(2)\times SU(2)`$ formulation are duality equivalent in the ungauged case, but differ after gauging. The first case has an AdS vacuum, the other not, see . Note, gauging does not preserve the duality symmetry.; in the quaternionic formulation we could not find a superpotential with at least two extrema. This may be related to the fact, that it is non-trivial to introduce global quaternionic coordinates on a curved space, which would be necessary for the discussion of domain walls; but this issue deserves further clarification . We do not want to go into details here about the relation of the two models and will instead come back to our parameterization given by the Kähler potential (4). So, we will treat this 4-d scalar manifold not as a quaternionic but as a special Kähler space and follow the formalism developed in the literatur . A disadvantage is however, that it is not clear whether this procedure yields necessarily supersymmetric solution with four unbroken supercharges (because in the case at hand $`z_1`$ and $`z_2`$ are in the same multiplet). A consistency check is, that the scaling dimensions coming from the sugra scalars fit into known representations in the dual field theory. We come back to this point below.
The two phase rotations of $`z_1`$ and $`z_2`$ are two Abelian isometries corresponding to the Killing vectors
$$k_1=z_1_{z_1}\overline{z}_1_{\overline{z}_1},k_2=z_2_{z_2}\overline{z}_2_{\overline{z}_2}.$$
(5)
It appears more convenient to introduce polar coordinates, as in , given by
$$z_1=r(\mathrm{cos}\theta /2)e^{i(\psi +\phi )/2},z_2=r(\mathrm{sin}\theta /2)e^{i(\psi \phi )/2}$$
(6)
and the scalar field metric $`ds^2=_{\overline{z}^i}_{z^j}Kd\overline{z}^idz^j`$ becomes
$$ds^2=\frac{dr^2}{(1r^2)^2}+\frac{r^2}{4(1r^2)^2}\left(d\psi +\mathrm{cos}\theta d\phi \right)^2+\frac{r^2}{4(1r^2)}\left(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2\right).$$
(7)
The Kähler form is exact and can be written as
$$K_{uv}dq^udq^v=d\widehat{\omega },\mathrm{where}\widehat{\omega }=\frac{i}{2}\frac{r^2}{1r^2}\tau _3,\tau _3=d\psi +\mathrm{cos}\theta d\phi $$
(8)
with $`q^u=(r,\theta ,\phi ,\psi )`$. Next, we gauge a general linear combination of the two Killing vectors $`k_1,k_2`$ with some constants $`a,b`$
$$k=a(k_1+k_2)/2+b(k_1k_2)/2=4i(a_\psi +b_\phi )$$
(9)
and the scalar derivative becomes $`D_\mu q^u=_\mu q^u+k^uA_\mu `$, where $`A_\mu `$ is the graviphoton (we do not consider vector multiplets). The superpotential is given by the Killing prepotential $`W=P`$ which is defined by the Kähler 2-from
$$K_{vu}k^u=_vP$$
(10)
and becomes
$$W=P=\frac{r^2}{2(1r^2)}\left(a+b\mathrm{cos}\theta \right).$$
(11)
In addition, the Killing spinor equations get corrected by $`W\mathrm{\Gamma }_\mu ϵ`$ for the gravitino variation and $`ik^uϵ`$ for the hyperino variation. There are numerous different conventions yielding different factors in these variations, but they can be fixed by the condition, that the vacuum is given by extrema of $`W`$ and that the difference of the extrema on each side of the wall gives the energy stored by the wall. Moreover, solutions of the BPS equations solve also the equations of motion for our Lagrangian (2); we come back to these equations in the next section.
## 3 BPS domain wall solution
In supergravity one refers to domain walls as kink solution interpolating between different extrema of the potential. As an Ansatz for the metric which is adapted to the RG flow interpreation and preserves 4-d Poincare invariance one can use
$$ds^2=\mu ^2\left(dt^2+d\stackrel{}{x}^2\right)+\frac{d\mu ^2}{\left(\mu \widehat{W}(\mu )\right)^2},$$
(12)
where the fifth coordinate $`\mu `$ will be identified with an energy scale in the dual 4-d field theory. In these coordinates the UV region (= large length scale in supergravity) corresponds to $`\mu \mathrm{}`$, whereas the IR is approached for $`\mu 0`$. For our purpose we are only interested in the dependence of the scalar fields on the fifth coordinate or, in terms of the dual field theory, on the energy scale $`\mu `$. Using this ansatz for the metric the first order equations that solve the equations of motion are
$$W=\pm \widehat{W},\beta ^M\mu \frac{d}{d\mu }\mathrm{\Phi }^M=3g^{MN}_N\mathrm{log}W.$$
(13)
Using the projector $`\mathrm{\Gamma }^5ϵ=\pm ϵ`$ the first equation is equivalent to the gravitino and the second to the hyperino variation. Oviously, supersymmetric extrema of $`V`$ occur at $`_MW=0`$ where also $`\beta ^M=0`$ holds. For $`W|_{W=0}0`$ one has an AdS space with $`W`$ being the cosmological constant while $`W|_{W=0}=0`$ corresponds to a flat space time. From the field theory point of view, $`\beta ^M`$ is a natural candidate for the $`\beta `$-functions of the couplings related to the supergravity scalar fields. For BPS solutions these functions determine the holographic RG flow .
The nature of the fixed point is determined by the derivatives of the $`\beta `$-functions<sup>3</sup><sup>3</sup>3Here we assume that the fixed point is non singular i.e. the metric non-degenerate.
$$_N\beta ^M|_{\beta =0}=3g^{MK}\frac{_N_KW}{W}|_{\beta =0}.$$
(14)
In the case that all scalars are in vector multiplets one finds $`_M\beta ^N|_{\beta =0}=2\delta _M^N`$ and all fixed points are necessarily UV attractive . Consequently, these models cannot describe a smooth RG flow or give a smooth RS scenario, where each side of the wall has to be IR attractive. This situation is very typical for many potentials coming from string or M-theory compactification, which do not allow for “good” domain walls with at least two smoothly connected extrema. The situation in four dimensions is better (for a review see ), but 5-d supergravity is more restrictive. In fact, generically one has no isolated extrema and instead a “run-away” potential giving rise to dilatonic domain walls. This behaviour is caused by the scalar field parameterizing the volume of the internal space. Whenever this scalar is dynamical it runs either to zero or infinite volume. There is no mechanism known to stabilize the volume at a finite value, at least not in a supersymmetric way. This happens also in our case, where the radial part allows only for a single extrema where the superpotential vanishes and the spacetime becomes flat. On the other hand the volume scalar can become non-dynamical. For example, quantum corrections in string or $`M`$-theory can provide a natural lower bound on the volume , which cut-off the radial flow, see . Or, adopting the procedure discussed for 4-d vacua , we can consider an infinite volume limit (non-compact Calabi-Yau), which becomes equivalent to treat the $`r`$-coordinate as non-dynamical (i.e. constant). In doing this we can normalize $`g_{\theta \theta }=1`$ and consider first the model for the superpotential
$$W=2(a+b\mathrm{cos}\theta ).$$
(15)
At the end we will also comment on the solution coming from a running radial coordinate. Coming back to our analysis from before, this superpotential has the properties we are looking for. It has two extrema where $`\mathrm{cos}\theta =\pm 1`$ with
$$_\theta \beta ^\theta =3b\frac{b+a\mathrm{cos}\theta }{(a+b\mathrm{cos}\theta )^2},\mathrm{\Delta }=_\theta \beta ^\theta |_{\beta =0}=\frac{3b}{b\pm a}$$
(16)
where $`\mathrm{\Delta }`$ corresponds to the scaling dimension of the corresponding field theory operator. Therefore, for $`a>b>0`$ it describes a flow from the UV at $`\mathrm{cos}\theta =1`$ towards the IR at $`\mathrm{cos}\theta =1`$. If $`b>a0`$ we have two IR fixed points and similarly for $`0>a>b`$ one finds on each side of the wall an UV fixed point. Finally for $`a=b>0`$ it is a domain wall between an IR point at $`\mathrm{cos}\theta =1`$ and flat spacetime at $`\mathrm{cos}\theta =1`$. Thus, this simple model describes all known types of supersymmetric domain wall solutions. If both sides have the same type of fixed points, $`W`$ has to change its sign and the solution cannot be interpreted as RG-flow, see also . The pole in the $`\beta `$-functions indicates a first order phase transition.
As we mentioned earlier, in order to trust the embedding of our model into $`N=2`$ supergravity, we have to ensure, that the scaling dimensions fit into short representations of superfields of the dual field theory, see e.g. , and there are some interesting cases. If we gauge e.g., only the $`_\phi `$ isometry we obtain $`a=0`$ and $`\mathrm{\Delta }=3`$; if we gauge only $`k_1`$ or $`k_2`$ we have $`a=\pm b`$ with $`\mathrm{\Delta }=\frac{3}{2}`$ or if we turn off the $`_\phi `$ gauging there are no running hyper scalars and we obtain a special case of the model discuss in with $`\mathrm{\Delta }=2`$ for the vector scalars.
Next, let us construct the explicit domain wall solution. The coordinate system used before was adapted for the discussion of the RG flow, but in order to find an explicit solution we write the metric as
$$ds^2=e^{2A(z)}\left(dt^2+d\stackrel{}{x}^2\right)+dz^2.$$
(17)
For these coordinates the BPS equations become
$$_zA=W,_z\theta =3g^{\theta \theta }_\theta W.$$
(18)
and inserting the superpotential (15) and $`g_{\theta \theta }=1`$, we find as solution
$$e^{2A}=e^{4az}(\mathrm{cosh}6bz)^{2/3},\mathrm{cos}\theta =\mathrm{tanh}6bz.$$
(19)
Approaching the two AdS vacua at $`z=\pm \mathrm{}`$, the warp factor becomes $`e^{2A}e^{4(\pm ab)|z|}`$ and as mentioned before if $`a>b>0`$ we have an UV fixed point at $`z=+\mathrm{}`$ and an IR fixed point on the other side. If $`b>a>0`$ there are IR fixed point on both sides, which becomes $`Z_2`$ symmetric if $`a=0`$. In this case, $`e^{2A}`$ is exponentially decreasing on both sides yielding a localization of gravity on the wall and our model describes a Randall-Sundrum scenario . Having this thick wall, one can also consider a thin-wall approximation ($`b\mathrm{}`$), where the scalar becomes constant and the spacetime is everywhere AdS, up to the discontinuity at $`z=0`$<sup>4</sup><sup>4</sup>4Supersymmetric versions of the Randall-Sundrum geometry have also been discussed in refs. .. Another interesting example is $`a=b`$, where we find
$$e^{2A}=(1+e^{12az})^{2/3}$$
(20)
and the domain wall represents a flow from an IR fixed point at $`z=\mathrm{}`$ to flat spacetime at $`z=+\mathrm{}`$.
Finally let us comment on the case of a running radius $`r`$, i.e. we consider the complete superpotential (11). In the $`\theta `$-equation in (18), the radial part drops out and we obtain the same solution $`\mathrm{cos}\theta =\mathrm{tanh}6bz`$ and in addition we have the equation for the radius
$$_zr=3g^{rr}_rW=3r(a+b\mathrm{cos}\theta )=3(ab\mathrm{tanh}6bz)$$
(21)
which is solved by $`r^2=e^{6az}\mathrm{cosh}6bz`$ for $`a>b`$ (keeping in mind, that $`0r<1`$). Therefore, $`z=0`$ corresponds to the singular point $`r=1`$, where the superpotential and the potential $`V`$ have a pole and the warp factor in the metric has a zero: $`e^{2A}z^{1/6}`$ for $`z0`$. Towards larger values of $`z`$ the warp factor increases and for $`z=+\mathrm{}`$ corresponding to $`r=0`$ the potentials vanish and the spacetime becomes flat. Hence, taking the radial part into account, there are neither UV nor IR fixed points related to an AdS spacetime. Notice that using our definition from before, the $`\beta `$-function $`\beta ^r=6(\frac{1}{r}r)`$ vanishes only at the singularity $`r=1`$, which is IR attractive, as the cases discussed in .
## 4 Some comments about the M-theory embedding
Perhaps the easiest way to understand the solution comes from $`M`$-theory, where the four scalars of the universal hypermultiplet can be understood as follows. One scalar parameterizes the volume of the internal space, one axion comes from the dualization of the 5-d 4-form field strength, and two scalars are related to membranes wrapping $`\mathrm{\Omega }_{0,3}`$ or $`\mathrm{\Omega }_{3,0}`$ cycle. In our parameterization we gauged the Killing vector $`ka_\psi +b_\phi `$. The case $`b=0`$ reproduces the results derived in and therefore corresponds to turning on 4-form fluxes in the $`M`$-theory compactification, which is equivalent to a non-trivial M5-brane background. On the other hand gauging the $`_\phi `$ isometry gives a mass term for the $`\theta `$ axion. In fact for this deformation (setting $`a=0`$) the supergravity potential can be written as (normalizing $`g_{\theta \theta }=3/4`$ and rescaling $`\theta \frac{2\theta }{\sqrt{3}}`$): $`V=6b\mathrm{cos}\frac{4\theta }{\sqrt{3}}`$ and we obtain a sine-Gordon model coupled to gravity. This form suggests, that it can be understood as coming from a non-trivial instanton background. M-theory compactified to 5 dimension may give a topological term $`d^5x\theta 𝑑G`$. If there are no sources present, the Bianchi identity implies $`dG=0`$ and such a term vanishes. If on the other hand magnetic sources for $`G`$ exist in 5 dimensions, then $`dG`$ is non-vanishing, and this term is the 5-d analogue of the familiar universal axionic coupling $`d^4x\psi tr(F\stackrel{~}{F})`$ in 4 dimensions. An obvious candidate for magnetic sources of $`G`$ are M5-branes. However in order to generate a cosine potential for $`\theta `$ such sources must be pointlike in 5 dimensions (like the instanton density $`tr(FF)`$ in 4 dimensions ) and therefore correspond to Euclidean M5-branes wrapping the whole CY manifold, in contrast to the 5-branes associated with the gauging of the $`_\psi `$ isometry, which wrap a holomorphic 2-cycle and become 3-branes upon compactification. As in 4 dimensions, see , we expect that the sum over an M5-brane instanton gas will reproduce the cosine potential. Clearly, this interpretation deserves a more detailed investigation.
Acknowledgements
The work is supported by a Heisenberg grant of the DFG. I am grateful to S. Ferrara, S. Gukov, C. Pope and A. Zaffaroni for discussions and comments. I would also like the thank C. Herrmann, J. Louis and S. Thomas for collaboration on part of this project.
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# The Long-Time Dynamics of Dirac Particles in the Kerr-Newman Black Hole Geometry
## 1 Introduction
It has recently been shown that the Dirac equation does not admit normalizable, time-periodic solutions in the non-extreme Kerr-Newman axisymmetric black hole geometry . This was interpreted as an indication that a Dirac particle either falls into the black hole or escapes to infinity, but that it cannot stay in a bounded region outside the event horizon. In this paper we make precise this interpretation in the general time-dependent setting. We thus consider the Cauchy problem for the Dirac equation with smooth initial data on the hypersurface $`t=0`$, compactly supported outside the event horizon. We study the probability for the Dirac particle to be inside a given annulus located outside the event horizon, and we prove that this probability tends to zero as $`t`$ goes to infinity. Hence, in contrast to the situation for a bounded orbit of a classical point particle, there exists no compact region outside the event horizon in which the quantum mechanical Dirac particle will remain for all time. In more precise form, our result is stated as follows. Recall that in Boyer-Lindquist coordinates $`(t,r,\vartheta ,\phi )`$ with $`r>0`$, $`0\vartheta \pi `$, $`0\phi <2\pi `$, the Kerr-Newman metric takes the form
$$\begin{array}{ccc}ds^2\hfill & =& g_{jk}dx^jx^k\hfill \\ & =& \frac{\mathrm{\Delta }}{U}(dta\mathrm{sin}^2\vartheta d\phi )^2U\left(\frac{dr^2}{\mathrm{\Delta }}+d\vartheta ^2\right)\hfill \\ & & \frac{\mathrm{sin}^2\vartheta }{U}(adt(r^2+a^2)d\phi )^2\hfill \end{array}$$
(1.1)
with
$$U(r,\vartheta )=r^2+a^2\mathrm{cos}^2\vartheta ,\mathrm{\Delta }(r)=r^22Mr+a^2+Q^2,$$
and the electromagnetic potential is given by
$$A_jdx^j=\frac{Qr}{U}(dta\mathrm{sin}^2\vartheta d\phi ),$$
where $`M`$, $`aM`$ and $`Q`$ denote the mass, the angular momentum and the charge of the black hole, respectively. We shall here restrict attention to the non-extreme case $`M^2>a^2+Q^2`$. Then the function $`\mathrm{\Delta }`$ has two distinct zeros,
$$r_0=M\sqrt{M^2a^2Q^2}\text{and}r_1=M+\sqrt{M^2a^2Q^2},$$
corresponding to the Cauchy and the event horizon, respectively. We will here consider only the region $`r>r_1`$ outside of the event horizon, and thus $`\mathrm{\Delta }>0`$.
###### Theorem 1.1.
Consider the Cauchy problem for the Dirac equation in the non-extreme Kerr-Newman black hole geometry outside the event horizon
$$(i\gamma ^jD_jm)\mathrm{\Psi }(t,x)=\mathrm{\hspace{0.33em}0},\mathrm{\Psi }(0,x)=\mathrm{\Psi }_0(x),$$
(1.2)
where the initial data $`\mathrm{\Psi }_0`$ is in $`L^2((r_1,\mathrm{})\times S^2,d\mu )^4`$, where $`d\mu `$ is the induced invariant measure on the hypersurface $`t=\text{const}`$. Then for any $`\delta >0`$ and $`R>r_1+\delta `$, the probability for the Dirac particle to be inside the annulus $`K_{\delta ,R}=\{r_1+\delta rR\}`$ tends to zero as $`t\mathrm{}`$, i.e.
$$\underset{t\mathrm{}}{lim}_{K_{\delta ,R}}(\overline{\mathrm{\Psi }}\gamma ^j\mathrm{\Psi })(t,x)\nu _j𝑑\mu =\mathrm{\hspace{0.33em}0},$$
(1.3)
where $`\nu `$ denotes the future directed normal.
The decay of probabilities in compact sets (1.3) can be stated equivalently that the Dirac wave function decays in $`L_{\text{loc}}^2`$ outside and away from the event horizon. Since the Dirac equation is linear and stationary, for smooth initial data we obtain immediately that also the time-derivatives $`_t^n\mathrm{\Psi }`$ decay in $`L_{\text{loc}}^2`$, and standard Sobolev methods yield that $`\mathrm{\Psi }`$ decays even in $`L_{\text{loc}}^{\mathrm{}}`$. We point out that the initial data is merely bounded (but not necessarily small), near the event horizon. Our assumptions include the case when the initial data is smooth and bounded in the maximal Kruskal extension up to the bifurcation $`2`$-sphere (as is considered in for the wave function in the Schwarzschild geometry). We note that the axisymmetric character of the background geometry makes the analysis significantly more delicate than in the spherically symmetric case, mainly because for $`a0`$ both the radial and angular ODEs depend on the energy, and thus for the study of the dynamics we must consider the system of these coupled equations.
The proof is organized as follows. We first bring the Dirac equation into the Hamiltonian form $`i_t\mathrm{\Psi }=H\mathrm{\Psi }`$ with a self-adjoint operator $`H`$. Our main technical tool is an integral representation for the Dirac propagator $`\mathrm{exp}(itH)`$ acting on wave functions with compact support outside the event horizon. This integral representation is stated in Theorem 3.6. To derive it, we first consider the Dirac equation in an annulus outside the event horizon with suitable Dirichlet-type boundary conditions. The Hamiltonian corresponding to this modified problem has a purely discrete spectrum, and thus the propagator can be decomposed into discrete eigenstates. We then take the infinite-volume limit by letting the inner boundary of the annulus tend to the event horizon and the outer boundary to infinity in a suitable way. Our construction is based on Chandrasekhar’s separation of variables for the Dirac equation in the Kerr-Newman metric together with estimates for the asymptotic behavior of the amplitudes and phases of the separated radial eigenfunctions (Lemmas 3.1 and 3.5), and for the spectral gaps (Lemma 3.3). The usefulness of our integral representation for the propagator is that it explicitly gives the continuous spectral measure of $`H`$ in terms of the solutions of the radial and angular ODEs arising in Chandrasekhar’s separation of variables. For initial data which is compactly supported outside the event horizon, the decay of the probabilities (1.3) then follows by standard results of Fourier analysis. The generalization to initial data in $`L^2`$ and $`L_{\text{loc}}^{\mathrm{}}`$ near the event horizon is done by approximation in our Hilbert space framework.
## 2 Separation of Variables, Hamiltonian Formulation
The Dirac equation in the Kerr-Newman geometry can be completely separated into ODEs by Chandrasekhar’s method . We here outline the separation, see for details. After the regular and time-independent transformation
$$\mathrm{\Psi }\widehat{\mathrm{\Psi }}=S\mathrm{\Psi }$$
(2.1)
with
$`S`$ $`=`$ $`\mathrm{\Delta }^{\frac{1}{4}}\text{diag}((ria\mathrm{cos}\vartheta )^{\frac{1}{2}},(ria\mathrm{cos}\vartheta )^{\frac{1}{2}},`$
$`(r+ia\mathrm{cos}\vartheta )^{\frac{1}{2}},(r+ia\mathrm{cos}\vartheta )^{\frac{1}{2}}),`$
the Dirac equation can be written as
$$(+𝒜)\widehat{\mathrm{\Psi }}=\mathrm{\hspace{0.33em}0}$$
(2.2)
with
$``$ $`=`$ $`\left(\begin{array}{cccc}imr& 0& \sqrt{\mathrm{\Delta }}𝒟_+& 0\\ 0& imr& 0& \sqrt{\mathrm{\Delta }}𝒟_{}\\ \sqrt{\mathrm{\Delta }}𝒟_{}& 0& imr& 0\\ 0& \sqrt{\mathrm{\Delta }}𝒟_+& 0& imr\end{array}\right)`$
$`𝒜`$ $`=`$ $`\left(\begin{array}{cccc}am\mathrm{cos}\vartheta & 0& 0& _+\\ 0& am\mathrm{cos}\vartheta & _{}& 0\\ 0& _+& am\mathrm{cos}\vartheta & 0\\ _{}& 0& 0& am\mathrm{cos}\vartheta \end{array}\right)`$
and the differential operators
$`𝒟_\pm `$ $`=`$ $`{\displaystyle \frac{}{r}}{\displaystyle \frac{1}{\mathrm{\Delta }}}\left[(r^2+a^2){\displaystyle \frac{}{t}}+a{\displaystyle \frac{}{\phi }}ieQr\right]`$
$`_\pm `$ $`=`$ $`{\displaystyle \frac{}{\vartheta }}+{\displaystyle \frac{\mathrm{cot}\vartheta }{2}}i\left[a\mathrm{sin}\vartheta {\displaystyle \frac{}{t}}+{\displaystyle \frac{1}{\mathrm{sin}\vartheta }}{\displaystyle \frac{}{\phi }}\right].`$
Employing the ansatz
$$\widehat{\mathrm{\Psi }}(t,r,\vartheta ,\phi )=e^{i\omega t}e^{i(k+\frac{1}{2})\phi }\left(\begin{array}{c}X_{}(r)Y_{}(\vartheta )\\ X_+(r)Y_+(\vartheta )\\ X_+(r)Y_{}(\vartheta )\\ X_{}(r)Y_+(\vartheta )\end{array}\right),k\text{Z},$$
(2.5)
we obtain the eigenvalue problems,
$$\widehat{\mathrm{\Psi }}=\lambda \widehat{\mathrm{\Psi }},𝒜\widehat{\mathrm{\Psi }}=\lambda \widehat{\mathrm{\Psi }},$$
(2.6)
under which the Dirac equation (2.2) decouples into the system of ODEs
$`\left(\begin{array}{cc}\sqrt{\mathrm{\Delta }}𝒟_+& imr\lambda \\ imr\lambda & \sqrt{\mathrm{\Delta }}𝒟_{}\end{array}\right)\left(\begin{array}{c}X_+\\ X_{}\end{array}\right)`$ $`=`$ $`0`$ (2.11)
$`\left(\begin{array}{cc}_+& am\mathrm{cos}\vartheta +\lambda \\ am\mathrm{cos}\vartheta +\lambda & _{}\end{array}\right)\left(\begin{array}{c}Y_+\\ Y_{}\end{array}\right)`$ $`=`$ $`0,`$ (2.16)
where $`𝒟_\pm `$ and $`_\pm `$ are the radial and angular operators
$`𝒟_\pm `$ $`=`$ $`{\displaystyle \frac{}{r}}\pm {\displaystyle \frac{i}{\mathrm{\Delta }}}\left[\omega (r^2+a^2)+\left(k+{\displaystyle \frac{1}{2}}\right)a+eQr\right]`$ (2.17)
$`_\pm `$ $`=`$ $`{\displaystyle \frac{}{\vartheta }}+{\displaystyle \frac{\mathrm{cot}\vartheta }{2}}\left[a\omega \mathrm{sin}\vartheta +{\displaystyle \frac{k+\frac{1}{2}}{\mathrm{sin}\vartheta }}\right].`$ (2.18)
We will in what follows also use the vector notation $`X=(X_+,X_{})`$, $`Y=(Y_{},Y_+)`$ and for clarity sometimes add indices for the parameters involved, e.g. $`X^{k\omega \lambda }X`$. We point out that (2.5) is an eigenfunction of the angular operator $`i_\phi `$ with eigenvalue $`k+\frac{1}{2}`$. The reason why we need to consider half odd integer eigenvalues is that the transformation from the usual single-valued wave function in space-time, to the wave function $`\widehat{\mathrm{\Psi }}`$ in (2.2) involves a sign flip at $`\phi =0`$ (see \[1, Section 2.1\]).
In this paper, we want to study time-dependent solutions of the Dirac equation. In the separation ansatz (2.5), the dynamics of the solution is encoded through the $`\omega `$-dependence in the ODEs (2.11) and (2.16). Unfortunately, both the radial and angular operators (2.17) and (2.18) depend on $`\omega `$, making the situation rather complicated. Therefore it is useful to bring the Dirac equation (2.2) into Hamiltonian form, in a way which is compatible with the separation of variables. We first bring the time derivative in (2.2) to one side of the equation and obtain
$$\left(\frac{r^2+a^2}{\sqrt{\mathrm{\Delta }}}B+a\mathrm{sin}\vartheta C\right)i\frac{}{t}\widehat{\mathrm{\Psi }}=(^3+𝒜^3)\widehat{\mathrm{\Psi }}$$
(2.19)
with
$$B=\left(\begin{array}{cccc}0& 0& i& 0\\ 0& 0& 0& i\\ i& 0& 0& 0\\ 0& i& 0& 0\end{array}\right),C=\left(\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& 1& 0\\ 0& 1& 0& 0\\ 1& 0& 0& 0\end{array}\right),$$
where the operators $`^3`$ and $`𝒜^3`$ are obtained from $``$ and $`𝒜`$ by setting the time derivatives to zero. The matrices $`B`$ and $`C`$ satisfy the relations $`B^2=\text{1 1}=C^2`$ and $`BC=CB`$. Thus the linear combination of these matrices which appears in (2.19) can be inverted with the formula $`(\alpha B+\beta C)^1=(\alpha ^2\beta ^2)^1(\alpha B\beta C)`$ (and $`\alpha ,\beta \text{I R}`$). Furthermore, we introduce a new radial variable $`u(\mathrm{},\mathrm{})`$ by
$$\frac{du}{dr}=\frac{r^2+a^2}{\mathrm{\Delta }}.$$
(2.20)
Omitting for simplicity the hat of the wave function, the Dirac equation (2.19) becomes
$$i\frac{}{t}\mathrm{\Psi }=H\mathrm{\Psi }$$
with the Hamiltonian
$`H`$ $`=`$ $`\left({\displaystyle \frac{(r^2+a^2)^2}{\mathrm{\Delta }}}a^2\mathrm{sin}^2\vartheta \right)^1\left({\displaystyle \frac{r^2+a^2}{\sqrt{\mathrm{\Delta }}}}Ba\mathrm{sin}\vartheta C\right)(^3+𝒜^3)`$ (2.21)
$`=`$ $`\left[\left(1{\displaystyle \frac{a^2\mathrm{\Delta }\mathrm{sin}^2\vartheta }{(r^2+a^2)^2}}\right)^1\left(\text{1 1}{\displaystyle \frac{a\sqrt{\mathrm{\Delta }}\mathrm{sin}\vartheta }{r^2+a^2}}BC\right)\right](\widehat{}+\widehat{𝒜}),`$
where $`r`$ is an implicit function of $`u`$, and
$`\widehat{}`$ $`=`$ $`{\displaystyle \frac{mr\sqrt{\mathrm{\Delta }}}{r^2+a^2}}\left(\begin{array}{cccc}0& 0& 1& 0\\ 0& 0& 0& 1\\ 1& 0& 0& 0\\ 0& 1& 0& 0\end{array}\right)+\left(\begin{array}{cccc}_{}& 0& 0& 0\\ 0& _+& 0& 0\\ 0& 0& _+& 0\\ 0& 0& 0& _{}\end{array}\right)`$ (2.30)
$`\widehat{𝒜}`$ $`=`$ $`{\displaystyle \frac{am\mathrm{cos}\vartheta \sqrt{\mathrm{\Delta }}}{r^2+a^2}}\left(\begin{array}{cccc}0& 0& i& 0\\ 0& 0& 0& i\\ i& 0& 0& 0\\ 0& i& 0& 0\end{array}\right)`$ (2.40)
$`+\left(\begin{array}{cccc}0& _+& 0& 0\\ _{}& 0& 0& 0\\ 0& 0& 0& _+\\ 0& 0& _{}& 0\end{array}\right)`$
with
$`_\pm `$ $`=`$ $`i{\displaystyle \frac{}{u}}\left({\displaystyle \frac{ia}{r^2+a^2}}{\displaystyle \frac{}{\phi }}+{\displaystyle \frac{eQr}{r^2+a^2}}\right)`$
$`_\pm `$ $`=`$ $`{\displaystyle \frac{\sqrt{\mathrm{\Delta }}}{r^2+a^2}}\left(i{\displaystyle \frac{}{\vartheta }}+i{\displaystyle \frac{\mathrm{cot}\vartheta }{2}}\pm {\displaystyle \frac{1}{\mathrm{sin}\vartheta }}{\displaystyle \frac{}{\phi }}\right).`$
The Hamiltonian (2.21) is an operator acting on the wave functions on the hypersurfaces $`t=\text{const}`$. The simplest scalar product on such a hypersurface is
$$(\mathrm{\Psi }|\mathrm{\Phi })=_{\mathrm{}}^{\mathrm{}}𝑑u_1^1d\mathrm{cos}\vartheta _0^{2\pi }𝑑\phi \overline{\mathrm{\Psi }}(t,u,\vartheta ,\phi )\mathrm{\Phi }(t,u,\vartheta ,\phi ),$$
(2.41)
where “$`\overline{\mathrm{\Psi }}`$” denotes the complex conjugated, transposed spinor. In the spherically symmetric case $`a=0`$, the Hamiltonian (2.21) is Hermitian (i.e. formally self-adjoint) with respect to this scalar product. However for $`a0`$, $`H`$ is not Hermitian. In order to get around this problem, we introduce a different scalar product as follows. Notice that the operators $`\widehat{}`$ and $`\widehat{𝒜}`$, (2.30),(2.40), are both Hermitian with respect to (2.41). The reason why the Hamiltonian (2.21) is not Hermitian is that, when the taking the adjoint of $`H`$ using integration by parts, one gets $`r`$\- and $`\vartheta `$-derivatives of the square bracket in (2.21). But we can compensate this square bracket by inserting its inverse into the scalar product. Thus we introduce on the four-component spinors the inner product
$$<\text{ }\mathrm{\Psi }|\mathrm{\Phi }>_{(t,u,\vartheta ,\phi )}=\overline{\mathrm{\Psi }}(t,u,\vartheta ,\phi )\left(\text{1 1}+\frac{a\sqrt{\mathrm{\Delta }}\mathrm{sin}\vartheta }{r^2+a^2}BC\right)\mathrm{\Phi }(t,u,\vartheta ,\phi )$$
(2.42)
and define the scalar product $`<\text{ }.|.>`$ by
$$<\text{ }\mathrm{\Psi }|\mathrm{\Phi }>=_{\mathrm{}}^{\mathrm{}}𝑑u_1^1d\mathrm{cos}\vartheta _0^{2\pi }𝑑\phi <\text{ }\mathrm{\Psi }|\mathrm{\Phi }>_{(t,u,\vartheta ,\phi )}.$$
(2.43)
Then the Hamiltonian $`H`$ is Hermitian with respect to (2.43). Let us verify that (2.43) is positive. In the region outside the event horizon under consideration, $`r>r_1>M`$. Also, since we are in the non-extreme case, $`M>Q,a`$, and as a consequence, $`\mathrm{\Delta }<r^2`$. We conclude that
$$\left|\frac{a\sqrt{\mathrm{\Delta }}\mathrm{sin}\vartheta }{r^2+a^2}\right|\frac{a\sqrt{\mathrm{\Delta }}}{r^2+a^2}<\frac{a}{r}\frac{\sqrt{\mathrm{\Delta }}}{r}<\mathrm{\hspace{0.33em}1}.$$
Combining this inequality with the fact that the matrix $`BC`$ has eigenvalues $`\pm 1`$, we obtain that the bracket in (2.42) is indeed a positive matrix.
We denote the Hilbert space of wave functions with scalar product (2.43) by $``$. Then the operator $`H`$, (2.21), is essentially self-adjoint on $``$ with domain of definition
$$D(H)=C_0^{\mathrm{}}(\text{I R}\times S^2)^4.$$
In Section 3, we shall consider the Dirac operator also with certain Dirichlet-type boundary conditions, which we now introduce. First for given $`u_2\text{I R}`$, we restrict $`u`$ to the half line $`u(\mathrm{},u_2]`$ and impose the boundary conditions
$$\mathrm{\Psi }_1(u_2,\vartheta ,\phi )=\mathrm{\Psi }_3(u_2,\vartheta ,\phi )\text{and}\mathrm{\Psi }_2(u_2,\vartheta ,\phi )=\mathrm{\Psi }_4(u_2,\vartheta ,\phi ).$$
(2.44)
Let $`_{u_2}`$ be the Hilbert space of square integrable wave functions $`\mathrm{\Psi }(u,\vartheta ,\phi )`$, $`uu_2`$ with the scalar product
$$<\text{ }\mathrm{\Psi }|\mathrm{\Phi }>_{u_2}:=_{\mathrm{}}^{u_2}𝑑u_1^1d\mathrm{cos}\vartheta _0^{2\pi }𝑑\phi <\text{ }\mathrm{\Psi }|\mathrm{\Phi }>_{(t,u,\vartheta ,\phi )}.$$
(2.45)
Then the Hamiltonian (2.21) on $`_{u_2}`$ with boundary conditions (2.44), which we denote for clarity by $`H_{u_2}`$, is Hermitian (the main point here is that the boundary values at $`u=u_2`$ vanish when the adjoint of $`H_{u_2}`$ is calculated using integration by parts). The operator $`H_{u_2}`$ is essentially self-adjoint on $`_{u_2}`$ with domain of definition
$$D(H_{u_2})=\left\{\mathrm{\Psi }C_0^{\mathrm{}}((\mathrm{},u_2]\times S^2)^4\text{ and (}\text{2.44}\text{) is satisfied}\right\}.$$
Similarly, for $`u_1,u_2\text{I R}`$, $`u_1<u_2`$, we restrict $`u`$ to the closed interval $`u[u_1,u_2]`$ with boundary conditions
$$\begin{array}{c}\mathrm{\Psi }_1(u_1)=\mathrm{\Psi }_3(u_1),\mathrm{\Psi }_2(u_1)=\mathrm{\Psi }_4(u_1)\hfill \\ \mathrm{\Psi }_1(u_2)=\mathrm{\Psi }_3(u_2),\mathrm{\Psi }_2(u_2)=\mathrm{\Psi }_4(u_2).\hfill \end{array}$$
(2.46)
We denote the Hilbert space of square integrable wave functions $`\mathrm{\Psi }(u,\vartheta ,\phi )`$, $`u_1uu_2`$, with the scalar product
$$<\text{ }\mathrm{\Psi },\mathrm{\Phi }>_{u_1,u_2}:=_{u_1}^{u_2}du_1^1d\mathrm{cos}\vartheta _0^{2\pi }d\phi <\text{ }\mathrm{\Psi }|\mathrm{\Phi }>_{(t,u,\vartheta ,\phi )}$$
(2.47)
by $`_{u_1,u_2}`$, and the Hamiltonian (2.21) on $`_{u_1,u_2}`$ by $`H_{u_1,u_2}`$. It is essentially self-adjoint with
$$D(H_{u_1,u_2})=\left\{\mathrm{\Psi }C_0^{\mathrm{}}([u_1,u_2]\times S^2)^4\text{ and (}\text{2.46}\text{) is satisfied}\right\}.$$
Our above Hamiltonian formulation of the Dirac equation is well-suited to Chandrasekhar’s separation of variables. Namely, the boundary conditions (2.44) reduce to simple boundary conditions for the radial functions,
$$X_+(u_2)=X_{}(u_2),$$
(2.48)
whereas (2.46) amounts to
$$X_+(u_1)=X_{}(u_1)\text{and}X_+(u_2)=X_{}(u_2).$$
(2.49)
The scalar product (2.41) splits into the product of a radial and an angular part, namely
$$(\widehat{\mathrm{\Psi }}^{k\omega \lambda }|\widehat{\mathrm{\Psi }}^{k^{}\omega ^{}\lambda ^{}})=(X^{k\omega \lambda }|X^{k^{}\omega ^{}\lambda ^{}})(Y^{k\omega \lambda }|Y^{k^{}\omega ^{}\lambda ^{}})$$
with
$`(X^{k\omega \lambda }|X^{k^{}\omega ^{}\lambda ^{}})`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\overline{X^{k\omega \lambda }}(u)X^{k^{}\omega ^{}\lambda ^{}}(u)𝑑u`$
$`(Y^{k\omega \lambda }|Y^{k^{}\omega ^{}\lambda ^{}})`$ $`=`$ $`2\pi \delta ^{kk^{}}{\displaystyle _1^1}\overline{Y^{k\omega \lambda }}(\vartheta )Y^{k^{}\omega ^{}\lambda ^{}}(\vartheta )d\mathrm{cos}\vartheta .`$
The scalar product (2.43), however, does not split into a product, more precisely
$`<\text{ }\mathrm{\Psi }|\mathrm{\Phi }>=(X^{k\omega \lambda }|X^{k^{}\omega ^{}\lambda ^{}})(Y^{k\omega \lambda }|Y^{k^{}\omega ^{}\lambda ^{}})`$ (2.50)
$`+a(X^{k\omega \lambda }|{\displaystyle \frac{\sqrt{\mathrm{\Delta }}}{r^2+a^2}}\sigma ^2|X^{k^{}\omega ^{}\lambda ^{}})(Y^{k\omega \lambda }|\mathrm{sin}\vartheta \sigma ^1|Y^{k^{}\omega ^{}\lambda ^{}}),`$
where $`\sigma ^i`$ are the Pauli matrices. This mixing of the radial and angular parts in the scalar product can be understood from the fact that the Kerr-Newman solution is only axisymmetric.
## 3 An Integral Representation for the Propagator
The propagator $`\mathrm{exp}(itH)`$ has the spectral decomposition
$$e^{itH}=_{\mathrm{}}^{\mathrm{}}e^{i\omega t}𝑑E_\omega ,$$
(3.1)
where $`dE_\omega `$ is the spectral measure of $`H`$. In this section, we shall bring this formula into a more explicit form. This will be done by expressing the spectral measure in terms of solutions of the radial and angular ODEs of the previous section. Since the spectrum of $`H`$ is continuous, it is not obvious how to relate the spectral measure to the solutions of our ODEs. To bypass this problem, we begin with the spectral decomposition of the operator $`H_{u_1,u_2}`$ (which has a purely discrete spectrum), and then deduce the desired integral representation for $`\mathrm{exp}(itH)`$ by taking suitable limits $`u_1\mathrm{}`$ and $`u_2\mathrm{}`$.
As an elliptic operator on a bounded domain, the Hamiltonian $`H_{u_1,u_2}`$ has a purely discrete spectrum with finite-dimensional eigenspaces (see ). In view of our separation of variables, the most convenient eigenvector basis is the following. First we can choose the basis vectors as eigenvectors of the operator $`i_\phi `$ with eigenvalue $`k+\frac{1}{2}`$, $`k\text{Z}`$. We denote this eigenspace of $`i_\phi `$ by $`_{u_1,u_2}^k`$, and the restriction of $`H_{u_1,u_2}`$ to $`_{u_1,u_2}^k`$ by $`H_{u_1,u_2}^k`$. Furthermore, the basis vectors can be chosen as eigenvectors of the angular operator $`𝒜`$. As is shown in the Appendix, the spectrum of $`𝒜`$ on $`_{u_1,u_2}^k`$ is discrete, non-degenerate, and depends smoothly on $`\omega `$. Thus the eigenvalues of $`𝒜`$ can be written as $`\lambda _n(\omega )`$, $`n\text{Z}`$, with $`\lambda _n<\lambda _{n+1}`$ and $`\lambda _n(.)C^{\mathrm{}}(\text{I R})`$. For any given $`k\text{Z}`$, $`\omega \sigma (H_{u_1,u_2}^k)`$, and $`n\text{Z}`$, the radial ODE (2.11) has at most one solution satisfying the boundary conditions (2.46). Hence we have for any $`k`$, $`\omega `$, and $`n`$ at most one eigenstate of $`H_{u_1,u_2}`$, which we denote by $`\mathrm{\Psi }_{u_1,u_2}^{k\omega n}`$. The set of $`n`$ for which such an eigenvector exists is denoted by $`N(k,\omega )`$. Thus our eigenvector basis is
$$(\mathrm{\Psi }_{u_1,u_2}^{k\omega n})_{k\text{Z},\omega \sigma (H_{u_1,u_2}^k),nN(k,\omega )}.$$
(3.2)
We normalize these eigenfunctions with respect to the scalar product (2.41); more precisely, we normalize both the radial and angular parts according to
$$(X_{u_1,u_2}^{k\omega n}|X_{u_1,u_2}^{k\omega n})=\mathrm{\hspace{0.33em}1},(Y^{k\omega n}|Y^{k\omega n})=\mathrm{\hspace{0.33em}1}$$
(3.3)
with $`X`$ and $`Y`$ as in (2.5). Since the angular operator $`𝒜`$ is self-adjoint with respect to the scalar product $`(.|.)`$, its eigenvectors are orthogonal, and thus the eigenfunctions for fixed $`k`$ and $`\omega `$ are even orthonormal,
$$(\mathrm{\Psi }_{u_1,u_2}^{k\omega n}|\mathrm{\Psi }_{u_1,u_2}^{k\omega n^{}})=\delta ^{nn^{}},n,n^{}N(k,\omega ).$$
(3.4)
We mention for clarity that for different values of $`\omega `$, the eigenfunctions are in general not orthogonal with respect to $`(.|.)`$, but since $`H_{u_1,u_2}`$ is self-adjoint with respect to $`<\text{ }.|.>`$, its eigenspaces are orthogonal with respect to the latter scalar product, and thus
$$<\text{ }\mathrm{\Psi }_{u_1,u_2}^{k\omega n}|\mathrm{\Psi }_{u_1,u_2}^{k^{}\omega ^{}n^{}}>=\mathrm{\hspace{0.33em}0}\text{for }\omega \omega ^{}.$$
These subtle differences between the two scalar products clearly become irrelevant in the spherically symmetric case $`a=0`$.
In the basis (3.2), the spectral decomposition (3.1) for $`H_{u_1,u_2}`$ can be written as
$`e^{itH_{u_1,u_2}}\mathrm{\Psi }`$ (3.5)
$`=`$ $`{\displaystyle \underset{k\text{Z}}{}}{\displaystyle \underset{\omega \sigma (H_{u_1,u_2}^k)}{}}e^{i\omega t}\left({\displaystyle \underset{n,n^{}N(k,\omega )}{}}c_{nn^{}}\mathrm{\Psi }_{u_1,u_2}^{k\omega n}<\text{ }\mathrm{\Psi }_{u_1,u_2}^{k\omega n^{}}|\mathrm{\Psi }>\right).`$
Here the coefficients $`c_{nn^{}}`$ must be chosen such that the bracket in (3.5) is the projection of $`\mathrm{\Psi }`$ onto the eigenspace of $`H_{u_1,u_2}^k`$ corresponding to the eigenvalue $`\omega `$; more precisely,
$$c_{nn^{}}=(A^1)_{nn^{}}\text{with}A_{nn^{}}=<\text{ }\mathrm{\Psi }_{u_1,u_2}^{k\omega n}|\mathrm{\Psi }_{u_1,u_2}^{k\omega n^{}}>.$$
(3.6)
Notice that the first two sums in (3.5) give a decomposition of $`\mathrm{\Psi }`$ into the orthogonal eigenstates of the operators $`i_\phi `$ and $`H`$, respectively, and thus converge in norm in $`_{u_1,u_2}`$. The bracket in (3.5) is the basis representation of the projector on the respective eigenspace.
Our first goal is to take the limit $`u_1\mathrm{}`$ in (3.5). We expect that in this infinite volume limit, the “energy gaps” $`\mathrm{\Delta }\omega _{kn}`$ between neighboring eigenvalues, defined by
$`\mathrm{\Delta }\omega _{kn}`$ $`=`$ $`\mathrm{min}\{\stackrel{~}{\omega }_{kn}\omega _{kn}|\stackrel{~}{\omega }_{kn}>\omega _{kn}\}\text{with}`$
$`\omega _{kn},\stackrel{~}{\omega }_{kn}\sigma (H_{u_1,u_2}^k)\text{ and }N(k,\omega _{kn}),N(k,\stackrel{~}{\omega }_{kn})0,`$
should tend to zero. The basic idea is to rewrite the sum over the spectrum in (3.5) as Riemann sums which converge to integrals as $`u_1\mathrm{}`$, yielding a formula for the propagator of the Hamiltonian $`H_{u_2}`$. For making this idea mathematically precise, it is essential to get good estimates for $`\mathrm{\Delta }\omega _{kn}`$ and to relate the eigenvectors $`\mathrm{\Psi }_{u_1,u_2}^{k\omega n}`$ in (3.5) to solutions
$$\mathrm{\Psi }_{u_2}^{k\omega n}(u)\text{with }k\text{Z},\omega \text{I R},n\text{Z},u(\mathrm{},u_2]$$
of the Dirac equation with boundary conditions (2.44). We note that the convergence of the series in $`n`$ is not a real issue. Indeed, using an $`L^2`$ approximation argument, we will show in the proof of Theorem 3.6 that we may restrict attention to a finite number of angular momentum modes.
We denote the radial and angular functions corresponding to $`\mathrm{\Psi }_{u_2}^{k\omega n}`$ by $`X_{u_2}^{k\omega n}`$ and $`Y^{k\omega n}`$, respectively. In the variable $`u`$, (2.20), the radial equation (2.11) becomes
$`\left[{\displaystyle \frac{d}{du}}+i\mathrm{\Omega }(u)\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\right]X={\displaystyle \frac{\sqrt{\mathrm{\Delta }}}{r^2+a^2}}\left(\begin{array}{cc}0& imr\lambda \\ imr\lambda & 0\end{array}\right)X`$ (3.11)
with
$$\mathrm{\Omega }(u)=\omega +\frac{(k+\frac{1}{2})a+eQr}{r^2+a^2},$$
and where for ease in notation the indices of $`X`$ were omitted. The next lemma describes the asymptotic behavior of $`X(u)`$ as $`u\mathrm{}`$.
###### Lemma 3.1.
Every nontrivial solution $`X`$ of (3.11) with boundary conditions (2.48) is asymptotically as $`u\mathrm{}`$ of the form
$$X(u)=\left(\begin{array}{c}e^{i\mathrm{\Omega }_0u}f_0^+\\ e^{i\mathrm{\Omega }_0u}f_0^{}\end{array}\right)+R_0(u)$$
(3.12)
with
$`f_0`$ $``$ $`0`$ (3.13)
$`\mathrm{\Omega }_0`$ $`=`$ $`\omega +{\displaystyle \frac{(k+\frac{1}{2})a+eQr_1}{r_1^2+a^2}}`$ (3.14)
$`|R_0|`$ $``$ $`ce^{du}`$ (3.15)
and suitable constants $`c,d>0`$, which can be chosen locally uniformly in $`\omega `$.
Proof. Substituting into (3.11) the ansatz
$$X(u)=\left(\begin{array}{c}e^{i\mathrm{\Omega }_0u}f^+(u)\\ e^{i\mathrm{\Omega }_0u}f^{}(u)\end{array}\right),$$
(3.16)
we obtain for $`f`$ the equation
$`{\displaystyle \frac{d}{du}}f`$ $`=`$ $`[i(\mathrm{\Omega }_0\mathrm{\Omega }(u))\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$ (3.22)
$`+{\displaystyle \frac{\sqrt{\mathrm{\Delta }}}{r^2+a^2}}\left(\begin{array}{cc}0& e^{2i\mathrm{\Omega }u}(imr\lambda )\\ e^{2i\mathrm{\Omega }u}(imr\lambda )& 0\end{array}\right)]f.`$
The square bracket vanishes on the event horizon $`r=r_1`$. In the variable $`u`$, this leads to exponential decay for $`u\mathrm{}`$, in the sense that there are constants $`c_1,d>0`$ such that
$$\left|\frac{d}{du}f\right|c_1e^{du}|f|.$$
(3.23)
Since $`X`$ is a nontrivial solution, $`|f|0`$. Thus we can divide (3.23) by $`|f|`$ and integrate from any $`u<u_2`$ to $`u_2`$ to obtain
$$\mathrm{log}|f||_u^{u_2}c_2e^{du}|_u^{u_2}$$
with $`c_2=c_1/d`$. Since the right side of this inequality stays finite when $`u\mathrm{}`$, we conclude that there is a constant $`L>0`$ with
$$\frac{1}{L}|f(u)|L\text{for all }u<u_2.$$
(3.24)
Using that $`\lambda `$ depends smoothly on $`\omega `$ (see the Appendix), the constants $`c_1,c_2,d`$, and $`L`$ clearly can be chosen locally uniformly in $`\omega `$.
We substitute (3.24) into (3.23),
$$\left|\frac{d}{du}f\right|c_1Le^{du}.$$
(3.25)
This inequality shows that $`f^{}`$ is integrable, and thus $`f(u)`$ converges for $`u\mathrm{}`$. Setting
$$f_0=\underset{u\mathrm{}}{lim}f(u)\stackrel{(\text{3.24})}{}\mathrm{\hspace{0.33em}0},$$
we can integrate (3.25) from $`\mathrm{}`$ to $`u<b`$ and get
$$|f(u)f_0|ce^{du}$$
with $`c=c_1L/d`$. Substituting in the ansatz (3.12), we get (3.15).
¿From (3.13) we see that $`X(u)`$ does not decay to zero for $`u\mathrm{}`$. As a consequence, the function $`\mathrm{\Psi }_{u_2}^{k\omega n}`$ cannot have finite norm and thus is not a vector in the Hilbert space $`_{u_2}`$. This shows that the Hamiltonian $`H_{u_2}`$ has no point spectrum. In contrast to (3.3), we normalize the functions $`\mathrm{\Psi }_{u_2}^{k\omega n}`$ according to
$$\underset{u\mathrm{}}{lim}|X_{u_2}^{k\omega n}|=\mathrm{\hspace{0.33em}1},(Y^{k\omega n}|Y^{k\omega n})=\mathrm{\hspace{0.33em}1}.$$
(3.26)
The next two lemmas describe the behavior of the normalization factors and the energy gaps as $`u_1\mathrm{}`$.
###### Lemma 3.2.
For fixed $`u_2`$ and asymptotically as $`u_1\mathrm{}`$,
$`X_{u_1,u_2}^{k\omega n}`$ $`=`$ $`g(u_1)X_{u_2}^{k\omega n}|_{[u_1,u_2]}\text{with}`$ (3.27)
$`|g(u_1)|^2`$ $`=`$ $`(u_2u_1)+𝒪(1).`$ (3.28)
Furthermore,
$$|<\text{ }\mathrm{\Psi }_{u_1,u_2}^{k\omega n}|\mathrm{\Psi }_{u_1,u_2}^{k\omega n^{}}>\delta ^{nn^{}}|\frac{c}{u_2u_1}<\text{ }Y^{k\omega n}|\mathrm{sin}\vartheta \sigma ^1|Y^{k\omega n^{}}>,$$
(3.29)
where the constant $`c`$ can be chosen locally uniformly in $`\omega `$.
Proof. Since $`X_{u_1,u_2}^{k\omega n}`$ and $`X_{u_2}^{k\omega n}`$ are solutions of the same ODE (3.22) with the same boundary conditions at $`u_2`$, they clearly coincide up to a normalization factor $`g`$, i.e. $`X_{u_1,u_2}^{k\omega n}=gX_{u_2}^{k\omega n}`$. Taking the norm on both sides and using the first part of (3.3), we obtain that
$$|g(u_1)|^2=_{u_1}^{u_2}\overline{X_{u_2}^{k\omega n}}(u)X_{u_2}^{k\omega n}(u)𝑑u.$$
We now substitute (3.12), multiply out, and use that $`|f_0|^2=1`$ according to the first part of (3.26), to obtain
$$|g(u_1)|^2=_{u_1}^{u_2}\left(1+\overline{X}R_0+\overline{R_0}X|R_0|^2\right)𝑑u.$$
(3.30)
Since $`X`$ is bounded and $`R_0`$ has exponential decay (3.15), the last three summands in (3.30) are integrable, and thus $`|g(u_1)|^2(u_2u_1)`$ is bounded uniformly in $`u_1`$. This proves (3.28).
The scalar product $`<\text{ }\mathrm{\Psi }_{u_1,u_2}^{k\omega n}|\mathrm{\Psi }_{u_1,u_2}^{k\omega n^{}}>`$ can be computed via (2.50). The orthonormality (3.4) yields that
$`<\text{ }\mathrm{\Psi }_{u_1,u_2}^{k\omega n}|\mathrm{\Psi }_{u_1,u_2}^{k\omega n^{}}>\delta ^{nn^{}}`$ (3.31)
$`=`$ $`a(X_{u_1,u_2}^{k\omega n}|{\displaystyle \frac{\sqrt{\mathrm{\Delta }}}{r^2+a^2}}\sigma ^2|X_{u_1,u_2}^{k\omega n^{}})(Y^{k\omega n}|\mathrm{sin}\vartheta \sigma ^1|Y^{k\omega n^{}}).`$
In order to estimate the radial scalar product, we first note that the factor $`\sqrt{\mathrm{\Delta }}`$ goes exponentially to zero for $`u\mathrm{}`$, and thus
$$(X_{u_1,u_2}^{k\omega n}|\frac{\sqrt{\mathrm{\Delta }}}{r^2+a^2}\sigma ^2|X_{u_1,u_2}^{k\omega n^{}})c_4_{u_1}^{u_2}e^{du}|X_{u_1,u_2}^{k\omega n}||X_{u_1,u_2}^{k\omega n^{}}|𝑑u$$
for some constant $`c_4>0`$. Substituting (3.28) and using that the integral is uniformly bounded due to the factor $`\mathrm{exp}du`$, we obtain the estimate
$$\left|(X_{u_1,u_2}^{k\omega n}|\frac{\sqrt{\mathrm{\Delta }}}{r^2+a^2}\sigma ^2|X_{u_1,u_2}^{k\omega n^{}})\right|c_5(u_2u_1)^1,$$
which together with (3.31) yields (3.29).
###### Lemma 3.3.
The following estimate holds asymptotically as $`u_1\mathrm{}`$,
$$\mathrm{\Delta }\omega _{kn}=\frac{\pi }{u_2u_1}+𝒪((u_2u_1)^2),$$
(3.32)
for fixed $`u_2`$ locally uniformly in $`\omega `$.
Proof. We consider solutions of (3.11) satisfying the boundary conditions at $`u_2`$ and ask for which values of $`\omega `$ and $`\lambda _n(\omega )`$ our boundary conditions are also fulfilled at $`u_1`$. As is immediately verified from (3.11),
$$\frac{d}{du}\left(|X_+|^2|X_{}|^2\right)=\mathrm{\hspace{0.33em}0}.$$
Thus $`|X_+|^2|X_{}|^2`$ is independent of $`u`$, and since it vanishes at $`u_2`$,
$$|X_+|^2=|X_{}|^2\text{for all }uu_2.$$
(3.33)
Hence for the boundary values at $`u_1`$, we need not be concerned about the absolute values of $`X_\pm `$; it suffices to consider the condition for the phases
$$\mathrm{arg}X_+(u_1)=\mathrm{arg}X_{}(u_1).$$
(3.34)
It is convenient to work again with the ansatz (3.16). In order to describe the dependence on $`\omega `$, we differentiate (3.22) with respect to $`\omega `$. Since $`\lambda `$ depends smoothly on $`\omega `$, we obtain the bound
$$\left|\frac{d}{du}_\omega f\right|c_1e^{du}|_\omega f|+c_3e^{du}|f|$$
(3.35)
with constants $`c_1`$, $`d`$ as in (3.23) and $`c_3>0`$. Using that $`|f|`$ is bounded from above (3.24), we get
$$\left|\frac{d}{du}(|_\omega f|+c_4)\right|c_1e^{du}(|_\omega f|+c_4)$$
with $`c_4=c_3L/c_1`$. Similar to the development after (3.23), dividing by $`(|_\omega f|+c_4)`$ and integrating yields
$$\mathrm{log}(|_\omega f|+c_4)|_u^{u_2}c_2e^{du}|_u^{u_2},$$
and since the right side of this inequality is uniformly bounded in $`u`$,
$$|_\omega f|c_5$$
(3.36)
for some constant $`c_5>0`$.
For the study of the phase shifts, we introduce the phase function
$$\rho (u)=\mathrm{arg}f_+(u)\mathrm{arg}f_{}(u)\mathrm{\hspace{0.25em}2}\mathrm{\Omega }u_2$$
(the last summand was included so that $`\rho (u_2)=0`$). The derivative of the argument of a complex-valued function $`h`$ is given by
$$\frac{d}{du}\mathrm{arg}h(u)=\text{Im}\frac{h^{}(u)}{h(u)}.$$
Using this formula together with the fact that, according to (3.24) and (3.33) both $`|f_+|`$ and $`|f_{}|`$ are bounded away from zero, we obtain that
$$\left|\frac{d}{du}_\omega \rho \right|\mathrm{\hspace{0.33em}4}L\left|\frac{d}{du}_\omega f\right|+\mathrm{\hspace{0.25em}8}L^2|_\omega f|\left|\frac{d}{du}f\right|.$$
Substituting in (3.35), (3.25) as well as the bounds (3.24) and (3.36), we conclude that
$$\left|\frac{d}{du}_\omega \rho \right|c_6e^{du}$$
with some constant $`c_6>0`$. We integrate this inequality from $`u<u_2`$ to $`u_2`$. Since $`\rho (u_2)=0`$ independently of $`\omega `$, the boundary term $`_\omega \rho (u_2)`$ drops out, and we obtain the bound
$$\left|_\omega \rho (u)\right|C\text{for all }uu_2$$
(3.37)
with a constant $`C>0`$. This means that the equation for $`f`$, (3.22) leads only to finite phase shifts.
The boundary conditions at $`u_1`$, (3.34), are fulfilled iff
$$\mathrm{\Phi }:=\mathrm{\hspace{0.33em}2}\mathrm{\Omega }(u_2u_1)+\rho =\mathrm{\hspace{0.33em}0}(\text{mod }2\pi ).$$
Differentiating with respect to $`\omega `$ and integrating again from $`\omega _I`$ to $`\omega _{II}`$, $`\omega _I<\omega _{II}`$, we obtain that
$`\left|\mathrm{\Phi }(\omega _{II})\mathrm{\Phi }(\omega _I)\mathrm{\hspace{0.25em}2}(\omega _{II}\omega _I)(u_2u_1)\right|`$
$``$ $`{\displaystyle _{\omega _I}^{\omega _{II}}}|_\omega \rho |𝑑\omega \stackrel{(\text{3.37})}{}C(\omega _{II}\omega _I),`$
and this proves (3.32).
We can now prove the integral representation for the propagator of $`H_{u_2}`$.
###### Proposition 3.4.
For every $`\mathrm{\Psi }C_0^{\mathrm{}}((\mathrm{},u_2])\times S^2)^4`$ and $`x=(u,\vartheta ,\phi )`$,
$$\left(e^{itH_{u_2}}\mathrm{\Psi }\right)(x)=\frac{1}{\pi }\underset{k\text{Z}}{}_{\mathrm{}}^{\mathrm{}}𝑑\omega e^{i\omega t}\underset{n\text{Z}}{}\mathrm{\Psi }_{u_2}^{k\omega n}(x)<\text{ }\mathrm{\Psi }_{u_2}^{k\omega n}|\mathrm{\Psi }>.$$
(3.38)
Proof. According to the bound (3.29), the operator $`A`$ in (3.6) converges uniformly in $`\omega `$ and $`k`$ to the identity as $`u_1\mathrm{}`$, and thus (3.5) simplifies asymptotically to
$`\left(e^{itH_{u_1,u_2}}\mathrm{\Psi }\right)(x)`$
$`=`$ $`{\displaystyle \underset{k\text{Z}}{}}{\displaystyle \underset{\omega \sigma (H_{u_1,u_2}^k)}{}}e^{i\omega t}{\displaystyle \underset{nN(k,\omega )}{}}\mathrm{\Psi }_{u_1,u_2}^{k\omega n}<\text{ }\mathrm{\Psi }_{u_1,u_2}^{k\omega n}|\mathrm{\Psi }>+𝒪((u_2u_1)^1).`$
Using (3.27) and (3.28), we can express $`\mathrm{\Psi }_{u_1,u_2}^{k\omega n}`$ by $`\mathrm{\Psi }_{u_2}^{k\omega n}`$,
$`\left(e^{itH_{u_1,u_2}}\mathrm{\Psi }\right)(x)={\displaystyle \underset{k\text{Z}}{}}{\displaystyle \frac{1}{u_2u_1}}{\displaystyle \underset{\omega \sigma (H_{u_1,u_2}^k)}{}}e^{i\omega t}`$
$`\times {\displaystyle \underset{nN(k,\omega )}{}}\mathrm{\Psi }_{u_1}^{k\omega n}<\text{ }\mathrm{\Psi }_{u_1}^{k\omega n}|\mathrm{\Psi }>_{u_1,u_2}+𝒪((u_2u_1)^1).`$
The gap estimate, Lemma 3.3, shows that the sum over the spectrum is a Riemann sum which converges as $`u_1\mathrm{}`$ to an integral.
The idea for proving an integral representation for $`\mathrm{exp}(itH)`$ is to take in (3.38) a suitable limit $`u_2+\mathrm{}`$. In preparation, we need to derive estimates which describe the asymptotics of solutions of the radial equation (3.11) for large $`u`$. In this regime,
$`{\displaystyle \frac{d}{du}}X`$ $`=`$ $`\left[\left(\begin{array}{cc}i\omega & im\\ im& i\omega \end{array}\right)+{\displaystyle \frac{1}{u}}\left(\begin{array}{cc}ieQ& imM\lambda \\ imM\lambda & ieQ\end{array}\right)\right]X`$ (3.44)
$`+𝒪(u^2)X.`$
Thus the matrix potential on the right converges for $`u\mathrm{}`$. If $`|\omega |<m`$, its eigenvalues $`\lambda =\pm \sqrt{m^2\omega ^2}`$ are real, and this leads to one fundamental solution of (3.44) which decays exponentially like $`\mathrm{exp}(\sqrt{m^2\omega ^2}u)`$, and the other solution has exponential growth $`\mathrm{exp}(\sqrt{m^2\omega ^2}u)`$. We denote these two fundamental solutions by $`\mathrm{\Psi }_1^{k\omega n}`$ and $`\mathrm{\Psi }_2^{k\omega n}`$, respectively, and normalize them according to
$$\underset{u\mathrm{}}{lim}|\mathrm{\Psi }_{1/2}^{k\omega n}(u)|=\mathrm{\hspace{0.33em}1},$$
(3.45)
where our notation $`1/2`$ means that the above equation is valid in both cases $`1`$ and $`2`$. For $`|\omega |>m`$, on the other hand, the eigenvalues of the matrix potential at $`u=\mathrm{}`$ are imaginary, $`\lambda =\pm i\sqrt{\omega ^2m^2}`$, and this leads to two fundamental solutions $`\mathrm{\Psi }_{1/2}^{k\omega n}`$ with oscillatory behavior $`\mathrm{exp}(\pm i\sqrt{\omega ^2m^2}u)`$. For the normalization, we are now free to choose both the amplitude and the phase. Our convention is that
$$f_{0,\mathrm{\hspace{0.25em}1}}^{k\omega n}=\left(\begin{array}{c}1\\ 0\end{array}\right)\text{and}f_{0,\mathrm{\hspace{0.25em}2}}^{k\omega n}=\left(\begin{array}{c}0\\ 1\end{array}\right)$$
(3.46)
with $`f_{0,\mathrm{\hspace{0.25em}1}/2}^{k\omega n}`$ as in the asymptotic expansion (3.12). The next lemma describes the asymptotics of the oscillatory solutions as $`u\mathrm{}`$.
###### Lemma 3.5.
Every nontrivial solution $`X`$ of (3.11) for $`|\omega |>m`$ has for large $`u`$ the asymptotic form
$$X(u)=A\left(\begin{array}{c}e^{i\mathrm{\Phi }(u)}f_{\mathrm{}}^+\\ e^{i\mathrm{\Phi }(u)}f_{\mathrm{}}^{}\end{array}\right)+R_{\mathrm{}}(u)$$
(3.47)
with
$`f_{\mathrm{}}`$ $``$ $`0`$ (3.48)
$`\mathrm{\Phi }`$ $`=`$ $`ϵ(\omega )\left(\sqrt{\omega ^2m^2}u+{\displaystyle \frac{\omega eQ+Mm^2}{\sqrt{\omega ^2m^2}}}\mathrm{log}u\right)`$ (3.49)
$`A`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cosh}\mathrm{\Theta }& \mathrm{sinh}\mathrm{\Theta }\\ \mathrm{sinh}\mathrm{\Theta }& \mathrm{cosh}\mathrm{\Theta }\end{array}\right),\mathrm{\Theta }={\displaystyle \frac{1}{4}}\mathrm{log}\left({\displaystyle \frac{\omega +m}{\omega m}}\right)`$ (3.52)
$`|R_{\mathrm{}}|`$ $``$ $`{\displaystyle \frac{C}{u}}`$ (3.53)
and a constant $`C>0`$.
Proof. We write (3.11) symbolically as
$$X^{}=VX$$
with a matrix potential $`V(u)`$. According to (3.44) and the hypothesis $`|\omega |>m`$, the eigenvalues of $`V`$ are, for sufficiently large $`u`$, purely imaginary. More precisely, there is a transformation matrix $`B(u)`$ with
$$B^1VB=i\mathrm{\Omega }\sigma ^3$$
(3.54)
and a suitable function $`\mathrm{\Omega }(u)`$. Since the matrix potential $`V`$ converges for $`u\mathrm{}`$ and has a regular expansion in powers of $`1/u`$, we can choose $`B`$ such that
$$|B(u)|c_0,|B^{}(u)|\frac{c_0}{u^2}$$
(3.55)
with a constant $`c_0>0`$. The transformed wave function $`(B^1X)`$ satisfies the equation
$$\frac{d}{du}(B^1X)=\left[i\mathrm{\Omega }(u)\sigma ^3B^1B^{}\right](B^1X).$$
(3.56)
Hence employing the ansatz
$$X=B\left(\begin{array}{c}e^{i\mathrm{\Phi }}f^+(u)\\ e^{i\mathrm{\Phi }}f^{}(u)\end{array}\right)\text{with}\mathrm{\Phi }^{}(u)=\mathrm{\Omega }(u)$$
(3.57)
and using the bound (3.55), we obtain the inequality
$$\left|\frac{d}{du}f\right|\frac{c_0^2}{u^2}|f|.$$
(3.58)
A short calculation shows that $`\mathrm{\Phi }`$ has the explicit form (3.49), and that $`B(u)=A+𝒪(\frac{1}{u})`$ with $`A`$ according to (3.52). The term of order $`𝒪(\frac{1}{u})`$ can be absorbed into $`R_{\mathrm{}}`$.
The inequality (3.58) can be used similar to (3.23) in Lemma 3.1 Namely, dividing by $`|f|`$ and integrating yields for sufficiently large $`u`$ the bounds
$$\frac{1}{L}|f(u)|L.$$
(3.59)
After substituting the upper bound for $`|f|`$ into (3.58), one sees that $`f^{}`$ is integrable. Thus $`f`$ has a finite and, according to (3.59), non-zero limit,
$$f_{\mathrm{}}:=\underset{u\mathrm{}}{lim}f(u)\mathrm{\hspace{0.33em}0}.$$
Finally, the $`1/u`$-decay (3.53) follows by integrating (3.58) backwards from $`u=\mathrm{}`$ and employing the resulting bound in the ansatz (3.57).
In analogy to potential wall problems for Schrödinger operators, we call the function $`f_{\mathrm{}}`$ in (3.47) corresponding to our fundamental solutions $`\mathrm{\Psi }_{1/2}^{k\omega n}`$ the transmission coefficients, and denote them by $`f_{\mathrm{}\mathrm{\hspace{0.25em}1}/2}^{k\omega n}`$.
###### Theorem 3.6.
For every $`\mathrm{\Psi }C_0^{\mathrm{}}(\text{I R}\times S^2)^4`$,
$$\left(e^{itH}\mathrm{\Psi }\right)(x)=\frac{1}{\pi }\underset{k,n\text{Z}}{}_{\mathrm{}}^{\mathrm{}}𝑑\omega e^{i\omega t}\underset{a,b=1}{\overset{2}{}}t_{ab}^{k\omega n}\mathrm{\Psi }_a^{k\omega n}(x)<\text{ }\mathrm{\Psi }_b^{k\omega n}|\mathrm{\Psi }>,$$
(3.60)
where the coefficients $`t_{ab}`$ are for $`|\omega |<m`$ given by
$$t_{ab}=\delta _{a,1}\delta _{b,1}.$$
(3.61)
For $`|\omega |>m`$, the $`t_{ab}`$ are given by the integrals
$$t_{ab}=\frac{1}{2\pi }_0^{2\pi }\frac{t_a\overline{t_b}}{|t_1|^2+|t_2|^2}𝑑\alpha ,$$
(3.62)
where the functions $`t_a`$ are related to the transmission coefficients by
$$t_1(\alpha )=f_{\mathrm{}\mathrm{\hspace{0.25em}2}}^+e^{i\alpha }f_{\mathrm{}\mathrm{\hspace{0.25em}2}}^{}e^{i\alpha },t_2(\alpha )=f_{\mathrm{}\mathrm{\hspace{0.25em}1}}^+e^{i\alpha }+f_{\mathrm{}\mathrm{\hspace{0.25em}1}}^{}e^{i\alpha }.$$
(3.63)
The integral and the series in (3.60) converge in norm in the Hilbert space $``$.
Proof. Our strategy is as follows. Choosing $`u_2`$ so large that $`\text{supp}\mathrm{\Psi }(\mathrm{},u_2)`$, Proposition 3.4 yields for $`t=0`$ the “completeness relation”
$$\mathrm{\Psi }(x)=\frac{1}{\pi }\underset{k\text{Z}}{}_{\mathrm{}}^{\mathrm{}}𝑑\omega \underset{n\text{Z}}{}\mathrm{\Psi }_{u_2}^{k\omega n}(x)<\text{ }\mathrm{\Psi }_{u_2}^{k\omega n}|\mathrm{\Psi }>.$$
This formula remains true when $`u_2`$ is further increased,
$$\mathrm{\Psi }(x)=\frac{1}{\pi }\underset{k\text{Z}}{}_{\mathrm{}}^{\mathrm{}}𝑑\omega \underset{n\text{Z}}{}\mathrm{\Psi }_{u_2+\tau }^{k\omega n}(x)<\text{ }\mathrm{\Psi }_{u_2+\tau }^{k\omega n}|\mathrm{\Psi }>,\tau 0.$$
(3.64)
Hence we can take the average over $`\tau `$ in the finite interval $`[0,T]`$ with $`T>0`$ and obtain, using Fubini’s theorem,
$$\mathrm{\Psi }=\frac{1}{\pi }\underset{k\text{Z}}{}_{\mathrm{}}^{\mathrm{}}𝑑\omega \underset{n\text{Z}}{}\left[\frac{1}{T}_0^T𝑑\tau \mathrm{\Psi }_{u_2+\tau }^{k\omega n}<\text{ }\mathrm{\Psi }_{u_2+\tau }^{k\omega n}|\mathrm{\Psi }>\right].$$
(3.65)
We shall first prove that the square bracket in (3.65) has a finite limit as $`T\mathrm{}`$. Then we will show that for $`T\mathrm{}`$, we can in (3.65) take the limit inside the integral and the series in (3.65). This will give a decomposition of the identity in terms of eigensolutions of $`H`$, from which the representation of the propagator (3.60) will follow immediately by inserting the phase factors $`\mathrm{exp}(i\omega t)`$.
Let us analyze the square bracket in (3.65). We can write $`\mathrm{\Psi }_{u_2+\tau }^{k\omega n}`$ as a linear combination of the fundamental solutions $`\mathrm{\Psi }_{1/2}^{k\omega n}`$,
$$\mathrm{\Psi }_{u_2+\tau }^{k\omega n}(x)=\underset{a=1}{\overset{2}{}}c_a(\tau )\mathrm{\Psi }_a^{k\omega n}(x),$$
(3.66)
where the coefficients $`c_{1/2}`$ must be chosen such that our Dirichlet-type boundary conditions are satisfied at $`u_2+\tau `$. Then the square bracket becomes
$$\frac{1}{T}_0^T𝑑\tau \mathrm{\Psi }_{u_2+\tau }^{k\omega n}<\text{ }\mathrm{\Psi }_{u_2+\tau }^{k\omega n}|\mathrm{\Psi }>=\underset{a,b=1}{\overset{2}{}}t_{ab}(T)\mathrm{\Psi }_a^{k\omega n}<\text{ }\mathrm{\Psi }_b^{k\omega n}|\mathrm{\Psi }>$$
(3.67)
with
$$t_{ab}(T)=\frac{1}{T}_0^Tc_a(\tau )\overline{c_b(\tau )}𝑑\tau .$$
(3.68)
In the case $`|\omega |<m`$, $`\mathrm{\Psi }_1^{k\omega n}`$ and $`\mathrm{\Psi }_2^{k\omega n}`$ are for large $`u`$ exponentially decaying and increasing, respectively. Hence in order to fulfill the boundary conditions at $`u=u_2+\tau `$, the quotient $`c_2(\tau )/c_1(\tau )`$ must go exponentially to zero. Moreover, our normalization conditions (3.26) and (3.45) imply that $`|c_1(\tau )|^2`$ must tend to one. We conclude that there is a constant $`c_1`$ with
$$|c_a(\tau )\delta _{a,1}|c_1e^{\sqrt{m^2\omega ^2}\tau },$$
and so (3.68) converges for $`T\mathrm{}`$ to (3.61). In the case $`|\omega |>m`$, the fundamental solutions are oscillating for large $`u`$, as described by Lemma 3.5. The boundary conditions at $`u_2`$, (2.48), imply that the following scalar product must vanish,
$$<\text{ }\left(\begin{array}{c}f_1^+e^{i\mathrm{\Phi }(u_2+\tau )}f_1^{}e^{i\mathrm{\Phi }(u_2+\tau )}\\ f_2^+e^{i\mathrm{\Phi }(u_2+\tau )}f_2^{}e^{i\mathrm{\Phi }(u_2+\tau )}\end{array}\right)+𝒪(\tau ^1),\left(\begin{array}{c}c_1(\tau )\\ c_2(\tau )\end{array}\right)>=\mathrm{\hspace{0.33em}0},$$
(3.69)
where $`f_{1/2}`$ are the transmission coefficients. Moreover, the normalization and phase conditions (3.26) and (3.46) yield that
$$|c_1|^2+|c_2|^2=\mathrm{\hspace{0.33em}1}.$$
(3.70)
The general solution to (3.69) is
$$\left(\begin{array}{c}c_1\\ c_2\end{array}\right)=\frac{1}{D}\left(\begin{array}{c}f_2^+e^{i\mathrm{\Phi }(u_2+\tau )}f_2^{}e^{i\mathrm{\Phi }(u_2+\tau )}\\ f_1^+e^{i\mathrm{\Phi }(u_2+\tau )}+f_1^{}e^{i\mathrm{\Phi }(u_2+\tau )}\end{array}\right)+𝒪(\tau ^1)$$
(3.71)
with a complex parameter $`D`$, which can be chosen so as to satisfy the normalization condition (3.70). We now substitute (3.71) into (3.68) and take $`\mathrm{\Phi }`$ as the integration parameter,
$$t_{ab}(T)=\frac{1}{T}_{\mathrm{\Phi }(0)}^{\mathrm{\Phi }(T)}c_a(\mathrm{\Phi })\overline{c_b(\mathrm{\Phi })}\frac{d\mathrm{\Phi }}{|\mathrm{\Phi }^{}|}.$$
(3.72)
Using (3.49), one sees that (3.72) converges for $`T\mathrm{}`$ to the average over one period, giving (3.62) and (3.63). We conclude that the bracket in (3.65) converges pointwise as $`T\mathrm{}`$.
Next we shall prove that in (3.65) we may take the limit $`T\mathrm{}`$ inside the series and the integral, and that for the resulting limit the series and the integral converge in norm. The sum over $`k`$ in (3.65) gives the decomposition into the eigenspaces of the angular operator $`i_\phi `$, which we denote by $`^k`$. We may consider the situation on each such eigenspace separately, and thus assume that $`\mathrm{\Psi }^k`$. For the integral and the $`n`$-summation in (3.65) the situation is more difficult because the spectral decomposition of the Hamiltonian depends on $`u_2`$, and because the eigenvalues $`\lambda _n`$ of $`𝒜`$ and corresponding eigenspaces depend on $`\omega `$. We first apply to (3.64) the operator product $`𝒜^{2p}H^{2q}`$ with $`p,q0`$ (with $`𝒜`$ as in (2.6), where the $`t`$-derivative in $`𝒜`$ is carried out by applying the operator $`iH`$) and take the inner product with $`\mathrm{\Psi }`$. This gives
$$<\text{ }\mathrm{\Psi }|𝒜^{2p}H^{2q}\mathrm{\Psi }>=_{\mathrm{}}^{\mathrm{}}\frac{d\omega }{\pi }\omega ^{2q}\underset{n\text{Z}}{}\lambda _n^{2p}(\omega )|<\text{ }\mathrm{\Psi }_{u_2+\tau }^{k\omega n}|\mathrm{\Psi }>|^2.$$
(3.73)
If we consider on the right side of (3.73) instead of $`|<\text{ }\mathrm{\Psi }_{u_2+\tau }^{k\omega n}|\mathrm{\Psi }>|^2`$ a mixed product with $`\mathrm{\Phi },\mathrm{\Psi }C_0^{\mathrm{}}((\mathrm{},u_2)\times S^2)^4^k`$, we can in the integrand use the inequality $`xy\frac{1}{2}(x^2+y^2)`$, and then apply (3.73) to obtain
$`{\displaystyle \underset{n\text{Z}}{}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\omega }{\pi }}\omega ^{2q}\lambda _n^{2p}(\omega )|<\text{ }\mathrm{\Phi }|\mathrm{\Psi }_{u_2+\tau }^{k\omega n}><\text{ }\mathrm{\Psi }_{u_2+\tau }^{k\omega n}|\mathrm{\Psi }>|`$ (3.74)
$``$ $`{\displaystyle \frac{1}{2}}\left(<\text{ }\mathrm{\Phi }|𝒜^{2p}H^{2q}\mathrm{\Phi }>+<\text{ }\mathrm{\Psi }|𝒜^{2p}H^{2q}\mathrm{\Psi }>\right),`$
and this bound holds for all $`\tau 0`$.
Let $`\epsilon >0`$. Then for any $`\omega _0>0`$, (3.74) yields for $`p=0`$ and $`q=1`$ that
$`{\displaystyle \underset{n\text{Z}}{}}{\displaystyle _{\text{I R}[\omega _0,\omega _0]}}{\displaystyle \frac{d\omega }{\pi }}|<\text{ }\mathrm{\Phi }|\mathrm{\Psi }_{u_2+\tau }^{k\omega n}><\text{ }\mathrm{\Psi }_{u_2+\tau }^{k\omega n}|\mathrm{\Psi }>|`$
$``$ $`{\displaystyle \frac{1}{\omega _0^2}}{\displaystyle \underset{n\text{Z}}{}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\omega }{\pi }}\omega ^2|<\text{ }\mathrm{\Phi }|\mathrm{\Psi }_{u_2+\tau }^{k\omega n}><\text{ }\mathrm{\Psi }_{u_2+\tau }^{k\omega n}|\mathrm{\Psi }>|`$
$``$ $`{\displaystyle \frac{1}{2\omega _0^2}}\left(H\mathrm{\Phi }^2+H\mathrm{\Psi }^2\right).`$
We choose $`\omega _0`$ so large that
$$\underset{n\text{Z}}{}_{\text{I R}[\omega _0,\omega _0]}\frac{d\omega }{\pi }|<\text{ }\mathrm{\Phi }|\mathrm{\Psi }_{u_2+\tau }^{k\omega n}><\text{ }\mathrm{\Psi }_{u_2+\tau }^{k\omega n}|\mathrm{\Psi }>|<\frac{\epsilon }{2}$$
(3.75)
for all $`\tau 0`$. This inequality allows us to restrict attention to $`\omega `$ in the finite interval $`[\omega _0,\omega _0]`$. Next for a constant $`n_0>0`$, we consider (3.74) for $`p=1`$ and $`q=0`$. This gives the inequality
$`{\displaystyle \underset{|n|>n_0}{}}{\displaystyle _{\omega _0}^{\omega _0}}{\displaystyle \frac{d\omega }{\pi }}|<\text{ }\mathrm{\Phi }|\mathrm{\Psi }_{u_2+\tau }^{k\omega n}><\text{ }\mathrm{\Psi }_{u_2+\tau }^{k\omega n}|\mathrm{\Psi }>|`$
$``$ $`{\displaystyle \frac{1}{2}}\left(<\text{ }\mathrm{\Phi }|𝒜^2\mathrm{\Phi }>+<\text{ }\mathrm{\Psi }|𝒜^2\mathrm{\Psi }>\right)\underset{\omega [\omega _0,\omega _0],|n|>n_0}{sup}\lambda _n^2(\omega ).`$
Clearly $`\lambda _n^2(\omega )\mathrm{}`$ for $`n\pm \mathrm{}`$ uniformly in $`\omega [\omega _0,\omega _0]`$, and thus we can by choosing $`n_0`$ sufficiently large arrange that
$$\underset{|n|>n_0}{}_{\omega _0}^{\omega _0}\frac{d\omega }{\pi }|<\text{ }\mathrm{\Phi }|\mathrm{\Psi }_{u_2+\tau }^{k\omega n}><\text{ }\mathrm{\Psi }_{u_2+\tau }^{k\omega n}|\mathrm{\Psi }>|<\frac{\epsilon }{2}.$$
(3.76)
Putting together the estimates (3.75) and (3.76), we conclude that
$$(\underset{n\text{Z}}{}_{\mathrm{}}^{\mathrm{}}\underset{|n|>n_0}{}_{\omega _0}^{\omega _0})\frac{d\omega }{\pi }|<\text{ }\mathrm{\Phi }|\mathrm{\Psi }_{u_2+\tau }^{k\omega n}><\text{ }\mathrm{\Psi }_{u_2+\tau }^{k\omega n}|\mathrm{\Psi }>|<\epsilon $$
(3.77)
for all $`\tau 0`$.
For $`\omega `$ in a finite interval and $`n`$ in a finite set, we can take the average over $`\tau [0,T]`$ and take the limit $`T\mathrm{}`$ using Lebesgue’s dominated convergence theorem (notice that according to Lemma 3.1, $`\mathrm{\Psi }_{u_2+\tau }^{k\omega n}`$ is on the support of $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ uniformly bounded in $`\tau `$). This gives
$`\underset{T\mathrm{}}{lim}{\displaystyle \underset{n=n_0}{\overset{n_0}{}}}{\displaystyle _{\omega _0}^{\omega _0}}{\displaystyle \frac{d\omega }{\pi }}[{\displaystyle \frac{1}{T}}{\displaystyle _0^T}d\tau <\text{ }\mathrm{\Phi }|\mathrm{\Psi }_{u_2+\tau }^{k\omega n}><\text{ }\mathrm{\Psi }_{u_2+\tau }^{k\omega n}|\mathrm{\Psi }>]`$ (3.78)
$`=`$ $`{\displaystyle \underset{n=n_0}{\overset{n_0}{}}}{\displaystyle _{\omega _0}^{\omega _0}}{\displaystyle \frac{d\omega }{\pi }}{\displaystyle \underset{a,b=1}{\overset{2}{}}}t_{ab}^{k\omega n}<\text{ }\mathrm{\Phi }|\mathrm{\Psi }_a^{k\omega n}><\text{ }\mathrm{\Psi }_b^{k\omega n}|\mathrm{\Psi }>`$
with $`t_{ab}`$ according to (3.61) and (3.62). Since $`\epsilon `$ in (3.77) can be chosen arbitrarily small and $`n_0\mathrm{}`$, $`\omega _0\mathrm{}`$ as $`\epsilon 0`$, we obtain that (3.78) is true also for $`n_0=\mathrm{}=\omega _0`$, with absolute convergence of the integral and the series. Since $`\mathrm{\Phi }`$ can be chosen arbitrarily, we conclude that
$$\mathrm{\Psi }=\underset{T\mathrm{}}{lim}(\text{3.65})=\underset{k,n\text{Z}}{}_{\mathrm{}}^{\mathrm{}}\frac{d\omega }{\pi }\underset{a,b=1}{\overset{2}{}}t_{ab}^{k\omega n}\mathrm{\Psi }_a^{k\omega n}<\text{ }\mathrm{\Psi }_b^{k\omega n}|\mathrm{\Psi }>.$$
(3.79)
The estimate (3.74) for $`p=0=q`$ remains true if we take the average over $`\tau [0,\mathrm{})`$, and a homogeneity argument (as in the proof of the Schwarz inequality in Hilbert spaces) yields that
$$\underset{n\text{Z}}{}_{\mathrm{}}^{\mathrm{}}\frac{d\omega }{\pi }|\underset{a,b=1}{\overset{2}{}}t_{ab}<\text{ }\mathrm{\Phi }|\mathrm{\Psi }_a^{k\omega n}><\text{ }\mathrm{\Psi }_b^{k\omega n}|\mathrm{\Psi }>|\mathrm{\Phi }\mathrm{\Psi }.$$
This bound shows that the integral and series in (3.79) converge in norm, and that $`\mathrm{\Psi }`$ need not be an eigenvector of $`i_\phi `$. We finally apply the unitary operator $`\mathrm{exp}(itH)`$ on both sides of (3.79) to obtain (3.60).
Notice that the coefficients $`t_{ab}`$ given by (3.62) are bounded,
$$|t_{ab}|\frac{1}{2}\text{for }|\omega |>m\text{.}$$
(3.80)
In the asymptotic region $`u\mathrm{}`$, (3.60) goes over to a Fourier representation in terms of the plane-wave solution (3.12). A careful analysis of this limiting case gives additional information on the coefficients $`t_{ab}`$, namely
$$t_{11}=\frac{1}{2}=t_{22}\text{for }|\omega |>m\text{.}$$
(3.81)
However, the non-diagonal elements $`t_{12}`$ and $`t_{21}`$ remain undetermined. We shall not derive the relations (3.81) here, and will not use them in what follows.
## 4 The Decay Estimates
Using the integral representation for the propagator of the previous section, we can now prove the decay of the probabilities.
Proof of Theorem 1.1. According to the hypotheses of the theorem, the initial data $`\mathrm{\Psi }_0`$ is in $`L^2((r_1,\mathrm{})\times S^2,d\mu )^4`$. Since the transformation of the spinors (2.1) is smooth and involves factors $`\mathrm{\Delta }^{\frac{1}{4}}`$ and $`\sqrt{r}`$, we obtain for the transformed initial data $`\widehat{\mathrm{\Psi }}_0`$ that
$$r^{\frac{1}{2}}\mathrm{\Delta }^{\frac{1}{4}}\widehat{\mathrm{\Psi }}_0L^2((r_1,\mathrm{})\times S^2,d\mu )^4.$$
Equivalently, computing the volume element on the hypersurface $`t=const`$,
$$\frac{r^2}{\mathrm{\Delta }}|\widehat{\mathrm{\Psi }}_0|^2L^1((r_1,\mathrm{})\times S^2,drd\mathrm{cos}\vartheta d\phi )^4$$
Transforming to the variable $`u`$, (2.20), one sees that $`\widehat{\mathrm{\Psi }}_0`$ is in the Hilbert space with scalar product (2.43), $`\widehat{\mathrm{\Psi }}_0`$. For simplicity, we again omit the hat in what follows.
Let $`\epsilon >0`$. Since the wave functions with compact support are dense in $``$, then for this $`\epsilon `$, there is a $`\mathrm{\Psi }_IC^{\mathrm{}}((r_1,\mathrm{})\times S^2)^4`$ such that
$$\mathrm{\Psi }_I\mathrm{\Psi }_0<\epsilon .$$
(4.1)
For the Cauchy problem with initial data $`\mathrm{\Psi }_0`$, we have the integral representation of Theorem 3.6. Since the series in (3.60) converge in norm, we can choose $`k_0`$ and $`n_0`$ such that
$$\mathrm{\Psi }_{k_0,n_0}\mathrm{\Psi }_I\epsilon ,$$
(4.2)
where $`\mathrm{\Psi }_{k_0,n_0}`$ is defined by
$$\mathrm{\Psi }_{k_0,n_0}(x)=\frac{1}{\pi }\underset{k=k_0}{\overset{k_0}{}}\underset{n=n_0}{\overset{n_0}{}}_{\mathrm{}}^{\mathrm{}}𝑑\omega \underset{a,b=1}{\overset{2}{}}t_{ab}^{k\omega n}\mathrm{\Psi }_a^{k\omega n}(x)<\text{ }\mathrm{\Psi }_b^{k\omega n}|\mathrm{\Psi }_I>.$$
(4.3)
Consider the integrand in (4.3) for fixed $`k`$ and $`n`$,
$$\underset{a,b=1}{\overset{2}{}}t_{ab}^{k\omega n}\mathrm{\Psi }_a^{k\omega n}(x)<\text{ }\mathrm{\Psi }_b^{k\omega n}|\mathrm{\Psi }_I>.$$
(4.4)
¿From (3.80) and the estimates of Lemma 3.1, one sees that (4.4) is bounded, locally uniformly in $`x`$ and $`\omega `$. Thus the norm convergence established in Theorem 3.6 implies that (4.4) is in $`L^1(\text{I R},\text{}\text{C}^4)`$ as a function of $`\omega `$, with an $`L^1`$-bound locally uniform in $`x`$. Hence its Fourier transform is $`L^{\mathrm{}}`$ in $`t`$, locally uniformly in $`x`$. Furthermore, the Riemann-Lebesgue lemma yields that its Fourier transform tends to zero as $`t\mathrm{}`$, pointwise in $`x`$. Since (4.3) involves only finitely many terms, we conclude that the solution of the Cauchy problem with initial data $`\mathrm{\Psi }_{k_0,n_0}`$,
$$\mathrm{\Psi }_{k_0,n_0}(t,x)=\frac{1}{\pi }\underset{k=k_0}{\overset{k_0}{}}\underset{n=n_0}{\overset{n_0}{}}_{\mathrm{}}^{\mathrm{}}𝑑\omega e^{i\omega t}\underset{a,b=1}{\overset{2}{}}t_{ab}^{k\omega n}\mathrm{\Psi }_a^{k\omega n}(x)<\text{ }\mathrm{\Psi }_b^{k\omega n}|\mathrm{\Psi }_I>,$$
(4.5)
is $`L^{\mathrm{}}`$ in $`t`$ locally uniformly in $`x`$, and $`lim_t\mathrm{}\mathrm{\Psi }_{n_0,k_0}(t,x)=0`$ for all $`x`$.
Choose $`K_{\delta ,R}`$ as in the statement of the theorem. Since the metric and Dirac matrices in the probability integral (1.3) are smooth and bounded on the compact set $`K_{\delta ,R}`$, the corresponding bilinear form is continuous on $``$, i.e. there is a constant $`c`$ depending only on $`\delta `$ and $`R`$ such that for all $`\mathrm{\Psi }_1,\mathrm{\Psi }_2`$,
$$_{K_{\delta ,R}}(\overline{\mathrm{\Psi }_1}\gamma ^j\mathrm{\Psi }_2)\nu _j𝑑\mu c\mathrm{\Psi }_1\mathrm{\Psi }_2.$$
(4.6)
The solution to our original Cauchy problem is obtained by applying the unitary operator $`\mathrm{exp}(itH)`$ to $`\mathrm{\Psi }_0`$,
$`\mathrm{\Psi }(t)`$ $`=`$ $`e^{itH}\mathrm{\Psi }_0`$
$`=`$ $`e^{itH}\mathrm{\Psi }_{k_0,n_0}+e^{itH}(\mathrm{\Psi }_I\mathrm{\Psi }_{k_0,n_0})+e^{itH}(\mathrm{\Psi }_0\mathrm{\Psi }_I)`$
$`=`$ $`\mathrm{\Psi }_{k_0,n_0}(t)+e^{itH}(\mathrm{\Psi }_I\mathrm{\Psi }_{k_0,n_0})+e^{itH}(\mathrm{\Psi }_0\mathrm{\Psi }_I),`$
where $`\mathrm{\Psi }_{k_0,n_0}(t)`$ has the integral representation (4.5). We substitute this formula for $`\mathrm{\Psi }(t)`$ into the probability integral, multiply out, and apply the estimate (4.6) as well as the unitarity of $`\mathrm{exp}(itH)`$ together with (4.1) and (4.2). This gives the inequality
$`{\displaystyle _{K_{\delta ,R}}}(\overline{\mathrm{\Psi }}\gamma ^j\mathrm{\Psi })(t,x)\nu _j𝑑\mu `$
$``$ $`{\displaystyle _{K_{\delta ,R}}}(\overline{\mathrm{\Psi }_{k_0,n_0}}\gamma ^j\mathrm{\Psi }_{k_0,n_0})(t,x)\nu _j𝑑\mu +\mathrm{\hspace{0.25em}4}c^2\epsilon ^2+\mathrm{\hspace{0.25em}4}c\epsilon \mathrm{\Psi }.`$
We showed above that the integrand in the last integral is uniformly bounded and tends to zero pointwise as $`t\mathrm{}`$. Thus the integral converges to zero according to Lebesgue’s dominated convergence theorem.
We remark that in the spherically symmetric case, the analytical method given above to prove that $`\widehat{\mathrm{\Psi }_0}`$ is an alternative to the nice geometric argument by Kay and Wald , who use the causal propagation property and a discrete symmetry of the maximally extended space-time at the bifurcation $`2`$-sphere.
## Appendix A Nondegeneracy and Regularity of the Angular Eigenfunctions
In this appendix, we shall consider the angular equations (2.16),(2.18). As explained in \[1, Appendix A\], it is useful to write (2.16) as an eigenvalue equation in $`\lambda `$,
$$𝒜Y=\lambda Y\text{with}𝒜=\left(\begin{array}{cc}am\mathrm{cos}\vartheta & _{}\\ _+& am\mathrm{cos}\vartheta \end{array}\right).$$
(A.1)
###### Proposition A.1.
For given $`k`$ and $`\lambda \sigma (𝒜)`$, there is at most one eigensolution of (A.1), which we denote by $`Y^k`$, i.e.
$$𝒜Y^k=\lambda Y^k.$$
(A.2)
By continuously varying the parameter $`\omega `$, the eigenvalue equation (A.2) can be extended to all values of $`\omega \text{I R}`$. Both $`\lambda `$ and $`Y^k`$ depend smoothly on $`\omega `$.
Proof. The two fundamental solutions of (A.2) behave near $`\vartheta =0`$ like
$$Y^k=(\vartheta ^k+o(\vartheta ^k),o(\vartheta ^k))\text{and}Y^k=(o(\vartheta ^{k1}),\vartheta ^{k1}+o(\vartheta ^{k1})),$$
respectively. Depending on whether $`k`$ is $`0`$ or negative, the second or first fundamental solution diverges in the limit $`\vartheta 0`$. In \[1, Appendix A\] it was shown that the eigenfunctions $`Y^k`$ are bounded on $`S^2`$ and smooth except at the poles. Thus we can rule out one of the fundamental solutions and conclude that (A.2) has at most one solution.
Note that the solutions of (A.2) are the eigenvectors of $`𝒜`$ restricted to the eigenspace of the operator $`i_\phi `$ with eigenvalue $`k`$, which we denote by $`^k`$. Since the terms involving $`\omega `$ in (A.1) are a relatively compact perturbation, standard perturbation theory yields that the spectrum of $`𝒜_{|^k}`$ depends continuously on $`\omega `$. As no degeneracies occur, each eigenvalue $`\lambda `$ gives rise to a unique continuous family of eigenvalues $`\lambda (\omega )`$. Standard perturbation theory without degeneracies then yields that $`\lambda (\omega )`$ and the corresponding eigenvector $`Y^k(\omega )`$ depend smoothly on $`\omega `$.
Acknowledgments: We would like to thank McGill University, Montréal, and the Max Planck Institute for Mathematics in the Sciences, Leipzig, for support and hospitality. We are grateful to the referee for helpful suggestions.
The research of FK was supported by NSERC grant # RGPIN 105490-1998, of S in part by the NSF, Grant No. DMS-G-9802370, and of Y in part by the NSF, Grant No. 33-585-7510-2-30.
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# Low and High Ionization Absorption Properties of Mgii Absorption–Selected Galaxies at Intermediate Redshifts. I. General PropertiesBased in part on observations obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among Caltech, the University of California, and NASA. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. Based in part on observations obtained with the NASA/ESA Hubble Space Telescope, which is operated by the STScI for the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5–26555.
## 1. Introduction
During the last decade a great deal of progress has been made toward understanding the physical properties of intervening metal–line absorption systems measured in the spectra of high redshift quasars. This is particularly true at intermediate redshifts, $`0.5z1.5`$, for absorbers selected by the presence of the resonant Mgii $`\lambda \lambda 2796,2803`$ doublet (e.g. Lanzetta, Turnshek, & Wolfe (1987); Tytler et al. (1987); Sargent, Boksenberg, & Steidel (1988); Petitjean & Bergeron (1990), Steidel & Sargent (1992)). One of the most notable achievements was the demonstration that Mgii absorbers with $`W_r(2796)0.3`$ Å are almost always associated with galaxies (Bergeron & Boissé (1991); Steidel (1995)). Those works substantiated the $`30`$–year standing hypothesis by Bahcall & Spitzer (1969) that metal–line absorption in quasar spectra arises in extended gaseous envelopes surrounding intervening galaxies.
The general picture today is that a wide variety of morphological types (from ellipticals to irregulars) have gaseous “halos” that extend to roughly 40 kpc, with the most common being Sbc–Scd types (Steidel, Dickinson, & Persson 1994; Guillemin & Bergeron (1997)). The line–of–sight gas kinematics of the absorbers is consistent with that expected for material bound in galactic potential wells (Petitjean & Bergeron (1990); Churchill, Steidel, & Vogt 1996; Charlton & Churchill (1998)). This picture, however, is not without its counter examples or ambiguities. In some cases there is evidence that compact star forming objects spread out over $`200`$ kpc are seen at the Mgii absorption redshift and there is no directly associable bright galaxy (Yanny (1992); Yanny & York (1992)). It is also not yet established whether the more numerous “weak” Mgii absorbers, those with $`W_r(2796)<0.3`$ Å, are related to galaxies similar in type to those associated with “strong” Mgii absorption. There is mounting evidence that a fair number of the weak systems do not arise within $`40`$ kpc of normal, bright galaxies (Churchill & Le Brun (1998); Churchill et al. 1999a ).
What are the typical low to high ionization absorption conditions in intermediate redshift Mgii absorption–selected galaxies? Do the majority of Mgii systems have an associated high ionization phase as seen in Siiv, Civ, Nv, and Ovi absorption? Are there any trends between the high ionization and low ionization absorption strengths? Are there other relationships (or lack of relationships!) that provide clues to the physical nature of galactic gas at intermediate redshifts?
Motivated by these and similar questions, we have undertaken a program to measure the absorption properties of a wide variety of chemical and ionization species associated with Mgii absorbers. Unique to our study is that the Mgii systems have been observed at high resolution ($`6`$ km s<sup>-1</sup>) with HIRES/Keck I (Churchill (1997)). These spectra also provide a population of weak systems, which are significantly more numerous in their redshift path density (Churchill et al. 1999a ). The HIRES spectra cover Mgii, several Feii transitions, Mgi, and depending upon redshift coverage, Caii, Tiii, Mnii, and Aliii. The remaining absorption properties, including neutral hydrogen and higher ionization species, have been measured in lower resolution ($`230`$ km s<sup>-1</sup>) spectra obtained from the Hubble Space Telescope archive of the Faint Object Spectrograph.
In this paper, we present the measurements of the absorption lines found in the FOS spectra. Additional analysis focused on the above motivational questions is presented in a parallel companion paper (Churchill et al. (2000), hereafter Paper II). In § 2, we outline our sample selection. Details of the data analysis are presented in § 3. In § 4, we provide a brief description of each system. The general absorption properties are presented in § 5 and a brief synopsis is given in § 6.
## 2. Sample Selection
The Mgii systems for our study were selected from the samples of Steidel & Sargent (1992, hereafter SS92) and Sargent, Boksenberg, & Steidel (1988, hereafter SBS). We obtained HIRES spectra (Vogt et al. (1994)) on the Keck I telescope of many of the brightest quasars from the SS92 and SBS database. The optical wavelength coverage for our HIRES spectra translates to a redshift interval of $`0.4z1.4`$ for 28 Mgii systems. An additional 23 systems with equivalent widths below the detection threshold in the SS92 and SBS quasar spectra were discovered in the HIRES spectra, which have a detection threshold of $`0.02`$ Å (Churchill et al. 1999a ). In order to study the wealth of ionization species with transitions further into the ultraviolet, we then searched the Hubble Space Telescope archive for Faint Object Spectrograph (FOS) observations of the quasars observed with HIRES/Keck. A number of the FOS spectra were obtained fully reduced from the HST QSO Absorption Line Key Project (hereafter KP, Bahcall et al. (1993); Bahcall et al. (1996); Jannuzi et al. (1998)). The remaining FOS spectra were obtained from the HST archive and reduced using the same methodology as for the KP spectra.
The journal of the HIRES/Keck observations is given in Table Low and High Ionization Absorption Properties of Mgii Absorption–Selected Galaxies at Intermediate Redshifts. I. General Properties<sup>1</sup><sup>1</sup>affiliation: Based in part on observations obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among Caltech, the University of California, and NASA. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. <sup>2</sup><sup>2</sup>affiliation: Based in part on observations obtained with the NASA/ESA Hubble Space Telescope, which is operated by the STScI for the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5–26555.. The spectral resolution was $`R=45,000`$, which corresponds to $`6.6`$ km s<sup>-1</sup>. From left to right, the columns are the quasar name, the visual magnitude, the emission redshift, the observation date, the total exposure time, and the observed wavelength range. In the continuum near the observed Mgii features, the spectra have signal–to–noise ratios ranging from 15 to 50, with the majority being around 30. The spectra do not have continuous wavelength coverage; there are small gaps redward of 5100 Å, where the free spectral range of the projected echelle format exceeds the width of the $`2048\times 2048`$ Tektronix CCD.
There are several FOS observational modes, including grating, polarimetry, and slit settings. We have chosen to limit our survey to the highest resolution spectra ($`R=1300`$), obtained with the G130H (1150–1600 Å), G190H (1600–2300 Å) and G270H (2225–3275 Å) gratings. We chose to exclude spectra obtained at lower resolution in order to maintain uniformity and because they are not very useful for narrow–line absorption work. However, we did use a lower resolution G160L spectrum (in one case) to search for a Lyman limit break. We also opted to exclude spectropolarimetry data because they required additional reduction and calibration steps that did not yield spectra of comparable quality to those obtained with non–spectropolarimetry mode. However, when available, we did use $`R=1300`$ spectropolarimetry data to search for Lyman limit breaks and damped Ly$`\alpha `$ lines at the absorber redshifts. Finally, we excluded spectra with slit widths greater than $`1`$″. Some of the FOS spectra in our sample were obtained before the COSTAR refurbishing mission. The instrumental spread function of pre–COSTAR spectra introduces broad wings in proportion to the aperture width of the observing mode (Jannuzi & Hartig (1994)). We studied only those FOS spectra obtained with fairly narrow apertures in order to mitigate systematics in the absorption line strengths between the pre–COSTAR and COSTAR spectra.
The final selection of FOS spectra are listed in Table Low and High Ionization Absorption Properties of Mgii Absorption–Selected Galaxies at Intermediate Redshifts. I. General Properties<sup>1</sup><sup>1</sup>affiliation: Based in part on observations obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among Caltech, the University of California, and NASA. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. <sup>2</sup><sup>2</sup>affiliation: Based in part on observations obtained with the NASA/ESA Hubble Space Telescope, which is operated by the STScI for the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5–26555.. From left to right the columns are the quasar name, the alias useful for searching the HST archive, the program identification numbers, and the names of the principal investigators for each of the gratings. The observing modes of the KP spectra are described in Jannuzi et al. (1998); the slit widths are all less than $`0.26`$″. The observing modes of the non–KP spectra vary and are noted in § 4, where each spectrum is individually discussed. We also excluded quasars with complicated broad absorption line features. We make no corrections to these narrow aperture pre–COSTAR spectra and assume a Gaussian instrumental spread function (e.g. Schneider et al. (1993)) for measuring equivalent widths in all FOS spectra.
These selection criteria resulted in a sample of 45 Mgii absorbing systems. These systems and their rest–frame equivalent widths are listed in Table Low and High Ionization Absorption Properties of Mgii Absorption–Selected Galaxies at Intermediate Redshifts. I. General Properties<sup>1</sup><sup>1</sup>affiliation: Based in part on observations obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among Caltech, the University of California, and NASA. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. <sup>2</sup><sup>2</sup>affiliation: Based in part on observations obtained with the NASA/ESA Hubble Space Telescope, which is operated by the STScI for the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5–26555.. The first three columns are the quasar name, the absorber redshift, and the Mgii $`\lambda 2796`$ rest–frame equivalent width, $`W_r(\text{Mg}\text{ii})`$. The remaining columns are explained in § 4. In Figure 1, we show a plot of $`W_r(\text{Mg}\text{ii})`$ vs. absorber redshift, $`z_{\mathrm{abs}}`$. Note that the sample is devoid of $`W_r(\text{Mg}\text{ii})0.5`$ Å systems for $`z1`$. This is not surprising because the total redshift path above $`z1`$ covered by the HIRES spectra is very small compared to the coverage from $`0.4z1.0`$ and because small equivalent width systems are very common.
The distribution of $`W_r(\text{Mg}\text{ii})`$ is shown in Figure 2. Using a maximum likelihood fit (e.g. Lanzetta, Turnshek, & Wolfe (1987)), we find that the distribution of our sample follows a power–law with $`f(W)W^{0.9\pm 0.6}`$ for $`0.03W_r(\text{Mg}\text{ii})<1.3`$ Å, which is consistent with that of an unbiased sample, $`f(W)W^{1.0\pm 0.1}`$, over this equivalent width range (Churchill et al. 1999a ). Between $`1.0`$$`1.3`$ Å the data are systematically below the fit, but are consistent with the maximum likelihood result to $`1\sigma `$ in all bins. Above $`1.3`$ Å, there is a slight overabundance of systems (not included in the maximum likelihood analysis), due to a selection bias toward larger equivalent width absorbers in the HIRES/Keck survey (Churchill (1997)). These are all damped Ly$`\alpha `$ absorbers. Overall, the sample studied here is consistent with an unbiased sample in its equivalent width distribution for $`W_r(\text{Mg}\text{ii})<1.3`$ Å and for $`z1`$.
## 3. Data Analysis
The HIRES data were processed in the standard manner using the IRAF<sup>1</sup><sup>1</sup>1IRAF is distributed by the National Optical Astronomy Observatories, which are operated by AURA, Inc., under contract to the NSF. Apextract package for echelle data. The details of the HIRES data reduction are given in Churchill (1995, 1997). The KP FOS spectra were reduced by the KP collaboration, as described in Schneider et al. (1993), Bahcall et al. (1996), and Jannuzi et al. (1998). The archival FOS spectra were reduced and calibrated using KP methods.
As a consistency check, we performed two separate, full analyses of the FOS data. The primary analysis invoked a priori knowledge of the Mgii redshifts for searching, identifying, and measuring the absorption lines (described in greater detail below). This analysis also used a priori knowledge of other metal–line systems along the line of sight (from an exhaustive search of the literature), in order to assess the possibility of misidentifications and blends from these other systems. We developed our own set of automated and graphically interactive software routines for objective feature finding, identifying blending due to other systems, and measuring equivalent widths (including Gaussian deblending). Many of our algorithms are based upon those described by Schneider et al. (1993) and Bahcall et al. (1996). The secondary analysis, which we call the “KP analysis”, was performed as a consistency check. This was an unbiased search for metal–line systems using the KP methodologies (Bahcall et al. (1993); Bahcall et al. (1996); Jannuzi et al. (1998)) and software, i.e. the ZSEARCH and JASON programs. The two full analyses were compared for each FOS absorption line measurement. Though the “KP analysis” yielded different results in detail from those of the primary analysis (i.e. small differences in equivalent widths and their errors), there were very few inconsistencies. In the few cases of differing results, both analyses were studied and a consensus was reached.
### 3.1. Continuum Fitting
Since the continuum fitting is the most subjective aspect of the data reduction and probably has the most significant impact on the analysis, we briefly elaborate. For the HIRES data, continuum fitting is fairly straight forward because the continuum is well sampled and the signal–to–noise ratio is quite high. We used the technique of Sembach & Savage (1992), which employs Legendre polynomials across a limited regions of spectrum around each of the transitions. The HIRES spectra were not flux calibrated, so the normalized flux is simply relative to the flux levels resulting from the HIRES sensitivity function. Following flux calibration, the FOS spectra were continuum fit as described in Schneider et al. (1993); this entails several iterations of interactive cubic spline fitting to the flux calibrated spectra.
#### 3.1.1 Line Finding
For both the HIRES and the FOS spectra, absorption features were detected using a method slightly modified from that of Schneider et al. (1993), which is optimal in the case of unresolved lines. One constructs a discrete model of the instrumental spread function (ISF) consisting of $`M=2J_0+1`$ elements, where $`J_0`$ is a non–negative integer and $`_1^MP_j=1`$. The term $`P_j`$ is the value at pixel $`j`$ of a symmetric ISF having a peak value at $`j=J_0+1`$. For both HIRES and FOS, the ISF is modeled as a Gaussian with $`R=45,000`$ and $`1300`$, respectively. We use $`J_0=6`$, which gives a model ISF over 13 pixels.
The equivalent width of an unresolved feature centered in pixel $`i`$ is calculated by centering the ISF on the pixel and then computing
$$w(\lambda _i)=\mathrm{\Delta }\lambda _i\frac{_{j_1}^{j_2}P_jF(\lambda _k)}{_{j_1}^{j_2}P_j^2,}$$
(1)
where $`F(\lambda _k)=1I(\lambda _k)/I_c(\lambda _k)`$, $`I(\lambda _k)`$ is the flux in pixel $`k`$, $`I_c(\lambda _k)`$ is its fitted continuum level, $`\mathrm{\Delta }\lambda _i=(\lambda _{i1}\lambda _{i+1})/2`$ is the wavelength interval spanned by the pixel, $`j_1`$ and $`j_2`$ are the minimum and the maximum points in the ISF, and the index<sup>2</sup><sup>2</sup>2The index $`k`$, as written here, provides a minor correction to a previously published version of this formula in Schneider et al. (1993). $`k=i+(j1)J_0`$, and where $`w(\lambda _i)<0`$ for an absorption feature. At the spectra ends the continuum is extrapolated over $`J_0`$ pixels (similar to a wrap around convolution). The uncertainty in $`w(\lambda _i)`$, is given by
$$\sigma _w(\lambda _i)=\mathrm{\Delta }\lambda _i\frac{\left\{_{j_1}^{j_2}P_j^2\sigma _F^2(\lambda _k)\right\}^{1/2}}{_{j_1}^{j_2}P_j^2,}$$
(2)
where $`\sigma _F(\lambda _k)=\sigma _I(\lambda _k)/I_c(\lambda _k)`$, and where $`\sigma _I(\lambda _k)`$ is the $`1\sigma `$ uncertainty in $`I(\lambda _k)`$ resulting from data reduction and calibration sources and Poisson noise in the quasar flux and sky. This uncertainty “spectrum” serves as the $`1\sigma `$ observed equivalent width detection threshold.
In general terms, an unresolved absorption line at pixel $`i`$ is defined when the pixel equivalent width, $`w(\lambda _i)`$, is less than $`N\sigma _w(\lambda _i)`$, where $`N`$ is an arbitrarily defined number giving the number of $`\sigma `$ beyond the equivalent width detection threshold. For resolved features, the detection is defined over a spectral region. The region extremes are defined at the pixels where the $`w(\lambda _i)`$, which are smooth over the scale length of the ISF, become greater then zero.
For the HIRES spectra, we enforce a $`5\sigma `$ detection threshold. In Figure 3, we show the cumulative distribution of the $`5\sigma `$ rest–frame equivalent width detection threshold of Mgii $`\lambda 2796`$. Our sample is 100% complete to a $`5\sigma `$ threshold of $`0.06`$ Å, 93% complete to $`0.03`$ Å, and 73% complete to $`0.02`$ Å. For the FOS spectra, we require only a $`3\sigma `$ detection threshold (whereas the “KP analysis” enforced a $`4.5\sigma `$ threshold). We applied a less stringent threshold because we used a priori knowledge of the expected location of the absorption lines, whereas the “KP analysis” was an unbiased search for absorption lines, and was therefore more conservative.
#### 3.1.2 Equivalent Widths
For the HIRES spectra, the equivalent widths, $`W`$, are measured directly by summing the quantity $`[1I(\lambda _i)/I_c(\lambda _i)]\mathrm{\Delta }\lambda _i`$ across the profiles. The measurement uncertainties are obtained from quadrature summing the quantity $`\sigma _I(\lambda _i)[W/I(\lambda _i)]`$. The equivalent width uncertainties do not account for subjectivity in the continuum placement.
For the FOS spectra, the equivalent widths and their uncertainties were measured by fitting Gaussians to the absorption features. The quoted equivalent widths and uncertainties were taken from our own measurements (i.e. not from the “KP analysis”). We have used an interactive $`\chi ^2`$ minimization scheme, where the minimization is performed by the NETLIB–slatec routine<sup>3</sup><sup>3</sup>3NETLIB is a collection of mathematical software, papers, and databases maintained by AT&T Bell Laboratories, the University of Tennessee, and Oak Ridge National Laboratory (www.netlib.org). dnls1 (More (1978)). The uncertainties in the fitted parameters, the Gaussian amplitudes, widths, and centers, are computed using a modified version of the routine dfridr (Press et al. (1992)). The equivalent width uncertainties are based upon standard error propagation, including correlated terms. For a given line, the minimum allowed Gaussian width is set by the instrumental resolution.
Due to the number density of absorption features and the limited resolution of the FOS spectra, Gaussian deblending was sometimes required. Deblending was used only in cases where the individual line centroids were clearly separated or where a weak blend in the wing of a stronger line showed clear “double” structure. Otherwise, a blend was quoted. Examples of the deblending cases are illustrated in Figure 4, which shows the Siiv doublet at $`z=0.9902`$ in PG $`1634+706`$; $`\lambda 1393`$ lies between two Ly$`\alpha `$ lines and $`\lambda 1402`$ resides in the wing of Galactic Mgii $`\lambda 2796`$.
#### 3.1.3 Wavelength Zero Point Shifting of FOS Spectra
The velocity zero points of the KP FOS spectra were defined by setting the mean of the singly ionized transitions from Galactic clouds to redshift zero (see Savage et al. (1993)). The archival FOS spectra, on the other hand, were not zero point shifted immediately following their reduction. In both the KP and archival spectra, analysis of the different absorption lines associated with the Mgii systems often revealed a systematic velocity shift with respect to the HIRES spectra, which yielded very precise ($`\sigma _z/z10^6`$), heliocentric, vacuum wavelength, absorption redshifts from Mgii.
We adjusted the zero points of the FOS spectra using a constant velocity shift, which was often determined from the mean velocity offsets of two or more Mgii systems along the line of sight. Since Siii and Cii have ionization states similar to Mgii, they were used to derive the shift. If neither Siii nor Cii was available, then we used Ly$`\alpha `$ and/or Civ. The largest shift we applied was no larger than $`70`$% of a FOS spectral resolution element.
#### 3.1.4 Line Identification Policy in FOS Spectra
Our philosophy for the detection and identification of lines in the crowded FOS spectra is as follows \[also see § 4 of Bahcall et al. (1996)\]. We generated an objective list of lines using a $`3\sigma `$ detection criterion (see § 3.1.1). This relatively “liberal” threshold was chosen because we are not producing a formal catalog of all absorbers, but have precise redshifts from Mgii in the HIRES spectra. In addition, the continuum fits to the FOS spectra were conservative in that they are probably systematically low in regions of dense line blending.
We examined, case by case, those lines that happen to be coincident with the predicted positions of the transitions from Si, Oi, Ni, Feii, Siii, Alii, Sii, Cii, Aliii, Nii, Feiii, Siiii, Siii, Siiv, Siv, Niii, Ciii, Civ, Svi, Nv, Ovi, and the Lyman series. For several of the above species, multiple transitions were covered, allowing either examination of a “clean” region of spectrum, or data that corroborated or ruled out identifications in “confused” regions. Some of these transitions were rarely detected at the $`3\sigma `$ level. We did not search regions of the spectra having an observed equivalent width detection threshold greater than $`3`$ Å.
A fair fraction of the time, line blending with transitions from other absorbing systems (some from our sample) was problematic. We attempted to identify all lines involved with blends as follows: We constructed a list of 100 transitions with accurate rest–frame wavelengths (Morton (1991)) and a list of all known absorption redshifts, including Ly$`\alpha `$ “systems” for which Ly$`\beta `$ could be confirmed, from an exhaustive literature search (optical and ultraviolet observations) for known systems. In most all cases, a viable candidate for a blend could be determined.
For a given system, a species/transition is identified under the following conditions: 1) there is at least one 3$`\sigma `$ detection that can be identified with a transition of the species; 2) all other covered transitions of the species are not inconsistent with the identification<sup>4</sup><sup>4</sup>4Additional transitions can be consistent if they: a) are not covered, b) have equivalent widths between the scaled $`f\lambda `$ value and the equivalent width of the strongest transitions, c) are in a blend for which the contribution of the transition of interest could be consistent.; 3) at least one transition is not blended with a possible transition from another system, or, if it is, can be unambiguously deblended using Gaussian fitting (meaning that it is on the wing of a line or has its own clear line center, see § 3.1.2). In the cases of doublets, if the weaker member of a doublet is not inconsistent with the first, then we quote it as a detection. This can occur when the second member is in a blend or is not covered by the spectrum. It is always possible that a doublet is a chance match with two “random” lines, but this type of false match is expected to have small probability in our spectra, based upon the simulations by Bahcall et al. (1996) and by Jannuzi et al. (1998).
If a species was detected in fewer than 10% of the systems, it was eliminated from our overall presentation. Each species has what we call a “flag transition”. The flag transitions are those with the strongest $`f\lambda `$ for each chemical/ionization species. For the non–doublet transitions, they are Feii $`\lambda 2600`$, Mgi $`\lambda 2853`$, Alii $`\lambda 1670`$, Cii $`\lambda 1334`$, Siii $`\lambda 1260`$, and Siiii $`\lambda 1206`$. If the flag transition for the species is not identifiable using the above criteria, but a weaker transition is, the equivalent width of the weaker transition is quoted (in Table Low and High Ionization Absorption Properties of Mgii Absorption–Selected Galaxies at Intermediate Redshifts. I. General Properties<sup>1</sup><sup>1</sup>affiliation: Based in part on observations obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among Caltech, the University of California, and NASA. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. <sup>2</sup><sup>2</sup>affiliation: Based in part on observations obtained with the NASA/ESA Hubble Space Telescope, which is operated by the STScI for the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5–26555.). We list these for completeness, but do not include them in our analysis.
In spectral regions with no detected absorption lines, the $`3\sigma `$ equivalent width limits were computed from Equation 2, assuming unresolved features. For any given undetected species, the transition with the most stringent limit was adopted. If there was a blend at the expected position of the transition, then an “upper limit” was obtained by quoting the equivalent width of the blended absorption line(s). Due to flat fielding uncertainties, the minimum limit quoted (observer frame) is $`0.13`$ Å (see Bahcall et al. 1996).
#### 3.1.5 Measuring Lyman Limit Breaks
Measurements of the Lyman limit breaks were made using KP techniques, as discussed in Schneider et al. (1993) and Jannuzi et al. (1998). In the cases of multiple or double breaks, we modeled the data employing the same technique applied by Churchill & Charlton (1999; see their Figure 3$`a`$). The redshifts of the Lyman limit break models were fixed at the redshifts of the Mgii absorbers. We varied the Hi column densities and superimposed synthetic spectra on the data. The quoted optical depths were obtained from the best matching model spectrum. Because the signal–to–noise ratio was often low, we did not use any fitting statistics, but simply attempted to make the depth of the break consistent with the model. In the cases of possible double or multiple breaks (due to redshift proximity), this modeling technique was successful at either singling out which Mgii system actually gives rise to the break, or placing constraints on the relative strengths of the double break. When a Lyman break is present, we quote a rough estimate of its optical depth in § 4.1, where each absorber is discussed individually. We also place a “$`+`$” in column six of Table Low and High Ionization Absorption Properties of Mgii Absorption–Selected Galaxies at Intermediate Redshifts. I. General Properties<sup>1</sup><sup>1</sup>affiliation: Based in part on observations obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among Caltech, the University of California, and NASA. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. <sup>2</sup><sup>2</sup>affiliation: Based in part on observations obtained with the NASA/ESA Hubble Space Telescope, which is operated by the STScI for the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5–26555. (described in § 4). If a break was not present, we tabulated a “$``$” and assign an upper limit on $`\mathrm{log}N(\text{H}\text{i})`$ of $`16.8`$ cm<sup>-2</sup>. When the location of the break was not covered in the spectra, we placed a “$`\mathrm{}`$”.
## 4. Presentation of System Properties
Line identifications of absorption features in FOS spectra is often plagued by line blending and other sources of confusion. Thus, we believe that a fairly complete description of the issues encountered with each system is warranted, especially in view of the possibility that any single system may be the subject of future detailed study as higher quality data become available.
The results from the HIRES/Keck spectra are presented in Table Low and High Ionization Absorption Properties of Mgii Absorption–Selected Galaxies at Intermediate Redshifts. I. General Properties<sup>1</sup><sup>1</sup>affiliation: Based in part on observations obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among Caltech, the University of California, and NASA. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. <sup>2</sup><sup>2</sup>affiliation: Based in part on observations obtained with the NASA/ESA Hubble Space Telescope, which is operated by the STScI for the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5–26555.. Tabulated are the quasar name, the absorber redshift, the Mgii $`\lambda 2796`$ rest–frame equivalent width, $`W_r(\text{Mg}\text{ii})`$, the Feii $`\lambda 2600`$ rest–frame equivalent width, $`W_r(\text{Fe}\text{ii})`$, and the Mgii $`\lambda 2853`$ rest–frame equivalent width, $`W_r(\text{Mg}\text{i})`$. Equivalent width limits are quoted at the $`3\sigma `$ level. In a few systems, we have detected Mnii, Caii, or Tiii; their detection is mostly an arbitrary function of wavelength coverage. We have not listed them in Table Low and High Ionization Absorption Properties of Mgii Absorption–Selected Galaxies at Intermediate Redshifts. I. General Properties<sup>1</sup><sup>1</sup>affiliation: Based in part on observations obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among Caltech, the University of California, and NASA. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. <sup>2</sup><sup>2</sup>affiliation: Based in part on observations obtained with the NASA/ESA Hubble Space Telescope, which is operated by the STScI for the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5–26555. (however, we show these data in Figure 5).
In Table Low and High Ionization Absorption Properties of Mgii Absorption–Selected Galaxies at Intermediate Redshifts. I. General Properties<sup>1</sup><sup>1</sup>affiliation: Based in part on observations obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among Caltech, the University of California, and NASA. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. <sup>2</sup><sup>2</sup>affiliation: Based in part on observations obtained with the NASA/ESA Hubble Space Telescope, which is operated by the STScI for the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5–26555., we present the results from the FOS/HST spectra. Tabulated are the quasar name, the absorber redshift, the Ly$`\alpha `$, Ly$`\beta `$, and Ly$`\gamma `$ equivalent widths, status of the Lyman break, and the equivalent widths of Alii, Aliii, Siii, Siiii, Siiv, Cii, Ciii, Civ, Nv, and Ovi. All equivalent widths are rest frame. The spectroscopic data are presented in Figures 5$`a`$$`ss`$. For each system, three sets of panels are shown: the HIRES detections, the FOS detections, and the FOS limits. The spectra are all normalized by the continuum fits. Both pre–COSTAR and COSTAR spectra are represented. We note the pre–COSTAR spectra in the discussion of individual systems (§ 4.1).
The panels with HIRES profiles show a velocity window of $`250`$ to $`250`$ km s<sup>-1</sup> centered on the Mgii $`\lambda 2796`$ optical depth mean. Ticks above the continuum give the velocity positions of the Voigt profiles components (the Voigt profile decomposition is described in Paper II). The panels with singlet FOS absorption lines show a velocity window of $`1250`$ to $`1250`$ km s<sup>-1</sup> centered on the line. The panels with doublet FOS absorption lines also show a velocity window of $`1250`$ to $`1250`$ km s<sup>-1</sup>, but the zero point is arbitrarily set half way between the doublet members in order to center the doublet in the panel. Ticks above the FOS spectra show the expected positions of the absorption based upon the Mgii Voigt profile components.
For each system we show the flag transition or both members of a doublet for the species presented in Tables Low and High Ionization Absorption Properties of Mgii Absorption–Selected Galaxies at Intermediate Redshifts. I. General Properties<sup>1</sup><sup>1</sup>affiliation: Based in part on observations obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among Caltech, the University of California, and NASA. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. <sup>2</sup><sup>2</sup>affiliation: Based in part on observations obtained with the NASA/ESA Hubble Space Telescope, which is operated by the STScI for the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5–26555. and Low and High Ionization Absorption Properties of Mgii Absorption–Selected Galaxies at Intermediate Redshifts. I. General Properties<sup>1</sup><sup>1</sup>affiliation: Based in part on observations obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among Caltech, the University of California, and NASA. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. <sup>2</sup><sup>2</sup>affiliation: Based in part on observations obtained with the NASA/ESA Hubble Space Telescope, which is operated by the STScI for the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5–26555.. If a weaker transition was used for a measurement, then it is labeled in Figures 5$`a`$$`ss`$ and a footnote is placed in Table Low and High Ionization Absorption Properties of Mgii Absorption–Selected Galaxies at Intermediate Redshifts. I. General Properties<sup>1</sup><sup>1</sup>affiliation: Based in part on observations obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among Caltech, the University of California, and NASA. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. <sup>2</sup><sup>2</sup>affiliation: Based in part on observations obtained with the NASA/ESA Hubble Space Telescope, which is operated by the STScI for the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5–26555. or Table Low and High Ionization Absorption Properties of Mgii Absorption–Selected Galaxies at Intermediate Redshifts. I. General Properties<sup>1</sup><sup>1</sup>affiliation: Based in part on observations obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among Caltech, the University of California, and NASA. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. <sup>2</sup><sup>2</sup>affiliation: Based in part on observations obtained with the NASA/ESA Hubble Space Telescope, which is operated by the STScI for the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5–26555.. For the HIRES spectra, we show the strongest Tiii, Caii, and/or Mnii, even though these species are not listed in Table Low and High Ionization Absorption Properties of Mgii Absorption–Selected Galaxies at Intermediate Redshifts. I. General Properties<sup>1</sup><sup>1</sup>affiliation: Based in part on observations obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among Caltech, the University of California, and NASA. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. <sup>2</sup><sup>2</sup>affiliation: Based in part on observations obtained with the NASA/ESA Hubble Space Telescope, which is operated by the STScI for the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5–26555..
### 4.1. Discussion of the Individual Systems
#### 4.1.1 $`0002+051`$, UM 18 ($`z_{\mathrm{em}}=1.90`$)
The FOS spectra of this quasar, which has four intervening Mgii absorbers, were previously studied by Jannuzi et al. (1998) and Koratkar et al. (1998). Jannuzi et al. obtained only the G270H grating spectrum, with pre–COSTAR optics, and using the $`0.25\mathrm{}\times 2\mathrm{}`$ slit. Koratkar et al. obtained both the G190H and G270H grating spectra in spectropolarimetry mode ($`1.0`$″ aperture). We have included only the Jannuzi et al. spectrum in our study, except for a constraint on the Lyman limit break obtained from the Koratkar et al. G190H spectrum.
$`z=0.5915`$ — This system has no detectable Feii or Mgi in the HIRES spectrum. In the FOS spectrum, Siii $`\lambda 1526`$ is blended with a Ly$`\gamma `$ line at $`z=1.5000`$ and we thus conservatively fit the entire blend and quote an upper limit for $`W_r(\text{Si}\text{ii})`$. Civ $`\lambda 1548`$ is blended with Siii $`\lambda 1260`$ at $`z=0.9560`$ and with Ly$`\gamma `$ at $`z=1.5359`$. The quoted upper limit for $`W_r(\text{C}\text{iv})`$ is taken to be the $`3\sigma `$ equivalent width limit at the expected position of $`\lambda 1550`$. Alii $`\lambda 1670`$ is blended with Ly$`\delta `$ at $`z=1.8016`$ so we measure that line and record it as the limit on $`W_r(\text{Al}\text{ii})`$. The Lyman limit was not covered.
$`z=0.8514`$ — This system is rich in Feii and Mgi in the HIRES spectrum. In the FOS spectrum, the quoted $`W_r(\text{Ly}\alpha )`$ may be very slightly overestimated due to the coincidence of Siiii $`\lambda 1206`$ at $`z=0.8665`$. We caution that the Siii $`\lambda 1526`$ transition is weak relative to our claimed $`\lambda 1260`$ detection, but is consistent within the permitted scaling. Jannuzi et al. (1998) identified a Ly$`\alpha `$ line at $`z=1.3592`$ at the location of Civ $`\lambda 1548`$, but our new knowledge of the presence of a Mgii system implies our preferred Civ $`\lambda 1548`$ identification. The unphysical doublet ratio of Civ $`\lambda \lambda 1548,1550`$ is due to a blend of $`\lambda 1550`$ with Ly$`\beta `$ at $`z=1.8016`$, which is part of a clear Lyman series. Both members of the Siiv $`\lambda \lambda 1393,1402`$ doublet are blended, the former with Ly$`\beta `$ at $`z=1.5160`$ and the latter with Galactic Feii $`\lambda 2586`$. A limit on $`W_r(\text{Si}\text{iv})`$ was recorded based on measuring the blend at the expected position of Siiv $`\lambda 1393`$. Nv $`\lambda 1242`$ is blended with Ly$`\beta `$ at $`z=1.2472`$. There is a strong Lyman limit break with $`\tau 1.4`$.
$`z=0.8665`$ — This system has no detectable Feii or Mgi in the HIRES spectrum. In the FOS spectrum, the blue wing of the Ly$`\alpha `$ line is blended with Ly$`\delta `$ at $`z=1.3866`$; we deblended the lines with a Gaussian fit to obtain the quoted $`W_r(\text{Ly}\alpha )`$. Jannuzi et al. (1998) favored the identification of a Ly$`\alpha `$ line at $`z=1.5653`$ at the location of Alii $`\lambda 1670`$. Furthermore, the profile is extremely broad (greater than 500 km s<sup>-1</sup> in the rest–frame), suggesting a blend; we thus quote an upper limit on $`W_r(\text{Al}\text{ii})`$. Siiii $`\lambda 1206`$ is coincident with the strong Ly$`\alpha `$ line at $`z=0.8514`$. Siiv $`\lambda 1393`$ is blended with Galactic Feii $`\lambda 2600`$ so an upper limit on $`W_r(\text{Si}\text{iv})`$ was measured at the position of $`\lambda 1402`$. Civ $`\lambda 1548`$ is blended with Ly$`\beta `$ at $`z=1.8184`$; we used the $`3\sigma `$ limit on $`\lambda 1550`$ to compute the upper limit on $`W_r(\text{C}\text{iv})`$. From a model of the Lyman limit, we determined that this system does not contribute to the observed break from the $`z=0.8514`$ system.
$`z=0.9560`$ — This system has no detectable Feii or Mgi in the HIRES spectrum. In the FOS spectrum, we claim a detection of Siiii $`\lambda 1206`$, but note that it could be blended with Ly$`\beta `$ at $`z=1.3011`$ (if Ly$`\alpha `$ is hidden in a complex blend at $`2797`$ Å). Siii $`\lambda 1260`$ is blended with Ly$`\gamma `$ at $`z=1.5359`$, and other members of that series are cleanly detected. The other Siii transitions are blended with other possible lines so cannot be used to corroborate or refute a Siii $`\lambda 1260`$ identification. We quote a limit on Siii from measurement of the blended line at the position of $`\lambda 1260`$. Siiv $`\lambda 1393`$ is blended in the red wing of Ly$`\gamma `$ at $`z=1.8016`$, but a clean upper limit can be derived from the position of Siiv $`\lambda 1402`$. The region where the Lyman limit is expected exhibits a complex shape, however, there is no apparent break.
#### 4.1.2 $`0058+019`$, PHL 938 ($`z_{\mathrm{em}}=1.96`$)
Only a fairly low signal–to–noise ratio, G190H FOS spectrum was obtained by Rao & Turnshek (1999) using the $`1.0`$″ aperture. The signal–to–noise ratio is significantly reduced blueward of the Lyman limit break at $`2250`$ Å.
$`z=0.6127`$ — This system is a damped Ly$`\alpha `$ absorber and is seen to have Feii, Mgi , Mnii, and Tiii absorption in the HIRES spectrum. In the FOS spectrum, the Siiv $`\lambda 1393`$ transition is clearly detected, though the $`\lambda 1402`$ could have a contribution from Ly$`11`$ from the strong Lyman series at $`z=1.4638`$. The Lyman limit was not covered.
$`z=0.7252`$ — This system has Mgi, but no Feii, in the HIRES spectrum. In the FOS spectrum, the Lyman series lines are somewhat ambiguous \[the data are quite noisy such that $`W_r(\text{Ly}\beta )`$ is greater than $`W_r(\text{Ly}\alpha )`$\]. The Ovi $`\lambda 1031`$ transition is in the red wing of a strong line, possibly Ly$`\beta `$ at $`z=0.7337`$; we determine an upper limit on $`W_r(\text{O}\text{vi})`$ at the position of $`\lambda 1037`$. The Lyman limit was not covered.
#### 4.1.3 PG $`0117+213`$ ($`z_{\mathrm{em}}=1.50`$)
A G270H FOS spectrum of this quasar was previously studied by Jannuzi et al. (1998) using the $`0.3`$″ aperture. Koratkar et al. (1998) obtained spectra in the spectropolarimetry mode using both the G190H and G270H gratings ($`1.0`$″ aperture). We have limited our study to the G270H spectrum from Jannuzi et al. , with the exception of analysis of Lyman limit breaks and of the damped Ly$`\alpha `$ line in the G190H spectrum. There are five intervening Mgii absorption systems.
$`z=0.5764`$ — This system is a damped Ly$`\alpha `$ absorber (Rao & Turnshek (1995)). In the HIRES spectrum, Mgi, Caii $`\lambda 3969`$ and several Tiii transition were detected (Caii $`\lambda 3934`$ was not covered). Feii, expected to be very strong, was not covered in the HIRES spectrum. In the FOS spectrum, Feii $`\lambda 1608`$ was detected. We quoted $`W_r(\text{Si}\text{ii})`$ using $`\lambda 1526`$, since all other Siii transitions were available only in the spectropolarimetry G190H spectrum. Alii $`\lambda 1670`$ is strong and could be affected by a blend with a Ly$`\alpha `$ line; however we note that $`W_r(\text{Al}\text{ii})`$ is consistent with a black–bottom Alii profile with the same kinematic spread as Mgii. We point out that the strong damped Ly$`\alpha `$ line was measured in a G190H spectrum obtained in spectropolarimetry mode. The Lyman limit was not covered.
$`z=0.7291`$ — This system has Feii and Mgi in the HIRES spectrum. We quoted $`W_r(\text{Si}\text{ii})`$ using $`\lambda 1526`$, since all other Siii transitions were available only in the spectropolarimetry G190H spectrum. We identify Cii $`\lambda 1334`$ in this system. The line was previously identified as Ly$`\beta `$ at $`z=1.2501`$ by Jannuzi et al. (1998), however, the ratio $`W_r(\text{Ly}\beta )/W_r(\text{Ly}\alpha )`$ would be large for the small $`W_r(\text{Ly}\alpha )`$ in that system and there are no other transitions to corroborate it. Ly$`\alpha `$ and the Lyman limit were only covered on the spectropolarimetry G190H spectrum, and could not be measured.
$`z=1.0480`$ — This system has Feii and very weak Mgi detected in the HIRES spectrum. In the FOS spectrum, we have quoted a detection for Siiii $`\lambda 1206`$, though we note Jannuzi et al. (1998) identify this absorption feature as Ly$`\alpha `$ at $`z=1.0326`$. We quote a Cii $`\lambda 1334`$ detection; though Jannuzi et al. have identified this feature as Feii $`\lambda 1144`$ at $`z=1.3868`$, their identification is not consistent with the absence of the Feii $`\lambda 2383`$ transition in the more sensitive HIRES spectrum. The quoted upper limit for $`W_r(\text{Si}\text{iv})`$ is obtained from $`\lambda 1402`$ because $`\lambda 1393`$ is blended with Galactic Mgi. $`W_r(\text{C}\text{iv})`$ is conservatively quoted as an upper limit because both the $`\lambda 1548`$ and $`\lambda 1550`$ are detected just below the $`3\sigma `$ level. There is a partial Lyman limit break with $`\tau 0.9`$ measured in the G190H spectropolarimetry spectrum.
$`z=1.3250`$ — This system has Feii, Mgi and Aliii $`\lambda 1863`$ in the HIRES spectrum. The Civ doublet is taken from the ground–based spectrum of SS92, which has resolution $`R860`$. In the FOS spectrum, Ly$`\alpha `$ is coincident with Siiii $`\lambda 1206`$ from the $`z=1.3430`$ Mgii system; the quoted $`W_r(\text{Ly}\alpha )`$ may be larger than the true Ly$`\alpha `$ absorption strength. Gaussian deblending was used to measure $`W_r(\text{Ly}\beta )`$ because the blue wing of the Ly$`\beta `$ line is blended with Galactic Feii $`\lambda 2383`$, but a clear asymmetry is apparent. $`W_r(\text{Si}\text{iv})`$ $`\lambda 1393`$ could have a contribution from a blend with Alii $`\lambda 1670`$ at $`z=0.9400`$, but this contribution should be small. The Siiv $`\lambda 1402`$ is not formally detected, but this is not inconsistent with $`\lambda 1393`$ due to an uncertain continuum fit near the spectrum edge; we claim a detection for Siiv. Siiii $`\lambda 1206`$ is blended with Galactic Mgii $`\lambda 2803`$ in its blue wing and possibly with an unidentified line in its red wing; the equivalent width of the full blend was taken as a limit on Siiii. The upper limit for $`W_r(\text{N}\text{v})`$ is obtained from the expected position of $`\lambda 1242`$ because $`\lambda 1238`$ is blended with Siiii $`\lambda 1206`$ at $`z=1.3868`$. Both Ovi $`\lambda 1031`$ and $`\lambda 1037`$ are blended with other lines, the former with Ly$`\beta `$ at $`z=1.3389`$, and the latter with Ovi $`\lambda 1031`$, also from $`z=1.3389`$. Absorption from Ovi might be present in this system, but it cannot be determined, so we quote an upper limit obtained by fitting the weaker blended feature at the position of $`\lambda 1037`$. The Lyman limit break shows structure suggestive of a “double break”. A model of the Lyman break revealed that both this system and the $`z=1.3430`$ system contribute roughly equally, ach having $`\tau 0.6`$.
$`z=1.3430`$ — This system has Feii and Aliii $`\lambda \lambda 1854,1862`$ in the HIRES spectrum. The Civ doublet is taken from the ground–based spectrum of SS92, which has resolution $`R860`$ at the position of the doublet. In the FOS spectrum, Nv $`\lambda 1238`$ is blended with Ly$`\alpha `$ at $`z=1.3868`$. We quote a detection of Nv based upon $`\lambda 1242`$, which is blueward of Ly$`\alpha `$ at $`z=1.3981`$. A significant feature was present at the location of Siiv $`\lambda 1393`$ redward of the Ly$`\alpha `$ forest, but $`\lambda 1402`$ was not covered. We tentatively claim a detection for the Siiv doublet (consistent with our policy on line identifications). Siiii $`\lambda 1206`$ is coincident with the strong Ly$`\alpha `$ line associated with the $`z=1.3250`$ Mgii absorber; we quote the equivalent width of this line as an upper limit on $`W_r(\text{Si}\text{iii})`$. There is a Lyman limit break with $`\tau 1.3`$ measured in the G190H spectropolarimetry spectrum. The Lyman limit break shows structure suggestive of a “double break”. Modeling showed that both this system and the $`z=1.3250`$ system contribute roughly equally to the break with each having $`\tau 0.6`$.
#### 4.1.4 PKS $`0454220`$ ($`z_{\mathrm{em}}=0.53`$)
G130H, G190H, and G270H FOS spectra of this quasar were obtained by M. Burbidge (see Cohen et al. (1991). Of the pre–COSTAR spectra included in our study, these are the only ones obtained with a “large” aperture ($`1`$″). As such, the quoted equivalent widths may be systematically small; we did not apply any correction factors.
$`z=0.4744`$ — This is a near–DLA system (also see Rao et al. 1995); the measured $`W_r(\text{Ly}\alpha )=5.5`$ Å places it on the logarithmic part of the curve of growth with $`N(\text{H}\text{i})10^{19.5}`$ cm<sup>-2</sup>. Due to a blend with Siiii $`\lambda 1206`$ at $`z=0.4833`$ in the red wing of the Ly$`\alpha `$ line, $`W_r(\text{Ly}\alpha )`$ could be slightly overestimated. Detected in the HIRES spectrum are several Feii transitions, Mgi , and Caii $`\lambda 3969`$ ($`\lambda 3934`$ was not covered). In the FOS spectrum, the Civ doublet is present on both the G190H and G270H gratings; the equivalent widths of the $`\lambda 1548`$ transition measured from the two spectra are consistent ($`0.93\pm 0.04`$ Å and $`1.09\pm 0.11`$ Å, respectively.) The red wing of Ciii is blended with Ly$`\gamma `$ at $`z=0.4833`$ and an unidentified line, thus a limit is quoted based on a fit to the blend. Ovi $`\lambda 1031`$ is blended with Ly$`\beta `$ at $`z=0.4833`$ and $`\lambda 1037`$ is coincident with both the $`\lambda 1031`$ member of Ovi at $`z=0.4833`$ and a dominating unidentified line. We conservatively quote an upper limit on $`W_r(\text{O}\text{vi})`$ based on the equivalent width of the blend at the expected position of Ovi $`\lambda 1031`$. Accurate measurement of the $`N(\text{H}\text{i})`$ optical depth from the Lyman limit is compromised by geocoronal emission; given that this system is a near–DLA, the flux below $`\lambda 1345`$ Å should be zero. Modeling showed that the break itself arises from the $`z=0.4833`$ system (see below).
$`z=0.4833`$ — In the HIRES spectrum, this system has detected Feii, Mgi, and Caii $`\lambda 3934`$ (Caii $`\lambda 3969`$ was not covered). In the FOS spectrum, Civ was covered by both the G190H and G270H gratings, but the profile shapes and equivalent widths are not consistent (the G190H appears compromised, but the G270H $`\lambda 1548`$ transition is large relative to $`\lambda 1550`$); we have quoted the average of the equivalent widths from the two gratings and present the G270H spectrum in Figure 5. Any Siiii $`\lambda 1206`$ that is present in the spectrum is blended in the red wing of the Ly$`\alpha `$ line at $`z=0.4744`$; we have quoted $`W_r(\text{Ly}\alpha )`$ for the upper limit on $`W_r(\text{Si}\text{iii})`$. The upper limit on $`W_r(\text{Si}\text{iv})`$ was obtained from $`\lambda 1402`$, since $`\lambda 1393`$ is blended with Siiv $`\lambda 1402`$ at $`z=0.4744`$. Ovi $`\lambda 1031`$ is blended with Ovi $`\lambda 1037`$ at $`z=0.4744`$ and a dominating unidentified line. However, a clean upper limit on $`W_r(\text{O}\text{vi})`$ was obtained from $`\lambda 1037`$. Modeling showed that the Lyman limit break arises from this system, however accurate measurement of the $`N(\text{H}\text{i})`$ optical depth is compromised by geocoronal emission.
#### 4.1.5 PKS $`0454+039`$ ($`z_{\mathrm{em}}=1.35`$)
G190H and G270H FOS spectra of this quasar (obtained by Bergeron with the $`1.0`$″ square aperture) were studied by Boissé et al. (1998), and a separate G270H spectrum ($`1.0`$″ round aperture) was obtained by Bowen. We independently reduced these data, and we present results below from the Bergeron G190H and from a combined G270H spectrum from the Bowen and Bergeron observations. Some slight differences between our analysis results and those of Boissé et al. are discussed on a system by system basis below.
$`z=0.6428`$ — This system has Feii in the HIRES spectrum. In the FOS spectrum, Ly$`\alpha `$ is blended with Ly$`7`$ at $`z=1.1532`$ and a Ly$`\alpha `$ line at $`z=0.6448`$, but it is distinct such that accurate deblending is possible. We claim a tentative detection of Siiv, which appears to be present in both the G190H and G270H spectra. We note however, that the blue wing of $`\lambda 1393`$ is blended with Ly$`\alpha `$ at $`z=0.8801`$ and that there may be a small contribution to $`\lambda 1402`$ from Nv $`\lambda 1242`$ at $`z=0.8596`$, though the latter is not formally detected. Though Churchill & Le Brun (1998) claimed a tentative upper limit on Civ, our reduction and continuum fit of the G190H spectrum yielded a detection above the $`3\sigma `$ level; we thus claim a detection for Civ. Siiii $`\lambda 1206`$ is blended with Ly$`9`$ at $`z=1.1532`$; we fit this feature to obtain an upper limit. The Lyman limit was not covered and Ly$`\beta `$ is below the Lyman break from the $`z=0.8596`$ system where the flux is zero.
$`z=0.8596`$ — This is a well studied damped Ly$`\alpha `$ system (e.g. Steidel et al. 1995; Lu et al. 1996b; Boissè et al. 1998). In the HIRES spectrum, a suite of five Feii transitions are detected, as are the Mnii $`\lambda 2600`$ triplet, Mgi and Tiii $`\lambda 3074`$ and $`\lambda 3243`$. In the FOS spectrum, numerous transitions are detected, and the optical depth of the Lyman limit is $`\tau 3.3`$.
$`z=0.9315`$ — This system has Feii in the HIRES spectrum. Our measurement of the very weak Ly$`\alpha `$ line was consistent with that reported by Churchill & Le Brun (1998). Ly$`\beta `$ is blended with the Lyman series at $`z=1.1532`$ and Ly$`\gamma `$ is blended with Ly$`\delta `$ at $`z=0.9780`$. Siiii $`\lambda 1206`$ is in the blue wing of a Ly$`\alpha `$ line at $`z=0.9185`$ so a non–restrictive limit was obtained by fitting this blend. Though it appears that a clean Civ doublet is present, $`\lambda 1548`$ is actually Feii $`\lambda 1608`$ from the $`z=0.8596`$ DLA and $`\lambda 1550`$ is Mgii $`\lambda 2796`$ from a dwarf galaxy at $`z=0.0714`$. A non–restrictive limit on $`W_r(\text{C}\text{iv})`$ is obtained by fitting the feature coincident with the expected position of $`\lambda 1548`$. Ovi $`\lambda 1031`$ is blended with the blue wing of Ly$`8`$ at $`z=1.1532`$ and $`\lambda 1037`$ is blended with Ly$`6`$, also at $`z=1.1532`$; we used $`W_r(\text{Ly}8)`$ to obtain an upper limit on $`W_r(\text{O}\text{vi})`$. There is no Lyman limit break.
$`z=1.1532`$ — In the HIRES spectrum, Feii, Mgi and the Aliii $`\lambda \lambda 1854,1862`$ are present. The Civ doublet is taken from the ground–based spectrum of SBS, which has resolution $`R2200`$ at the observed wavelength of Civ. In the FOS spectrum, the quoted $`W_r(\text{Si}\text{iv})`$ is taken from $`\lambda 1402`$ because $`\lambda 1393`$ is blended with Mgii $`\lambda 2796`$ from the dwarf galaxy at $`z=0.0714`$. Boissè et al. (1998) reported detection of Nv $`\lambda \lambda 1238,1242`$, but the feature located where $`\lambda 1238`$ is expected may be a Ly$`\alpha `$ line at $`z=1.1941`$, since $`\lambda 1242`$ would be in the extreme blue wing of the possible Ly$`\alpha `$ line at $`z=1.2025`$. We quote an upper limit on $`W_r(\text{N}\text{v})`$ from measuring the blend at the position of $`\lambda 1238`$ . Though Ovi may be present, $`\lambda 1031`$ is blended with Ly$`\beta `$ at $`z=1.1672`$ and Ovi is in a complex region in the wing of the damped Ly$`\alpha `$ line at $`z=0.8596`$. The optical depth of the Lyman limit is $`\tau 1.1`$.
#### 4.1.6 PKS $`0823223`$ ($`z_{\mathrm{em}}>0.92`$)
This quasar, a luminous BL Lacertae object, has only a lower limit on its emission redshift. Only a G270H FOS spectrum was obtained by Rao & Turnshek (1999) using the $`1.0`$″ aperture.
$`z=0.7055`$ — This system has no detectable Feii or Mgi in the HIRES spectrum. Ly$`\alpha `$ was not covered in the FOS spectrum. The upper limit on $`W_r(\text{Si}\text{ii})`$ is obtained from the $`\lambda 1526`$ transition because the other Siii transitions are not covered. Cii $`\lambda 1334`$, if present, is blended with strong Siii $`\lambda 1190`$ at $`z=0.9110`$, so this blend is fit to obtain an upper limit on $`W_r(\text{C}\text{iv})`$. For the upper limit on $`W_r(\text{Si}\text{iv})`$, we use $`\lambda 1402`$ because $`\lambda 1393`$ is blended with Galactic Feii $`\lambda 2374`$ and $`\lambda 2383`$. The Lyman limit was not covered.
$`z=0.9110`$ — This system was reported to be a “double” system by Veron–Cetty et al. (1990). Several Feii transitions and Mgi are present in the HIRES spectrum. The FOS spectrum is very rich in many, cleanly identified, transitions. The $`\lambda 1242`$ member of the Nv doublet is blended with Galactic Feii $`\lambda 2374`$, but $`\lambda 1238`$ is unambiguously detected. The Lyman break was not covered.
#### 4.1.7 $`0958+551`$, MRK 132 ($`z_{\mathrm{em}}=1.76`$)
The G270H FOS spectrum, obtained with the $`0.3`$″ round aperture, was studied by Jannuzi et al. (1998). The FOS spectrum covers only the Ly$`\alpha `$ forest down to a strong Lyman limit break at 2200 Å from a $`z=1.7327`$ system that wipes out all flux below this wavelength.
$`z=1.2113`$ — This system has no detectable Feii or Mgi in the HIRES spectrum. In the FOS spectrum, Ly$`\alpha `$ is blended with $`\lambda 1550`$ from a possible Civ doublet at $`z=0.7330`$. There are no corroborating transitions for the Civ redshift in the HIRES spectrum, but there is a possible Feii $`\lambda 1608`$ line in the FOS spectrum. We quote a non–restrictive upper limit on $`W_r(\text{Ly}\alpha )`$. Siiii $`\lambda 1206`$ is blended with Ciii $`\lambda 977`$ at $`z=1.7325`$, the redshift of the strong Lyman limit break. There is a strong absorption feature coincident with Siii $`\lambda 1260`$ (could be Feii $`\lambda 1608`$ at $`z=0.7330`$), but its strength is not consistent with detection limits based to the Siii $`\lambda \lambda 1190,1193`$; we have used $`\lambda 1193`$ to place an upper limit on $`W_r(\text{Si}\text{ii})`$. The Lyman limit break was not covered.
$`z=1.2724`$ — This system has Feii $`\lambda 2383`$ in the HIRES spectrum. The Civ doublet is taken from the ground–based spectrum of SBS, which has resolution $`R840`$ at this wavelength. In the FOS spectrum, only Ly$`\alpha `$ is clearly detected above the $`3\sigma `$ threshold. A possible detection of Siiv is not clear; there could be coincident Ly$`\alpha `$ lines blended with both members of the doublet. We conservatively quote an upper limit on $`W_r(\text{Si}\text{iv})`$. The Lyman limit break was not covered. Nv $`\lambda 1242`$ is possibly blended with Ovi $`\lambda 1031`$ at $`z=1.7327`$, so $`\lambda 1238`$ provides an upper limit on $`W_r(\text{N}\text{v})`$.
#### 4.1.8 PG $`1206+459`$ ($`z_{\mathrm{em}}=1.16`$)
Both the G190H and G270H FOS spectra of this quasar have been studied by Jannuzi et al. (1998). They were obtained with pre–COSTAR optics using the $`0.25\mathrm{}\times 2\mathrm{}`$ slit. The metal line properties of the Mgii systems were studied in detail by Churchill & Charlton (1999).
$`z=0.9276`$ — This system was reported as a double system by Jannuzi et al. (1998) and by Churchill & Charlton (1999), with redshifts $`z=0.9254`$ and $`z=0.9276`$. In the HIRES spectrum, Feii and Mgi are present. The FOS spectrum has the richest complement of transitions of all the systems in our sample. Cii $`\lambda 1334`$ is blended with a possible Ly$`\alpha `$ line at $`z=1.1166`$. Cii $`\lambda 1036`$ is blended with Ovi $`\lambda 1037`$ and with Ovi $`\lambda 1031`$ at $`z=0.9343`$. The quoted $`W_r(\text{C}\text{ii})`$ is based upon $`\lambda 1334`$ using a Gaussian deblending, which likely yielded a slightly large value. The Civ profile required a double Gaussian fit and the $`\lambda 1550`$ from $`z=0.9254`$ is blended with $`\lambda 1548`$ from $`z=0.9276`$. We approximately measured $`W_r(\text{C}\text{iv})`$ (estimated uncertainty of $`1`$ Å) as the sum of half of the equivalent width of this blend and the equivalent width of the unblended $`\lambda 1548`$ from the $`z=0.9254`$ system. Ovi also required a double Gaussian fit, but $`W_r(\text{O}\text{vi})`$ could be determined more accurately due to the larger separation of Ovi $`\lambda \lambda 1031,1037`$. Ciii $`\lambda 977`$ is clearly detected but is blended with Ly$`\alpha `$ at $`z=0.5482`$ and a smaller Ly$`\gamma `$ line from the $`z=0.9343`$ Mgii absorber; our upper limit is the equivalent width of the blend. Modeling of the Lyman limit showed that the break is entirely due to the $`z=0.9276`$ system with $`\tau 0.9`$ \[see Churchill & Charlton (1999)\].
$`z=0.9343`$ — This system has no detectable Feii or Mgi in the HIRES spectrum. $`W_r(\text{O}\text{vi})`$ was measured using $`\lambda 1037`$ because $`\lambda 1031`$ is blended with Ovi $`\lambda 1037`$ at $`z=0.9276`$. In the FOS spectrum, Ly$`\gamma `$ is blended with with Ly$`\alpha `$ at $`z=0.5482`$ and Ciii $`\lambda 977`$ at $`z=0.9276`$ such that only a non–restrictive limit can be derived from a fit to the blend. Cii $`\lambda 1334`$ is blended with strong Ly$`\alpha `$ at $`z=1.1223`$. This system makes no contribution to the Lyman limit break (see notes on the $`z=0.9276`$ system).
#### 4.1.9 PG $`1241+176`$ ($`z_{\mathrm{em}}=1.27`$)
The pre–COSTAR G270H FOS spectrum of this quasar, obtained with the $`0.25\mathrm{}\times 2\mathrm{}`$ slit, has been studied by Jannuzi et al. (1998).
$`z=0.5505`$ — In the HIRES spectrum, Feii, Mgi and Caii $`\lambda 3969`$ were detected ($`\lambda 3934`$ was not covered). In the FOS spectrum, detected Siii $`\lambda 1526`$ was used to measure $`W_r(\text{Si}\text{ii})`$ because no other Siii transitions are covered. The Lyman limit was not covered.
$`z=0.5584`$ — This system has no detectable Feii or Mgi in the HIRES spectrum. In the FOS spectrum, Siii $`\lambda 1526`$ was used to measure $`W_r(\text{Si}\text{ii})`$ because no other Siii transitions are covered. Civ $`\lambda 1550`$ was not detected at the $`3\sigma `$ level, however, it is consistent with the weak, but significant $`\lambda 1548`$ detection; we have quoted a detection for Civ. The Lyman limit was not covered.
$`z=0.8955`$ — This system has no detectable Feii or Mgi in the HIRES spectrum. In the FOS spectrum, only Ly$`\alpha `$ was unambiguously detected. Siiii $`\lambda 1206`$ is blended with Ovi $`\lambda 1031`$ at $`z=1.2154`$; it is difficult to evaluate the relative contributions so we quote an upper limit on $`W_r(\text{Si}\text{iii})`$. The Siiv doublet is somewhat ambiguous; if real, it is offset redward from the predicted location. The $`\lambda 1393`$ transition may be Ly$`\alpha `$ at $`z=1.1742`$ and the $`\lambda 1402`$ transition is blended with Cii $`\lambda 1334`$ at $`z=0.9927`$, a system with other corroborating lines. The limit on $`W_r(\text{Si}\text{iv})`$ was determined by fitting the blend at the expected position of $`\lambda 1393`$. The location of Nv $`\lambda 1242`$ is blended with an unknown transition in the blue wing of Ovi $`\lambda 1037`$ at $`z=1.2720`$, but a clean limit on $`W_r(\text{N}\text{v})`$ is determined at the position of $`\lambda 1393`$. (Cii $`\lambda 1036`$ at $`z=1.2720`$ is ruled out as the blend based upon the absence of stronger Cii $`\lambda 1334`$). The Lyman limit was not covered.
#### 4.1.10 PG $`1248+401`$ ($`z_{\mathrm{em}}=1.03`$)
Both the G190H and G270H FOS spectra ($`0.3`$″ aperture) of this quasar have been studied by Jannuzi et al. (1998).
$`z=0.7730`$ — This system has detected Feii and Mgi in the HIRES spectrum. The FOS spectrum is very rich with clearly detected transitions. As seen in Ly$`\alpha `$ and Ly$`\beta `$, there is a higher redshift system at $`z=0.7760`$ with no detectable Mgii in the HIRES spectrum. There are no detected metal lines from that system that might contaminate those of the $`z=0.7730`$ Mgii absorber. $`W_r(\text{Si}\text{ii})`$ was obtained by Gaussian deblending of $`\lambda 1260`$ with Siiii $`\lambda 1206`$ at $`z=0.8546`$. We have tentatively claimed Nv as a detection, noting that $`\lambda 1242`$ is not formally detected, but is consistent with the expected value from $`\lambda 1238`$. Nv was also detected by the systematic procedures of Jannuzi et al. (1998). $`W_r(\text{Ly}\gamma )`$ is quoted as an upper limit, since its measured value is inconsistent with $`W_r(\text{Ly}\alpha )`$ and $`W_r(\text{Ly}\beta )`$. The Lyman limit is on the edge of the G190H spectrum where the signal–to–noise ratio is hopeless.
$`z=0.8546`$ — This system has several Feii transitions but no Mgi in the HIRES spectrum. In the FOS spectrum, this system is part of a group of $`\text{Ly}\alpha +\text{Ly}\beta `$ absorbers at redshifts $`0.8524`$, $`0.8585`$, and $`0.8614`$. This is one of the examples in the KP data set of extensive metal–line systems occurring in overdensities in the distribution of Ly$`\alpha `$ absorbers (Bahcall et al. (1996); Jannuzi (1998)). The blend in the blue wing of Siiii $`\lambda 1206`$ is Siii $`\lambda 1260`$ at $`z=0.7730`$. $`W_r(\text{Si}\text{iii})`$ was obtained by Gaussian deblending. We quote an upper limit on $`W_r(\text{Si}\text{iv})`$ because $`\lambda 1393`$ is blended with Galactic Feii $`\lambda 2586`$ and $`\lambda 1402`$ is blended with Galactic Feii $`\lambda 2600`$. There is no Lyman limit break in a very clean portion of the spectrum.
#### 4.1.11 $`1317+277`$, TON 153 ($`z_{\mathrm{em}}=1.02`$)
Both the G190H and G270H FOS spectra of this quasar have been studied by Bahcall et al. (1996). These spectrum were obtained with pre–COSTAR optics using the $`0.25\mathrm{}\times 2\mathrm{}`$ slit. The G160L FOS spectrum, covering the Lyman limit, was presented by Bahcall et al. (1993).
$`z=0.6601`$ — In the HIRES spectrum, Feii, Mgi are present. In the FOS spectrum, there is an unknown blend with Nv $`\lambda 1242`$ but a clean upper limit on $`W_r(\text{N}\text{v})`$ is taken from the position of $`\lambda 1238`$. Ovi $`\lambda 1031`$ is blended with Ly$`\beta `$ at $`z=0.6691`$; we used $`\lambda 1037`$ to place an upper limit on $`W_r(\text{O}\text{vi})`$. There is a strong Lyman limit break, with $`\tau 5.4`$, associated with this system.
#### 4.1.12 PG $`1329+412`$ ($`z_{\mathrm{em}}=1.94`$)
The G190H grating spectrum ($`1.0`$″ aperture) was obtained by K. Lanzetta. The G270H grating was obtained by Rao & Turnshek (1999), also with the $`1.0`$ aperture. This is a complex spectrum with at least 10 known metal–line absorption redshifts, several of which give rise to a Lyman series. In the G190H spectrum, there is a strong Lyman limit break at $`2083`$ Å from a $`z=1.2841`$ system that wipes out all flux below this wavelength. The two strongest Lyman series systems at $`z=1.8408`$ and $`z=1.8367`$, give rise to a partial Lyman limit break. The former is associated with a Civ absorber (SBS) and has strong Ovi $`\lambda \lambda 1031,1037`$ absorption in the FOS spectrum.
$`z=0.5008`$ — This system had no detectable Feii nor Mgi in the HIRES spectrum. In the FOS spectrum, the upper limit on $`W_r(\text{Si}\text{ii})`$ was obtained using the $`\lambda 1526`$ transition since it was the only Siii line redward of the strong Lyman limit break at $`z=1.2841`$. The Siiv $`\lambda \lambda 1393,1402`$ is sitting on the shoulder of this break, and is thus fairly noisy, with $`\lambda 1402`$ blended with a possible Ly$`\alpha `$ line at $`z=0.7308`$; a non–restrictive upper limit is quoted from the noisy region at the expected position of $`\lambda 1393`$. The Lyman limit was not covered.
$`z=0.8933`$ — The HIRES spectrum is relatively noisy so that some weak kinematic outlying clouds could be missed. Near Mgii $`\lambda 2796`$ there is an unidentified line at $`v=+80`$ km s<sup>-1</sup> (which cannot be Mgii because the $`\lambda 2803`$ transition is absent). Feii was detected, but Mgi was not present. In the FOS spectrum, only Ly$`\alpha `$ is detected. Alii $`\lambda 1670`$ is blended in the red wing of strong Ly$`\alpha `$ at $`z=1.6008`$. Cii $`\lambda 1334`$ is blended with Ly$`\delta `$, also at $`z=1.6008`$. The Siiv doublet is on the shoulder of the “double”, partial Lyman limit break systems; $`\lambda 1393`$ is blended with Ly$`6`$ at $`z=1.8336`$ and $`\lambda 1402`$ is blended with the blue wing of a complex blend, including two Ly$`ϵ`$ lines at $`z=1.8366`$ and $`1.8408`$, and Ly$`\beta `$ at $`z=1.6008`$. We quote an upper limit for $`W_r(\text{Si}\text{iv})`$. The upper limit on $`W_r(\text{C}\text{iv})`$ is obtained using $`\lambda 1550`$ because $`\lambda 1548`$ is blended with Ovi $`\lambda 1031`$ at $`z=1.8408`$. The Lyman limit is wiped out by the break at $`2083`$ Å.
$`z=0.9739`$ — This system has no detectable Feii transitions nor Mgi in the HIRES spectrum. In the FOS spectrum, the Ly$`\alpha `$ line is the central member of a triple blend with Ly$`\gamma `$ at $`z=1.4716`$ being the red–most member and an unidentified line (possibly Ly$`\alpha `$ at $`z=0.9729`$) to the blue. The quoted $`W_r(\text{Ly}\alpha )`$ was obtained by Gaussian deblending, but because the profile is complicated the error is likely underestimated. $`W_r(\text{Si}\text{ii})`$ was also obtained using Gaussian deblending, due to an unidentified line, possibly Ly$`\alpha `$ at $`z=1.0499`$. The Cii $`\lambda 1334`$ transition is detected in a crowded region of the spectrum, but is unblended between Ly$`6`$ at $`z=1.8366`$ and Ly$`7`$ at $`z=1.8408`$. Siiv $`\lambda 1393`$ is clearly detected, not being blended with any Lyman series lines, however, $`\lambda 1402`$ is blended with Ly$`\alpha `$ at $`z=1.2830`$, which is corroborated by a Ly$`\beta `$ line. For Civ, $`\lambda 1550`$ is stronger at the $`4\sigma `$ level, suggesting a blend. The upper limit on $`W_r(\text{N}\text{v})`$ is obtained from $`\lambda 1242`$ because $`\lambda 1238`$ is blended with Niii $`\lambda 989`$ at $`z=1.4714`$. The Lyman limit is wiped out by the break at $`2083`$ Å.
$`z=0.9984`$ — This system has several Feii transitions, but Mgi was not covered in the HIRES spectrum. In the FOS spectrum, the Ly$`\alpha `$ is blended with an unidentified line; we measured $`W_r(\text{Ly}\alpha )`$ using Gaussian deblending, but note that the value is uncertain. Cii is blended with a strong Ly$`\alpha `$ line at $`z=1.6008`$. Siiv $`\lambda 1402`$ is blended with Galactic Mgii $`\lambda 2803`$. The Civ doublet, which brackets a possible Ly$`\alpha `$ line at $`z=1.5477`$, is undetected. There is also a possible Ly$`\alpha `$ line at $`z=1.0371`$ slightly to the blue of Nv $`\lambda 1238`$. The Lyman limit is wiped out by the break at $`2083`$ Å.
#### 4.1.13 PKS $`1354+195`$ ($`z_{\mathrm{em}}=0.72`$)
The G160L, G190H, and G270H FOS spectra of this quasar have been studied by Jannuzi et al. (1998) and Bergeron et al. (1994). These spectra were obtained pre–COSTAR, using the $`0.25\mathrm{}\times 2\mathrm{}`$ slit.
$`z=0.4566`$ — The HIRES spectrum is fairly noisy in the region of the Mgii doublet and Feii transitions, so it is possible that some kinematic outlying clouds were missed. Feii $`\lambda 2600`$ and Mgi were detected. In the FOS spectrum, the detections are all unambiguous.and the upper limits are not compromised by blends. We note that there is a Lyman limit break with $`\tau 0.9`$ in the G160L spectrum (Jannuzi et al. 1998).
$`z=0.5215`$ — This system has no detected Feii transitions nor Mgi in the HIRES spectrum. Siiii $`\lambda 1206`$ is blended with Siii $`\lambda 1260`$ at $`z=0.4566`$, that latter being corroborated by the presence of a proper–strength Siii $`\lambda 1526`$ line. The line at the expected wavelength of Cii $`\lambda 1334`$ is Siiv $`\lambda 1393`$ at $`z=0.4566`$. Based upon the G160L spectrum, there is no evidence for a Lyman limit break at this redshift.
#### 4.1.14 $`1622+238`$, 3C 336 ($`z_{\mathrm{em}}=0.93`$)
Both the G190H and G270H FOS spectra ($`1.0`$″ aperture) of this quasar have been studied by Steidel et al. (1997). The background level for the G190H is not certain. The spectrum has slightly negative flux in the core of the damped Ly$`\alpha `$ line at $`z=0.6561`$ and below the Lyman limit break. Thus, the equivalent widths of lines measured in this grating may be systematically small. The HIRES spectrum is quite noisy, such that the detection threshold is low for accompanying transitions.
$`z=0.4720`$ — In the HIRES spectrum, only the Mgii doublet is detected. In the FOS spectrum, the Civ $`\lambda 1550`$ is blended with Siiii $`\lambda 1206`$ at $`z=0.8913`$, but the $`\lambda 1548`$ appears to be clear. The upper limit on $`W_r(\text{Si}\text{ii})`$ was obtained from $`\lambda 1526`$ on the G270H grating because the stronger $`\lambda 1260`$ transition was in a noisy region of the G190H spectrum. Cii is ambiguous, being in a very noisy part of the spectrum and possibly being blended with Ovi $`\lambda 1037`$ at $`z=0.8913`$; we quote the $`3\sigma `$ upper limit on $`W_r(\text{C}\text{ii})`$. The upper limit on $`W_r(\text{Si}\text{iv})`$ is obtained using $`\lambda 1402`$ because $`\lambda 1393`$ is blended with a strong line, which could be Ly$`\alpha `$ at $`z=0.6868`$ with a contribution from Nv $`\lambda 1238`$ at $`z=0.6561`$. The Lyman limit break was not covered.
$`z=0.6561`$ — This system is a damped Ly$`\alpha `$ absorber (Steidel et al. (1997)). In the HIRES spectrum, Feii is strong and Tiii $`\lambda 3385`$ is detected. In the FOS spectrum, the upper limit on $`W_r(\text{N}\text{v})`$ is obtained using $`\lambda 1242`$ because $`\lambda 1238`$ is blended with a strong line, which could be Ly$`\alpha `$ at $`z=0.6868`$ with a contribution from Siiv $`\lambda 1393`$ at $`z=0.4720`$. The Lyman limit break was not covered.
$`z=0.7971`$ — This system was first seen in Civ absorption. Only the Mgii $`\lambda 2803`$ transition was detected because $`\lambda 2796`$ was wiped out by the pen mark on the HIRES CCD. Neither Feii nor Mgi was detected. Siiv $`\lambda 1402`$ is blended with Cii $`\lambda 1334`$, but $`\lambda 1393`$ is cleanly measured. A robust measurement of $`W_r(\text{Ly}\beta )`$ was not possible due to the signal to noise of the spectrum. Also, Ly$`\beta `$ may suffer from blending so we have quoted an upper limit on $`W_r(\text{Ly}\beta )`$. Though Siiii $`\lambda 1206`$ is in a busy part of the FOS spectrum, there are no other candidates for the line; the line immediately to the blue is Feii $`\lambda 1145`$ at $`z=0.8913`$. The Ovi was not formally detected in a noisy region of the spectrum. The region of the spectrum corresponding to the Lyman limit is very noisy due to the break at $`z=0.8913`$.
$`z=0.8913`$ — In this system, Feii, Mgi , and Mnii $`\lambda 2577`$ and $`\lambda 2594`$ were detected in the HIRES spectrum. In the FOS spectrum, Siii $`\lambda 1260`$ is blended with Galactic Feii $`\lambda 2383`$, so the $`\lambda 1193`$ transition of the $`\lambda \lambda 1190,1193`$ doublet is used to measure $`W_r(\text{Si}\text{ii})`$. $`W_r(\text{C}\text{ii})`$ was measured using Gaussian deblending, where the blue wing is Siiv $`\lambda 1402`$ at $`z=0.7971`$. Nv $`\lambda 1238`$ is blended with Galactic Feii $`\lambda 2344`$; we obtained the limit on $`W_r(\text{N}\text{v})`$ using $`\lambda 1242`$. The Ovi doublet is ambiguous, being in a very noisy part of the spectrum and possibly being blended with Cii $`\lambda 1334`$ at $`z=0.4720`$; we quote a limit on $`W_r(\text{O}\text{vi})`$. There is a Lyman limit break with a large optical depth whose accurate measurement is difficult.
#### 4.1.15 PG $`1634+706`$ ($`z_{\mathrm{em}}=1.34`$)
The G270H FOS spectrum of this quasar were studied by (Bahcall et al. 1996). This spectrum was taken with the pre–COSTAR optics using a $`0.25\mathrm{}\times 2\mathrm{}`$ slit. A pre–COSTAR spectropolarimetry G190H spectrum was obtained with the wide, $`4.3`$″ by Impey et al. (1996). We did not include the G190H spectrum in our study except to investigate the Lyman limit breaks.
$`z=0.8182`$ — This system has no detected Feii transitions nor Mgi in the HIRES spectrum. There were no detection in the FOS spectrum. The upper limit on $`W_r(\text{Si}\text{ii})`$ is taken from Siii $`\lambda 1526`$, which provided the best limit despite being blended with Siiv $`\lambda 1393`$ at $`z=0.9902`$. The apparent, but slightly offset, Alii $`\lambda 1670`$ is Siii $`\lambda 1526`$ at $`z=0.9902`$, as corroborated by Siii $`\lambda 1260`$. The upper limit on $`W_r(\text{Si}\text{iv})`$ is fairly clean, given the density of lines near the $`\lambda 1393`$ transition; $`\lambda 1402`$ is likely a blend, possibly with a Ly$`\alpha `$ line at $`z=1.0984`$. There is no flux at the position of the Lyman limit due to the strong break from the $`z=0.9902`$ system.
$`z=0.9056`$ — This system has no detected Feii transitions nor Mgi in the HIRES spectrum. In the FOS spectrum, the upper limit on $`W_r(\text{Si}\text{ii})`$ is obtained using $`\lambda 1193`$ because $`\lambda 1260`$ is coincident with Siiii $`\lambda 1206`$ at $`z=0.9902`$. Cii $`\lambda 1334`$ may be present, but the observed equivalent width is $`0.11\pm 0.04`$ Å, which is below the suggested threshold of $`0.13`$ Å for a detection limit in order to avoid flat fielding residuals in the reduced spectrum (Bahcall et al. 1996). (The signal–to–noise ratio of this spectrum is higher than many of the available flat fields for the epoch of these observations of PG $`1634+706`$.) Siiv $`\lambda 1393`$ is coincident with Cii $`\lambda 1334`$ at $`z=0.9902`$; we obtained $`W_r(\text{Si}\text{iv})`$ using $`\lambda 1402`$, which lies in the extended red wing of a possible Ly$`\alpha `$ line at $`z=1.1979`$. The Nv $`\lambda 1238`$ line includes a contribution from Siii $`\lambda 1190`$ at $`z=0.9902`$. This implies that the candidate $`\lambda 1238`$ line is unrealistically strong relative to $`\lambda 1242`$, and also the separation of the centroids of the two lines do not agree. We measure $`W_r(\text{N}\text{v})`$ for the candidate $`\lambda 1238`$ line, but cautiously quote this value only as an upper limit. There is no flux at the position of the Lyman limit due to the strong break from the $`z=0.9902`$ system.
$`z=0.9902`$ — In the HIRES spectrum, this system has several detected Feii transitions and detected Mgi. In the FOS spectrum, there is a Ly$`\alpha `$ line at $`z=0.9785`$ in the red wing of Siiii $`\lambda 1206`$; $`W_r(\text{Si}\text{iii})`$ was obtained using Gaussian deblending. Both members of the Siiv doublet are blended, but easily modeled with Gaussian fitting. The $`\lambda 1393`$ transition is the central member of a three line blend with Ly$`\alpha `$ at $`z=1.2788`$ and $`z=1.2852`$ and the $`\lambda 1393`$ transition is in the blue wing of Galactic Mgii $`\lambda 2796`$ (see Figure 4). The feature just blueward of where Nv $`\lambda 1238`$ is expected is Siiii $`\lambda 1206`$ at $`z=1.0414`$; a $`3\sigma `$ limit was determined at the expected position of $`\lambda 1242`$. There is a strong Lyman limit break with $`\tau >6.0`$ that wipes out all flux blueward of $`1815`$ Å.
$`z=1.0414`$ — This system has no detected Feii transitions nor Mgi in the HIRES spectrum. In the FOS spectrum, the detections are without ambiguity, expect perhaps Cii $`\lambda 1334`$, which could be an artifact of flat fielding problems. The relatively strong feature between Siiv $`\lambda 1393`$ and $`\lambda 1402`$ is Galactic Mgi. Nv $`\lambda 1242`$ is blended with Ly$`\alpha `$ at $`z=1.0881`$, but a $`3\sigma `$ limit was measured at the position of $`\lambda 1238`$. There is a Lyman limit break with $`\tau 1.4`$.
#### 4.1.16 PKS $`2128123`$, PHL 1598 ($`z_{\mathrm{em}}=0.50`$)
Both the G190H and G270H FOS spectra of this quasar have been studied by Jannuzi et al. (1998). The optics were pre–COSTAR and a $`0.25\mathrm{}\times 2\mathrm{}`$ slit was used.
$`z=0.4297`$ — This system is a near–DLA, with $`W_r(\text{Ly}\alpha )=2.92`$ Å. In the HIRES spectrum, Feii was not covered, whereas Mgi and both members of the Caii $`\lambda \lambda 3934,3969`$ were detected. In the FOS spectrum, there is no confusion with Ly$`\alpha `$ absorbers due to the low redshift of the quasar. We quote a detection for the Siiv doublet, since $`\lambda 1393`$ is detected and $`\lambda 1402`$ is consistent with its expected strength (recall that there is no Ly$`\alpha `$ forest confusion, nor are there any other metal–line systems). We obtained an upper limit (albeit not very constraining) on $`W_r(\text{Fe}\text{ii})`$ using $`\lambda 1145`$. The Lyman limit was not covered.
#### 4.1.17 PKS $`2145+067`$ ($`z_{\mathrm{em}}=1.00`$)
Both the G190H and G270H FOS spectra of this quasar have been studied by Bahcall et al. (1993) and by Bergeron et al. (1994). These were obtained with pre–COSTAR optics using the $`0.25\mathrm{}\times 2\mathrm{}`$ slit.
$`z=0.7908`$ — In the HIRES spectrum, Feii was detected, but Mgi was not. From the FOS spectrum, we quote a detection for Siii $`\lambda 1260`$ based upon its correct strength relative to detected Siii $`\lambda 1193`$, although we note Bahcall et al. (1996) identified the line as Ly$`\alpha `$ at $`z=0.8577`$. The line to the red of Siiii $`\lambda 1206`$ is Ly$`\alpha `$ at $`z=0.7810`$, which has a corroborating Ly$`\beta `$ line, thus this is not a doublet at another redshift, and a Siiii detection is claimed. Nv is detected, but we note that $`\lambda 1242`$ is probably blended with another absorption feature. Cii $`\lambda 1334`$ is not detected at $`3\sigma `$, which means that the strong feature to the blue of Ovi $`\lambda 1037`$ is not Cii $`\lambda 1036`$; it is likely Ly$`\alpha `$ at $`z=0.5266`$. The Lyman limit is located in an extremely noisy region of the G190H spectrum and no measurement is possible. See Bergeron et al. (1994) for an independent study of the FOS data for this system.
## 5. Overall Absorption Properties
For most systems, the measured absorption properties are not without some ambiguity in one transition or another. We do not anticipate that highly detailed system by system analyses can be made. That is not to say that, in some cases where robust measurements of extreme properties are measured, physical information cannot be extracted \[e.g. the absorbers toward PG $`1206+459`$ (Churchill & Charlton (1999))\]. However, we are confident that our analysis of the data has yielded measurements that provide a sound representation of the absorption properties of the sample as a whole, given the consistency between our two independent methods for identifying absorption lines (see § 3). Here, we present the general overall properties of the systems and discuss a few of the more immediately obvious trends in the data. We reserve further analysis for Paper II.
In Figure 6, we present the rest–frame absorption strengths of the species listed in Tables Low and High Ionization Absorption Properties of Mgii Absorption–Selected Galaxies at Intermediate Redshifts. I. General Properties<sup>1</sup><sup>1</sup>affiliation: Based in part on observations obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among Caltech, the University of California, and NASA. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. <sup>2</sup><sup>2</sup>affiliation: Based in part on observations obtained with the NASA/ESA Hubble Space Telescope, which is operated by the STScI for the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5–26555. and Low and High Ionization Absorption Properties of Mgii Absorption–Selected Galaxies at Intermediate Redshifts. I. General Properties<sup>1</sup><sup>1</sup>affiliation: Based in part on observations obtained at the W. M. Keck Observatory, which is operated as a scientific partnership among Caltech, the University of California, and NASA. The Observatory was made possible by the generous financial support of the W. M. Keck Foundation. <sup>2</sup><sup>2</sup>affiliation: Based in part on observations obtained with the NASA/ESA Hubble Space Telescope, which is operated by the STScI for the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5–26555. vs. that of Mgii. Only the “flag transitions” are plotted. If a transition other than the “flag transition” was used to measure an equivalent width or an upper limit on an equivalent width, the value was not plotted on any of the figures. However, this results in a negligible difference to the general appearance of Figure 6. For all data, the median errors are roughly the size of the data points and those points with downward pointing arrows represent upper limits. Though there are very few Caii $`\lambda \lambda 3934,3969`$ data points (due to redshift coverage), we have presented Caii absorption strengths for comparison with the numerous studies focused on the Galaxy, and on local and low redshift galaxies (e.g. Morton & Blades (1986); Robertson et al. (1988); Bowen 1991a , 1999b; Vallerga et al. (1993); Sembach, Danks, & Savage (1993); Welty, Morton, & Hobbs (1996)).
Spearman–Kendall non–parametric rank correlation tests, including upper limits (Isobe, Feigelson, & Nelson (1986); LaValley, Isobe, & Feigelson (1992)), show that the absorption strengths of all transitions with ionization potentials less than that of Cii ($`25`$ eV) are correlated with $`W_r(\text{Mg}\text{ii})`$ at a greater than 97% confidence level (with the exception of Caii). We take this as additional confirmation that our measurements, for the most part, provide an accurate representation of the overall sample, since it is expected that these absorption strengths should be correlated. Note, however, that since most of the low ionization transitions are saturated, their equivalent widths are a better measure of the overall kinematic spread of the gas than of column densities (e.g. Petitjean & Bergeron (1990), 1994). An interesting trend with ionization level is that, as $`W_r(\text{Mg}\text{ii})`$ is increased above $`1`$ Å, the low ionization strengths correspondingly increase whereas the higher ionization strengths have considerable scatter and a relatively low mean. Most of the latter data represent the damped Ly$`\alpha `$ systems, as can be seen by their large $`W_r(\text{Ly}\alpha )`$ values. Accurate measurements of Nv and Ovi were fairly rare, but from the few data points available, it would appear that $`W_r(\text{N}\text{v})`$ is often less than $`0.4`$ Å, and that Ovi may have a range of absorption strengths similar to Civ.
Based upon photoionization models (Bergeron & Stasińska (1986); SS92), it has been widely accepted that Mgii absorbers are optically thick at the Lyman limit in neutral hydrogen, i.e. they have neutral hydrogen column densities of $`N(\text{H}\text{i})2\times 10^{17}`$ cm<sup>-2</sup> in one or more of the kinematic components. In Figure 7, we present the rest–frame equivalent widths of Mgii and Civ vs. that of Ly$`\alpha `$. Three data point types are shown. Those systems with a measured Lyman limit break are solid circles, those with no break are open circles, and those for which the break was not covered are open squares. The majority of Mgii absorbers have Ly$`\alpha `$ strengths in the rough range $`0.5`$$`2`$ Å, with a sparsely populated tail extending out to $`10`$ Å. It is not clear if there is a gap in the distribution of $`W_r(\text{Ly}\alpha )`$; this would occur if damped Ly$`\alpha `$ systems \[by definition those with $`N(\text{H}\text{i})2\times 10^{20}`$ cm<sup>-2</sup>, which corresponds to $`W_r(\text{Ly}\alpha )9`$ Å\] are more common than are “Hi rich” systems with $`W_r(\text{Ly}\alpha )`$ in the intermediate range $`2`$$`8`$ Å.. Since these Hi–rich, but non–damped, systems have strong Mgii absorption (most above $`0.6`$ Å), their co–moving number density probably evolves in close step with the strongest Mgii absorbers (SS92).
The Ly$`\alpha `$ transition is easily saturated and is thus not a good measure of the neutral hydrogen column density. The Lyman limit break can provide an estimate of $`N(\text{H}\text{i})`$ and is the unambiguous signature of optically thick neutral hydrogen (Tytler (1982)). The location of Lyman limit break was covered for 17 of the 45 absorbers. All but two of the systems with $`W_r(\text{Ly}\alpha )>1`$ Å have breaks; it is likely that nearly all metal–line systems in this equivalent width range would also have breaks (also see Jannuzi et al. (1998)).
Since the redshift number density of the weakest Mgii systems is several times greater than that of systems with Lyman limit breaks (Churchill et al. 1999a ), it is expected that most Mgii absorbers with $`W_r(\text{Mg}\text{ii})<0.3`$ Å will not have a break. We find two weak systems with breaks, six without, and nine undetermined. We note that the majority of the systems have $`W_r(\text{Ly}\alpha )1.0`$ Å, and none of these for which the location of the Lyman limit was covered has a break. This suggests, not surprisingly, that the weakest of the weak Mgii absorbers are the ones lacking breaks. Note that the presence of a Lyman break is independent of Civ strength in that the full range of observed $`W_r(\text{C}\text{iv})`$ exhibit breaks. Also, note that the Civ strengths are relatively small in the damped Ly$`\alpha `$ systems, implying that their Civ kinematics are probably similar to a more “typical” strong Mgii absorber.
In Figure 8, we present the rest–frame equivalent widths of Siiv and Cii vs. that of Civ. As with Figure 6, we plot “flag transitions” only. In the plot of $`W_r(\text{C}\text{ii})`$ vs. $`W_r(\text{C}\text{iv})`$, we illustrate the combined effects of a range of kinematics and ionization conditions. We also present $`W_r(\text{Si}\text{iv})`$ vs. $`W_r(\text{C}\text{iv})`$, which has been central for studying the ionization conditions in the ISM and Halo of the Galaxy (e.g. Savage, Sembach, & Lu (1997)). Future high resolution spectra of these particular transitions will be essential for a detailed understanding of the kinematics and ionization conditions of the ISM and halos of these intermediate redshift galaxies.
For future detailed studies using higher resolution data, silicon has the added virtue of sampling three ionization states. Thus, we also present the rest–frame equivalent widths of Siiii vs. Siii, and of Siiii and Siii vs. Siiv, in Figure 9. Again, we plot “flag transitions” only. Since both silicon and magnesium are $`\alpha `$–process elements, and Siii and Mgii have nearly identical ionization potentials and transitions with similar oscillator strengths, it is expected that Siii will tightly trace Mgii. The Siiii $`\lambda 1206`$ transition, however, is very strong and is expected to be saturated, even at high resolution.
It is evident that the relative kinematics between the low and high ionization phases vary from absorber to absorber as seen in Figures 8 and 9. For equivalent widths greater than $`0.4`$ Å there is a large range of values (lack of correlations) for Cii and Siiv vs. Civ and for Siii and Siiii vs. Siiv.
## 6. Conclusion
We have measured the absorption properties of 45 intermediate redshift Mgii absorbers in FOS spectra from the HST archive and from the database of the HST QSO Absorption Line Key Project. The sample was selected from the 51 Mgii systems observed with HIRES/Keck for which FOS spectra of the same quasars existed. The Mgii profiles, and other transitions observed in the optical, have been resolved at $`6`$ km s<sup>-1</sup> resolution. The UV FOS spectra have resolution $`230`$ km s<sup>-1</sup>. In this paper, we presented the data, the data analysis, and a brief description of the properties of each system. We present additional analysis of the data in a parallel companion paper (Paper II).
We have found evidence for a high ionization gaseous phase in intermediate redshift Mgii absorbing galaxies. Mostly, the high ionization species detected are Civ and Siiv, which are commonly seen in absorption in the Galaxy (e.g. Savage & Sembach (1996); Savage, Sembach, & Lu (1997)). These data lead us to suggest that these galaxies have multiphase interstellar media and halos similar to those observed locally (Dahlem (1998); also see Churchill et al. 1999b ).
What is the kinematic spread of the Civ and what is its line of sight velocity structure? How is this structure related to that seen in Mgii? High resolution spectroscopic observation are sorely needed for sorting out the physical nature of this high ionization material and its relation to the kinematically complex low ionization gas. Only if the high ionization profiles are resolved at resolutions comparable to the HIRES/Keck data, can the relative kinematics of the low and high ionization gas can be quantified.
Support for this work was provided by the NSF (AST–9617185), and NASA (NAG 5–6399 and AR–07983.01–96A) the latter from the Space Telescope Science Institute, which is operated by AURA, Inc., under NASA contract NAS5–26555. BTJ acknowledges support from NOAO, which is operated by AURA, Inc., under cooperative agreement with the NSF. We thank Blair Savage for providing a table of Galactic absorption properties in electronic form.
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# 1 Fermi and Luttinger liquids
## 1 Fermi and Luttinger liquids
In two and three dimensions, many systems of interacting fermions at low temperatures are described by the Fermi liquid theory developed by Landau (see Ref. 1 for a brief review). According to this theory, at zero temperature, the ground state of each species of fermions has a Fermi surface in momentum space located at an energy called the Fermi energy $`E_F`$, such that all the states within that surface ($`i.e.`$, with energies less than $`E_F`$) are occupied while all the states outside it are unoccupied. An elementary low-energy excitation is one in which a particle is added (annihilated) in a state just outside (inside) the Fermi surface; these are called particle and hole excitations respectively. In an interacting system, these one-particle excitations are accompanied by a cloud of particle-hole pairs, and they are more commonly called quasiparticles; these carry the same charge as a single particle (or hole). If the particle number is held fixed, the low-energy excitations of the system consist of particle-hole pairs in which a certain number of particles are excited from states within the Fermi surface to states outside it. A few of these excitations have both low wave numbers and low energies with the energy being proportional to the wave number; such excitations can be thought of as sound waves. But most of the particle-hole excitations do not have such a linear relationship between energy and wave number; in fact, for most such excitations, a given energy can correspond to many possible momenta.
Another interesting property of a Fermi liquid in two and three dimensions is that the one-particle momentum distribution function $`n(\stackrel{}{k})`$, obtained by Fourier transforming the one-particle equal-time correlation function, has a finite discontinuity at the Fermi surface as shown in Fig. 1 (a). This discontinuity is called the quasiparticle renormalization factor $`z_\stackrel{}{k}`$; it is also equal to the residue of the pole in the one-particle propagator. For non-interacting fermions, $`z_\stackrel{}{k}=1`$; but for interacting fermions, $`0<z_\stackrel{}{k}<1`$ because a quasiparticle is a superposition of many states, only some of which are one-particle excitations. To compute $`z_\stackrel{}{k}`$, we consider the one-particle Green’s function $`G(\stackrel{}{x},t)`$ defined as the expectation value of the time-ordered product of the fermion operator $`\psi (\stackrel{}{x},t)`$ in the ground state $`|0`$, namely,
$$G(\stackrel{}{x},t)=0|T\psi (\stackrel{}{x},t)\psi ^{}(\stackrel{}{0},0)|0.$$
(1)
(We will ignore the spin label here). The Fourier transform of this function can be written as
$$𝒢(\stackrel{}{k},\omega )=\frac{i}{\omega ϵ_\stackrel{}{k}\mathrm{\Sigma }(\stackrel{}{k},\omega )},$$
(2)
where $`ϵ_\stackrel{}{k}`$ is the dispersion relation for the non-interacting theory; we absorb the chemical potential $`\mu `$ in the definition of $`ϵ_\stackrel{}{k}`$ so that $`ϵ_\stackrel{}{k}=0`$ for $`\stackrel{}{k}`$ lying on the Fermi surface. (We will set $`\mathrm{}=1`$). The self-energy $`\mathrm{\Sigma }(\stackrel{}{k},\omega )`$ contains the effects of all the interactions as well as any prescription necessary to shift the pole slightly off the real axis in $`\omega `$. For a Fermi liquid, $`𝒢(\stackrel{}{k},\omega )`$ has a pole near the real axis of $`\omega `$ for any value of $`\stackrel{}{k}`$ on the Fermi surface. In addition, $`\mathrm{\Sigma }`$ is sufficiently analytic at all such points so that the derivative $`\mathrm{\Sigma }/\omega `$ has a finite value. The quasiparticle renormalization factor $`z_\stackrel{}{k}`$ is then given by the residue at the pole, $`i.e.`$,
$$z_\stackrel{}{k}=\left(1\frac{\mathrm{\Sigma }}{\omega }\right)^1.$$
(3)
This gives the discontinuity in $`n(\stackrel{}{k})`$ at the Fermi surface.
Finally, in a Fermi liquid, the various correlation functions decay asymptotically at long distances as power laws, with the exponents being independent of the strength of the interactions. Thus non-interacting and interacting systems have the same exponents and there is a universality.
The discussion above does not apply if the ground state of the system spontaneously breaks some symmetry, for instance, if it is superconducting, or forms a crystal or develops charge or spin density ordering.
In contrast to a Fermi liquid, interacting fermion systems in one dimension behave quite differently ; we will assume again that the ground state breaks no symmetry. Such systems are called Luttinger liquids and they have the following general properties. First of all, there are no single particle or quasiparticle excitations. Thus all the low-energy excitations can be thought of as particle-hole excitations; further, all of these take the form of sound waves with a linear dispersion relation. (As we will see below, there are also excitations of another kind possible which correspond to adding a small number of particles $`N_R`$ and $`N_L`$ to the right and left Fermi points. However, these correspond to only two oscillator degrees of freedom, and therefore do not contribute to thermodynamic properties like the specific heat). Secondly, there is no discontinuity in the momentum distribution function at the Fermi momentum, as indicated in Fig. 1 (b). Rather, there is a cusp there whose form is determined by a certain exponent. Finally, this exponent depends on the strength of the interactions in a non-universal manner, and it also governs the power-law fall-offs of the correlation functions at large space-time distances .
Let us be more specific about the nature of the low-energy excitations in a one-dimensional system of fermions. Assume that we have a system of length $`L`$ with a boundary condition to be specified later. The translation invariance and the finite length make the one-particle momenta discrete. Suppose that the system has $`N_0`$ particles with a ground state energy $`E_0(N_0)`$ and a ground state momentum $`P_0=0`$; we are assuming that the system conserves parity. We will be interested in the thermodynamic limit $`N_0,L\mathrm{}`$ keeping the particle density $`\rho _0=N_0/L`$ fixed. Let us first consider a single species of non-interacting fermions which have two possible directions of motion, right-moving with $`dϵ_k/dk=v_F`$ and left-moving with $`dϵ_k/dk=v_F`$. Here $`ϵ_k`$ is the energy of a low-lying one-particle excitation, $`k`$ is its momentum measured with respect to a right Fermi momentum $`k_F`$ and a left Fermi momentum $`k_F`$ respectively, and $`v_F`$ is called the Fermi velocity. (See Fig. 2 for a typical picture of the momentum states of a lattice model). The values of $`k_F`$ and $`v_F`$ are defined for the non-interacting system; hence they depend on the density $`\rho _0`$ but not on the strength of the interaction. Then a low-lying excitation consists of two pieces ,
(i) a set of bosonic excitations each of which can have either positive momentum $`q`$ or negative momentum $`q`$ with an energy $`ϵ_q=v_Fq`$, where $`0<q<<k_F`$, and
(ii) a certain number of particles $`N_R`$ and $`N_L`$ added to the right and left Fermi points respectively, where $`N_R,N_L<<N_0`$. (Note that $`N_R`$ and $`N_L`$ can be positive, negative or zero. It is convenient to assume that $`N_R\pm N_L`$ are even integers; then the total number of particles $`N_0+N_R+N_L`$ is always even or always odd. We can choose the boundary condition (periodic or antiperiodic) to ensure that the ground state is always non-degenerate).
It turns out that the Hamiltonian and momentum operators for a one-dimensional system (which may have interactions) have the general form
$`H=`$ $`E_0(N_0)+{\displaystyle \underset{q>0}{}}vq[\stackrel{~}{b}_{R,q}^{}\stackrel{~}{b}_{R,q}+\stackrel{~}{b}_{L,q}^{}\stackrel{~}{b}_{L,q}]`$
$`+\mu (N_R+N_L)+{\displaystyle \frac{\pi v}{2LK}}(N_R+N_L)^2+{\displaystyle \frac{\pi vK}{2L}}(N_RN_L)^2,`$
$`P=`$ $`{\displaystyle \underset{q>0}{}}q[\stackrel{~}{b}_{R,q}^{}\stackrel{~}{b}_{R,q}\stackrel{~}{b}_{L,q}^{}\stackrel{~}{b}_{L,q}]+[k_F+{\displaystyle \frac{\pi }{L}}(N_R+N_L)](N_RN_L),`$ (4)
where $`v`$ is the sound velocity, $`q`$ is the momentum of the low-energy bosonic excitations created and annihilated by $`\stackrel{~}{b}_q^{}`$ and $`\stackrel{~}{b}_q`$, $`K`$ is a positive dimensionless number, and $`\mu `$ is the chemical potential of the system. We will see later that $`v`$ and $`K`$ are the two important parameters which determine all the low-energy properties of a system. Their values generally depend on both the strength of the interactions and the density. If the fermions are non-interacting, we have
$$v=v_F\mathrm{and}K=1.$$
(5)
Note that one can numerically find the values of $`v`$ and $`K`$ by studying the $`1/L`$ dependence of the low-energy excitations of finite size systems.
It is interesting that the expression for the momentum operator in Eq. (4) is independent of the interaction strength. We can understand the last term in the momentum as follows. For a continuum system, the Fermi momentum $`k_F(N)`$ is related to the density by the relation
$$L_{k_F(N)}^{k_F(N)}\frac{dk}{2\pi }=N.$$
(6)
Thus a system of $`N_0`$ particles has a Fermi momentum
$$k_F=\frac{\pi N_0}{L}=\pi \rho _0,$$
(7)
while a system of $`N=N_0+N_R+N_L`$ particles has a Fermi momentum equal to $`k_F+(\pi /L)(N_R+N_L)`$. If the $`N`$ particles occupy the momenta states symmetrically about zero momentum, the total momentum of that state is zero; in this state, both the right and left Fermi points have $`(N_R+N_L)/2`$ particles more than the original ground state. Now let us shift $`(N_RN_L)/2`$ particles from the left Fermi point to the right Fermi point, so that the right Fermi point has $`N_R`$ particles more and the left Fermi point has $`N_L`$ particles more than the original system. We then see that the total momentum has changed from zero to $`[k_F+(\pi /L)(N_R+N_L)](N_RN_L)`$; this is the last term in the expression for the momentum operator.
The form of the parameterization of the last two terms in the Hamiltonian in Eq. (4) can be understood as follows. (Note that these two terms vanish in the thermodynamic limit and do not contribute to the specific heat. However they are required for the completeness of the theory up to terms of order $`1/L`$, and for a comparison with conformal field theory). Specifically, we will prove that if the coefficients of $`(\pi /2L)(N_R+N_L)^2`$ and of $`(\pi /2L)(N_RN_L)^2`$ in Eq. (4) are denoted by $`A`$ and $`B`$ respectively, then
$$AB=v^2.$$
(8)
It will then follow that if $`A`$ is equal to $`v/K`$, $`B`$ must be equal to $`vK`$. Although the expressions in Eq. (4) are valid for lattice models also, let us for simplicity consider a continuum model which is invariant under Galilean transformations. First, let us set $`N_R=N_L`$, so that we have added $`\mathrm{\Delta }N=2N_R`$ particles to the system. The sound velocity $`v`$ of a one-dimensional system is related to the density of particles $`\rho =N/L`$ (where $`N=N_0+N_R+N_L`$), the particle mass $`m`$, and the pressure $`𝒫`$ as
$$m\rho v^2=L\left(\frac{𝒫}{L}\right)_N.$$
(9)
The pressure is related to the ground state energy by $`𝒫=(E_0/L)_N`$. Hence
$$m\rho v^2=L\left(\frac{^2E_0}{L^2}\right)_N=\frac{N^2}{L}\left(\frac{^2E_0}{N^2}\right)_L,$$
(10)
where the second equality follows from the first because $`E_0`$ depends on $`N`$ and $`L`$ only through the combination $`N/L`$. Comparing Eqs. (4) and (10), we see that the coefficient of $`(\pi /2L)(\mathrm{\Delta }N)^2`$ is given by
$$A=\frac{mv^2}{\pi \rho _0}.$$
(11)
(In certain expressions such as Eq. (11), we have ignored the difference between $`\rho `$ and $`\rho _0`$ since $`\mathrm{\Delta }NN_0`$). Next, let us take $`N_L=N_R`$; this corresponds to moving $`N_R`$ particles from the left Fermi point $`k_F`$ to the right Fermi point $`k_F`$ keeping the total number of particles equal to $`N_0`$. The change in momentum is therefore given by $`\mathrm{\Delta }P=2\pi \rho _0N_R`$. Since we can also view such an excitation as a center of mass excitation with momentum $`\mathrm{\Delta }P`$, the change in energy is given by $`\mathrm{\Delta }E=(\mathrm{\Delta }P)^2/(2mN)`$ since the total mass of the system is $`mN`$. It follows from this that the coefficient of $`(\pi /2L)(N_RN_L)^2`$ satisfies
$$B=\frac{\pi \rho _0}{m}.$$
(12)
We thus see that $`AB=v^2`$ independently of the nature of the interactions between the particles.
We now consider the other important property of a Luttinger liquid, namely, the absence of a discontinuity in $`n(k)`$ at the Fermi momenta or, equivalently, the absence of a pole in the one-particle propagator. Thus the effect of interactions is so drastic in one dimension that the self-energy $`\mathrm{\Sigma }`$ in Eq. (2) becomes non-analytic at the Fermi points. As a result, $`n(k)`$ becomes continuous at $`k=\pm k_F`$ with the form
$$n(k)=n(k_F)+\mathrm{constant}\mathrm{sign}(kk_F)|kk_F|^\beta ,$$
(13)
where $`\mathrm{sign}(z)1`$ if $`z>0`$, $`1`$ if $`z<0`$ and $`0`$ if $`z=0`$. The exponent $`\beta `$ is a positive number whose value depends on the strength of the interactions; for a non-interacting system, $`\beta =0`$ and we recover the discontinuity in $`n(k)`$. Similarly, the density of states (DOS) is obtained by integrating Eq. (2) over all momenta; near zero energy it vanishes with a power-law form
$$\stackrel{~}{n}(\omega )|\omega |^\beta ,$$
(14)
which signals the absence of one-particle states in the low-energy spectrum. We will see later how the exponent $`\beta `$ can be calculated in an interacting system called the Tomonaga-Luttinger model.
## 2 Bosonization
The basic idea of bosonization is that there are certain objects which can be calculated either in a fermionic theory or in a bosonic theory, and the two calculations give the same answer . Further, a particular quantity may seem very difficult to compute in one theory and may be easily calculable in the other theory. Bosonization works best in two space-time dimensions although there have been some attempts to extend it to higher dimensions.
In two dimensions, bosonization can be studied in either real time (Minkowski space) or in imaginary time (Euclidean space). In both cases, there is a one-to-one correspondence between the correlation functions of some fermionic and bosonic operators. We will work in real time here because bosonization has an added advantage in that case, namely, that there is a direct relationship between the creation and annihilation operators for a boson in terms of the corresponding operators for a fermion . To show this, we just need to consider a bosonic and a fermionic Fock space. A Hamiltonian is not needed at this stage; we need to introduce a Hamiltonian only when discussing interactions and time-dependent correlation functions.
### 2.1 Bosonization of a fermion with one chirality
Let us begin by considering just one component, say, right-moving, of a single species of fermions on a circle of length $`L`$ with the following boundary condition on the one-particle wave functions $`\stackrel{~}{\psi }(x)`$,
$$\stackrel{~}{\psi }(L)=e^{i\pi \sigma }\stackrel{~}{\psi }(0).$$
(15)
Thus $`\sigma =0`$ and $`1`$ correspond to periodic and antiperiodic boundary conditions, but any value of $`\sigma `$ lying in the range $`0\sigma <2`$ is allowed in principle. (If we assume that the particles are charged, then $`\pi \sigma `$ can be identified with an Aharonov-Bohm phase and can be varied by changing the magnetic flux through the circle). The normalized one-particle wave functions are then given by
$`\stackrel{~}{\psi }_{n_k}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{L}}}e^{ikx},`$
$`k`$ $`=`$ $`{\displaystyle \frac{2\pi }{L}}(n_k{\displaystyle \frac{\sigma }{2}}),`$ (16)
where $`n_k=0,\pm 1,\pm 2,\mathrm{}`$ is an integer. We now introduce a second quantized Fermi field
$$\psi _R(x)=\frac{1}{\sqrt{L}}\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}c_{R,k}e^{ikx},$$
(17)
where the subscript $`R`$ stands for right-moving, and
$$\{c_{R,k},c_{R,k^{}}\}=0,\mathrm{and}\{c_{R,k},c_{R,k^{}}^{}\}=\delta _{kk^{}}.$$
(18)
Using the identity
$$\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}e^{iny}=2\pi \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\delta (y2\pi m),$$
(19)
we obtain
$`\{\psi _R(x),\psi _R(x^{})\}`$ $`=`$ $`0,`$
$`\mathrm{and}\{\psi _R(x),\psi _R^{}(x^{})\}`$ $`=`$ $`\delta (xx^{})\mathrm{for}0x,x^{}L.`$ (20)
We define the vacuum or Fermi sea of the system to be the state $`|0`$ satisfying
$`c_{R,k}|0`$ $`=`$ $`0\mathrm{for}k>0,`$
$`c_{R,k}^{}|0`$ $`=`$ $`0\mathrm{for}k0,`$ (21)
as shown in Fig. 3. (Following this definition of the vacuum state, some people prefer to write the particle annihilation operator $`c_{R,k}`$ as a hole creation operator $`d_{R,k}^{}`$ for $`k0`$). Given any operator $`A`$ which can be written as a product of a string of $`c`$’s and $`c^{}`$’s, we denote its normal ordered form by the symbol $`:A:`$. This new operator is defined by moving all the $`c_k`$ with $`k>0`$ and $`c_k^{}`$ with $`k0`$ to the right of all the $`c_k`$ with $`k0`$ and $`c_k^{}`$ with $`k>0`$. This is achieved by transposing as many pairs of creation and annihilation operators as necessary, remembering to multiply by a factor of $`1`$ for each transposition. (It is sometimes claimed that $`:A:=A0|A|0`$. This is true if $`A`$ is quadratic in the $`c`$’s and $`c^{}`$’s, but it is not true in general).
Next we define the fermion number operator
$$\widehat{N}_R=\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}:c_{R,k}^{}c_{R,k}:=\underset{k>0}{}c_{R,k}^{}c_{R,k}\underset{k0}{}c_{R,k}c_{R,k}^{}.$$
(22)
Thus $`\widehat{N}_R|0=0`$. Now consider all possible states $`|\mathrm{\Psi }`$ satisfying $`\widehat{N}_R|\mathrm{\Psi }=0`$. Clearly, any such state can only differ from $`|0`$ by a certain number of particle-hole excitations, $`i.e.`$, it must be of the form
$$|\mathrm{\Psi }=c_{R,k_1}^{}c_{R,k_2}c_{R,k_3}^{}c_{R,k_4}c_{R,k_5}^{}c_{R,k_6}\mathrm{}|0,$$
(23)
where the $`k_i`$ are all different from each other, $`k_1,k_3,\mathrm{}>0`$, and $`k_2,k_4,\mathrm{}0`$. Two such excitations are shown in Fig. 4. We will now see that all such states can be written in terms of certain bosonic creation operators acting on the vacuum. Let us define the operators
$`b_{R,q}^{}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{n}_q}}{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}c_{R,k+q}^{}c_{R,k},`$
$`b_{R,q}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{n}_q}}{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}c_{R,kq}^{}c_{R,k},`$
$`q`$ $`=`$ $`{\displaystyle \frac{2\pi }{L}}n_q,`$ (24)
where $`n_q=1,2,3,\mathrm{}`$. Note that we have defined the boson momentum label $`q`$ to be positive. Also, the fermion boundary condition parameter $`\sigma `$ does not appear in the definitions in Eq. (24). We can check that
$`[\widehat{N}_R,b_{R,q}]`$ $`=`$ $`[\widehat{N}_R,b_{R,q}^{}]=0,`$
$`[b_{R,q},b_{R,q^{}}]`$ $`=`$ $`0,`$
$`[b_{R,q},b_{R,q^{}}^{}]`$ $`=`$ $`\delta _{qq^{}}.`$ (25)
Checking the last identity for $`q=q^{}`$ is slightly tricky due to the presence of an infinite number of fermion momenta $`k`$. One way to derive the commutators is to multiply each $`c_k`$ and $`c_k^{}`$ by a factor of $`\mathrm{exp}[\alpha |k|/2]`$ in Eq. (24), and to let $`\alpha 0`$ at the end of the calculation. We should emphasize that the length scale $`\alpha `$ is not to be thought of as a short-distance cut-off like a lattice spacing; if we had introduced a lattice, the number of fermion modes would have been finite, and the bosonization formulas in Eq. (24) would not have given the correct commutation relations.
We see that the vacuum defined above satisfies $`b_{R,q}|0=0`$ for all $`q`$. If we consider any operator $`A`$ consisting of a string of $`b`$’s and $`b^{}`$’s, we can define its bosonic normal ordered form $`:A:`$ by taking all the $`b_q`$’s to the right of all the $`b_q^{}`$’s by suitable transpositions. Given an operator $`A`$ which can be written in terms of either fermionic or bosonic operators, normal ordering it in the fermionic and bosonic ways do not always give the same result. However, it will always be clear from the context which normal ordering we mean.
We can now begin to understand why bosonization works. First of all, note that there is a one-to-one correspondence between the particle-hole excitations described in Eq. (23) and the bosonic excitations created by the $`b^{}`$’s . For instance, consider a bosonic excitation in which states with the momenta labeled by the integers $`n_1n_2\mathrm{}n_j>0`$ (following the convention in Eq. (24)) are excited. Some of these integers may be equal to each other; that would mean that particular momenta has an occupation number greater than $`1`$. Now we can map this excitation to a fermionic excitation in which $`j`$ fermions occupying the states labeled by the momenta integers $`0,1,2,\mathrm{},j+1`$ (following the convention in Eq. (16)) are excited to momenta labeled by $`n_1,n_21,n_32,\mathrm{},n_jj+1`$ respectively. This is clearly a one-to-one map, and we can reverse it to uniquely obtain a bosonic excitation from a given fermionic excitation. This mapping allows us to show, once an appropriate Hamiltonian is defined, that thermodynamic quantities like the specific heat are identical in the fermionic and bosonic models.
The above mapping makes it plausible, although it requires more effort to prove, that all particle-hole excitations can be produced by combinations of $`b^{}`$’s acting on the vacuum. For instance, the state in Fig. 4 (a) is given by $`b_{R,1}^{}|0`$. However the state in Fig. 4 (b) has a more lengthy expression in terms of bosonic operators, although it is also a single particle-hole excitation just like (a); to be explicit, it is given by the linear combination $`(1/6)[2b_{R,3}^{}+3b_{R,2}^{}b_{R,1}^{}+(b_{R,1}^{})^3]|0`$.
Next, we define bosonic field operators and show that some bilinears in fermionic fields, such as the density $`\rho _R(x)`$, have simple expressions in terms of bosonic fields. Define the fields
$`\chi _R(x)`$ $`=`$ $`{\displaystyle \frac{i}{2\sqrt{\pi }}}{\displaystyle \underset{q>0}{}}{\displaystyle \frac{1}{\sqrt{n}_q}}b_{R,q}e^{iqx\alpha q/2},`$
$`\chi _R^{}(x)`$ $`=`$ $`{\displaystyle \frac{i}{2\sqrt{\pi }}}{\displaystyle \underset{q>0}{}}{\displaystyle \frac{1}{\sqrt{n}_q}}b_{R,q}^{}e^{iqx\alpha q/2},`$
$`\varphi _R(x)`$ $`=`$ $`\chi _R(x)+\chi _R^{}(x){\displaystyle \frac{\sqrt{\pi }x}{L}}\widehat{N}_R.`$ (26)
The last term in the definition of $`\varphi _R(x)`$ has been put in for later convenience; it simplifies the expressions for the Hamiltonian and the fermion density in terms of $`\varphi _R`$. (Some authors prefer not to include that term in the definition of $`\varphi _R`$ but add it separately in the Hamiltonian and density). Note that $`\widehat{N}_R`$ commutes with both $`\chi _R`$ and $`\chi _R^{}`$. From the commutation relations in Eq. (25), we see that
$`[\chi _R(x),\chi _R(x^{})]`$ $`=`$ $`0,`$
$`[\chi _R(x),\chi _R^{}(x^{})]`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}\mathrm{ln}[1\mathrm{exp}({\displaystyle \frac{2\pi }{L}}(\alpha +i(xx^{}))],`$ (27)
$`=`$ $`{\displaystyle \frac{1}{4\pi }}\mathrm{ln}[{\displaystyle \frac{2\pi }{L}}(\alpha +i(xx^{}))]\mathrm{in}\mathrm{the}\mathrm{limit}L\mathrm{}.`$
Henceforth, the limit $`L\mathrm{}`$ will be assumed wherever convenient. We find that
$`[\varphi _R(x),\varphi _R(x^{})]`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}\mathrm{ln}\left[{\displaystyle \frac{\alpha i(xx^{})}{\alpha +i(xx^{})}}\right],`$ (28)
$`=`$ $`{\displaystyle \frac{i}{4}}\mathrm{sign}(xx^{})\mathrm{in}\mathrm{the}\mathrm{limit}\alpha 0.`$
Thus the commutator of two $`\varphi `$’s looks like a step function which is smeared over a region of length $`\alpha `$.
Now we use the operator identity
$$\mathrm{exp}A\mathrm{exp}B=\mathrm{exp}(A+B+\frac{1}{2}[A,B]),$$
(29)
if $`[A,B]`$ commutes with both $`A`$ and $`B`$. It follows that
$$\mathrm{exp}[i2\sqrt{\pi }\chi _R^{}(x)]\mathrm{exp}[i2\sqrt{\pi }\chi _R(x)]\mathrm{exp}[i\frac{2\pi x}{L}\widehat{N}_R]=\left(\frac{L}{2\pi \alpha }\right)^{1/2}\mathrm{exp}[i2\sqrt{\pi }\varphi _R(x)].$$
(30)
We observe that the left hand side of this equation is normal ordered while the right hand side is not; that is why the two sides are related through a divergent factor involving $`L/\alpha `$.
We can show that the fermion density operator is linear in the bosonic field, namely,
$`\rho _R(x)`$ $`=`$ $`:\psi _R^{}(x)\psi _R(x):`$ (31)
$`=`$ $`{\displaystyle \frac{1}{L}}{\displaystyle \underset{q>0}{}}\sqrt{n}_q(b_{R,q}e^{iqx}+b_{R,q}^{}e^{iqx})+{\displaystyle \frac{1}{L}}{\displaystyle \underset{k}{}}:c_{R,k}^{}c_{R,k}:`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{\pi }}}{\displaystyle \frac{\varphi _R}{x}}.`$
We now go in the opposite direction and construct fermionic field operators from bosonic ones. To do this, we first define the Klein factors $`\eta _R`$ and $`\eta _R^{}`$ which are unitary operators satisfying
$`[\widehat{N}_R,\eta _R^{}]`$ $`=`$ $`\eta _R^{},[\widehat{N}_R,\eta _R]=\eta _R,`$
$`[\eta _R,b_{R,q}]`$ $`=`$ $`[\eta _R,b_{R,q}^{}]=0.`$ (32)
Pictorially, in terms of Figs. 3 and 4, the action of $`\eta _R^{}`$ is to raise all the occupied fermion states by one unit of momentum, while the action of $`\eta _R`$ is to lower all the fermion occupied states by one unit of momentum. Although these actions are easy to describe in words, the explicit expressions for $`\eta _R`$ and $`\eta _R^{}=\eta _{R}^{}{}_{}{}^{1}`$ in terms of the $`c`$’s and $`c^{}`$’s are rather complicated. The Klein factors will be needed to ensure the correct anticommutation relations between the fermionic operators constructed below.
We observe that
$`[b_{R,q},\psi _R(x)]`$ $`=`$ $`{\displaystyle \frac{e^{iqx}}{\sqrt{n}_q}}\psi _R(x),`$
$`[b_{R,q}^{},\psi _R(x)]`$ $`=`$ $`{\displaystyle \frac{e^{iqx}}{\sqrt{n}_q}}\psi _R(x).`$ (33)
Since $`b_{R,q}`$ annihilates the vacuum, we have
$$b_{R,q}\psi _R(x)|0=\frac{e^{iqx}}{\sqrt{n}_q}\psi _R(x)|0.$$
(34)
Thus $`\psi _R(x)|0`$ is an eigenstate of $`b_{R,q}`$ for every value of $`q`$, namely, it is a coherent state. We therefore make the ansatz
$`\psi _R(x)|0`$ $`=`$ $`Q(x)\mathrm{exp}[{\displaystyle \underset{q>0}{}}{\displaystyle \frac{e^{iqx}}{\sqrt{n}_q}}b_{R,q}^{}|0]|0,`$ (35)
$`=`$ $`Q(x)\mathrm{exp}[i2\sqrt{\pi }\chi _R^{}(x)]|0,`$
where $`Q(x)`$ is some operator which commutes with all the $`b`$’s and $`b^{}`$’s. Since $`\psi _R`$ reduces the fermion number by $`1`$, $`Q`$ must contain a factor of $`\eta _R`$. Let us try the form $`Q(x)=F(x)\eta _R`$, where $`F(x)`$ is a $`c`$-number function of $`x`$. The form of $`F`$ is determined by computing
$`F(x)`$ $`=`$ $`0|\eta _R^{}\eta _RF(x)|0`$ (36)
$`=`$ $`0|\eta _R^{}\psi _R(x)|0`$
$`=`$ $`{\displaystyle \frac{e^{i\pi \sigma x/L}}{\sqrt{L}}}.`$
(The last line in Eq. (36) has been derived by using the actions of $`\eta _R^{}`$ above and of $`\psi _R`$ in Eq. (17). To see this explicitly, note that $`0|\eta _R^{}`$ is the conjugate of the state in which the top most fermion has been removed from the vacuum. Hence, in $`\psi _R|0`$, we only have to consider the state in which the top most fermion has been removed; so we require the wave function of the state with $`n_k=0`$ in Eq. (16)). We now obtain
$`\psi _R(x)|0`$ $`=`$ $`{\displaystyle \frac{e^{i\pi \sigma x/L}}{\sqrt{2\pi \alpha }}}\eta _Re^{i2\sqrt{\pi }\varphi _R(x)}|0,`$ (37)
where we have used Eq. (30) and the fact that $`\chi _R(x)`$ and $`\widehat{N}_R`$ annihilate the vacuum. We are thus led to the plausible conjecture
$$\psi _R(x)=\frac{e^{i\pi \sigma x/L}}{\sqrt{2\pi \alpha }}\eta _Re^{i2\sqrt{\pi }\varphi _R(x)}.$$
(38)
To prove this, we need to show that the two sides of this equation have the same action on all states, not just the vacuum. Such a proof is given in Ref. 7. Eq. (38) is one of the most important identities in bosonization.
We next introduce a non-interacting Hamiltonian by defining the energy of the fermion mode with momentum $`k`$ to be
$$ϵ_k=v_Fk$$
(39)
for all values of $`k`$. The Hamiltonian is
$`H_0`$ $`=`$ $`v_F{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}k:c_{R,k}^{}c_{R,k}:+{\displaystyle \frac{\pi v_F}{L}}\widehat{N}_R^2`$ (40)
$`=`$ $`v_F{\displaystyle _0^L}𝑑x:\psi _R^{}i_x\psi _R:+{\displaystyle \frac{\pi v_F}{L}}\widehat{N}_R^2.`$
This defines the chiral Luttinger model. (The term proportional to $`\widehat{N}_R^2`$ has been introduced in Eq. (40) so as to reproduce similar terms in Eq. (4) after we introduce left-moving fields in the next section). We can check that $`H_0|0=0`$ and
$`[H_0,b_{R,q}]`$ $`=`$ $`v_Fqb_{R,q},`$
$`[H_0,b_{R,q}^{}]`$ $`=`$ $`v_Fqb_{R,q}^{},`$ (41)
To reproduce these relations in the bosonic language, we must have
$`H_0`$ $`=`$ $`v_F{\displaystyle \underset{q>0}{}}qb_{R,q}^{}b_{R,q}+{\displaystyle \frac{\pi v_F}{L}}\widehat{N}_R^2`$ (42)
$`=`$ $`v_F{\displaystyle _0^L}𝑑x:(_x\varphi )^2:.`$
We can introduce an interaction in this model which is quadratic in the fermion density. Let us consider the interaction
$$V=\frac{1}{2}_0^Lg_4\rho _{R}^{}{}_{}{}^{2}(x)=\frac{g_4}{2\pi }\underset{q>0}{}qb_{R,q}^{}b_{R,q}.$$
(43)
Physically, such a term could arise if there is a short-range ($`i.e.`$, screened) Coulomb repulsion or a phonon mediated attraction between two fermions. We will therefore not make any assumptions about the sign of the interaction parameter $`g_4`$. If we add Eq. (43) to Eq. (42), we see that the only effect of the interaction in this model is to renormalize the velocity from $`v_F`$ to $`v_F+(g_4/2\pi )`$.
In the next section, we will consider a model containing fermions with opposite chiralities; we will then see that a density-density interaction can have more interesting effects than just renormalizing the velocity.
### 2.2 Bosonization of a fermion with two chiralities
Let us consider a fermion with both right- and left-moving components as depicted in Figs. 3 and 5 respectively. For the left-moving fermions in Fig. 5, we define the momentum label $`k`$ as increasing towards the left; the advantage of this choice is that the vacuum has the negative $`k`$ states occupied and the positive $`k`$ states unoccupied for both chiralities. We introduce a chirality label $`\nu `$, such that $`\nu =R`$ and $`L`$ refer to right- and left-moving particles respectively. Sometimes we will use the numerical values $`\nu =1`$ and $`1`$ for $`R`$ and $`L`$; this will be clear from the context. Let us choose periodic boundary conditions on the circle so that $`\sigma =0`$. Then the Fermi fields are given by
$`\psi _\nu (x)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{L}}}{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}c_{\nu ,k}e^{i\nu kx},`$
$`k`$ $`=`$ $`{\displaystyle \frac{2\pi }{L}}n_k,`$ (44)
where $`n_k=0,\pm 1,\pm 2,\mathrm{}`$, and
$`\{c_{\nu ,k},c_{\nu ^{},k^{}}\}`$ $`=`$ $`0,`$
$`\{c_{\nu ,k},c_{\nu ^{},k^{}}^{}\}`$ $`=`$ $`\delta _{\nu \nu ^{}}\delta _{kk^{}}.`$ (45)
The vacuum is defined as the state satisfying
$`c_{\nu ,k}|0=0\mathrm{for}k>0,`$
$`c_{\nu ,k}^{}|0=0\mathrm{for}k0.`$ (46)
We can then define normal ordered fermion number operators $`\widehat{N}_\nu `$ in the usual way.
Next we define bosonic operators
$`b_{\nu ,q}^{}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{n}_q}}{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}c_{\nu ,k+q}^{}c_{\nu ,k},`$
$`b_{\nu ,q}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{n}_q}}{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}c_{\nu ,kq}^{}c_{\nu ,k}.`$ (47)
Note that $`b_{R,q}^{}`$ and $`b_{L,q}^{}`$ create excitations with momenta $`q`$ and $`q`$ respectively, where the label $`q`$ is always taken to be positive. We can show as before that
$$[b_{\nu ,q},b_{\nu ^{},q^{}}]=0,\mathrm{and}[b_{\nu ,q},b_{\nu ^{},q^{}}^{}]=\delta _{\nu \nu ^{}}\delta _{qq^{}}.$$
(48)
The unitary Klein operators $`\eta _\nu `$ ($`\eta _\nu ^{}`$) are defined to be operators which raise (lower) the momentum label $`k`$ of all the occupied states for fermions of type $`\nu `$. We then have
$`\{\eta _R,\eta _L\}`$ $`=`$ $`\{\eta _R,\eta _L^{}\}=0,`$
$`[\widehat{N}_\nu ,\eta _\nu ^{}^{}]`$ $`=`$ $`\delta _{\nu \nu ^{}}\eta _\nu ^{}^{},[\widehat{N}_\nu ,\eta _\nu ^{}]=\delta _{\nu \nu ^{}}\eta _\nu ^{},`$
$`[\eta _\nu ,b_{\nu ^{},q}]`$ $`=`$ $`[\eta _\nu ,b_{\nu ^{},q}^{}]=0.`$ (49)
We now define the chiral creation and annihilation fields
$`\chi _\nu (x)`$ $`=`$ $`{\displaystyle \frac{i\nu }{2\sqrt{\pi }}}{\displaystyle \underset{q>0}{}}{\displaystyle \frac{1}{\sqrt{n}_q}}b_{\nu ,q}e^{i\nu qx\alpha q/2},`$
$`\chi _\nu ^{}(x)`$ $`=`$ $`{\displaystyle \frac{i\nu }{2\sqrt{\pi }}}{\displaystyle \underset{q>0}{}}{\displaystyle \frac{1}{\sqrt{n}_q}}b_{\nu ,q}^{}e^{i\nu qx\alpha q/2}.`$ (50)
Then
$`[\chi _\nu (x),\chi _\nu ^{}(x^{})]`$ $`=`$ $`0,`$
$`[\chi _\nu (x),\chi _\nu ^{}^{}(x^{})]`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}\delta _{\nu \nu ^{}}\mathrm{ln}[{\displaystyle \frac{2\pi }{L}}(\alpha i\nu (xx^{}))]\mathrm{in}\mathrm{the}\mathrm{limit}L\mathrm{}.`$ (51)
The chiral fields
$$\varphi _\nu (x)=\chi _\nu (x)+\chi _\nu ^{}(x)\frac{\sqrt{\pi }x}{L}\widehat{N}_\nu ,$$
(52)
satisfy
$$[\varphi _\nu (x),\varphi _\nu ^{}(x)]=\frac{i\nu }{4}\delta _{\nu \nu ^{}}\mathrm{sign}(xx^{})$$
(53)
in the limit $`\alpha 0`$. Finally, we can define two fields dual to each other
$`\varphi (x)`$ $`=`$ $`\varphi _R(x)+\varphi _L(x),`$
$`\theta (x)`$ $`=`$ $`\varphi _R(x)+\varphi _L(x),`$ (54)
such that $`[\varphi (x),\varphi (x^{})]=[\theta (x),\theta (x^{})]=0`$, while
$$[\varphi (x),\theta (x^{})]=\frac{i}{2}\mathrm{sign}(xx^{}).$$
(55)
The fermion density operators $`\rho _\nu (x)=:\psi _\nu ^{}(x)\psi _\nu (x):`$ satisfy $`\rho _\nu =_x\varphi _\nu /\sqrt{\pi }`$. Hence the total density and current operators are given by
$`\rho (x)`$ $`=`$ $`\rho _R+\rho _L={\displaystyle \frac{1}{\sqrt{\pi }}}_x\varphi ,`$
$`j(x)`$ $`=`$ $`v_F(\rho _R\rho _L)={\displaystyle \frac{v_F}{\sqrt{\pi }}}_x\theta ,`$ (56)
where $`v_F`$ is a velocity to be introduced shortly.
We can again show that the fermionic fields are given in terms of the bosonic ones as
$`\psi _R(x)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi \alpha }}}\eta _Re^{i2\sqrt{\pi }\varphi _R},`$
$`\psi _L(x)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi \alpha }}}\eta _Le^{i2\sqrt{\pi }\varphi _L}.`$ (57)
As before, we introduce a linear dispersion relation $`ϵ_{\nu ,k}=v_Fk`$ for the fermions. The non-interacting Hamiltonian then takes the form
$`H_0`$ $`=`$ $`v_F{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}k[:c_{R,k}^{}c_{R,k}+c_{L,k}^{}c_{L,k}:]+{\displaystyle \frac{\pi v_F}{L}}(\widehat{N}_R^2+\widehat{N}_L^2)`$ (58)
$`=`$ $`v_F{\displaystyle _0^L}dx[:\psi _R^{}(x)i_x\psi _R(x)\psi _L^{}(x)i_x\psi _L(x):]+{\displaystyle \frac{\pi v_F}{L}}(\widehat{N}_R^2+\widehat{N}_L^2)`$
in the fermionic language, and
$`H_0`$ $`=`$ $`v_F{\displaystyle \underset{q>0}{}}q(b_{R,q}^{}b_{R,q}+b_{L,q}^{}b_{L,q})+{\displaystyle \frac{\pi v_F}{L}}(\widehat{N}_R^2+\widehat{N}_L^2)`$ (59)
$`=`$ $`v_F{\displaystyle _0^L}dx[:(_x\varphi _R)^2+(_x\varphi _L)^2:]`$
$`=`$ $`{\displaystyle \frac{v_F}{2}}{\displaystyle _0^L}dx[:(_x\varphi )^2+(_x\theta )^2:]`$
in the bosonic language. If we use this Hamiltonian to transform all the fields to time-dependent Heisenberg fields, we find that $`\psi _R,\varphi _R`$ become functions of $`x_R=xv_Ft`$ while $`\psi _L,\varphi _L`$ become functions of $`x_L=x+v_Ft`$.
From Eq. (55), we see that the field canonically conjugate to $`\varphi `$ is given by
$$\mathrm{\Pi }=_x\theta .$$
(60)
Thus
$$[\varphi (x),\mathrm{\Pi }(x^{})]=i\delta (xx^{}),$$
(61)
and
$$H_0=\frac{v_F}{2}_0^L𝑑x[\mathrm{\Pi }^2+(_x\varphi )^2].$$
(62)
We now study the effects of four-fermi interactions. In the beginning it is simpler to work in the Schrödinger representation in which the fields are time-independent; we will transform to the Heisenberg representation later when we compute the correlation functions. Let us consider an interaction of the form
$$V=\frac{1}{2}_0^L𝑑x[2g_2\rho _R(x)\rho _L(x)+g_4(\rho _R^2(x)+\rho _L^2(x))].$$
(63)
Physically, we may expect an interaction such as $`g:\rho ^2(x):`$, so that $`g_2=g_4=g`$. However, it is instructive to allow $`g_2`$ to differ from $`g_4`$ to see what happens. For reasons explained before, we will again not assume anything about the signs of $`g_2`$ and $`g_4`$. In the fermionic language, the interaction takes the form
$`V={\displaystyle \frac{1}{2L}}{\displaystyle \underset{k_1,k_2,k_3=\mathrm{}}{\overset{\mathrm{}}{}}}[`$ $`2g_2c_{R,k_1+k_3}^{}c_{R,k_1}c_{L,k_2+k_3}^{}c_{L,k_2}`$ (64)
$`+g_4(c_{R,k_1+k_3}^{}c_{R,k_1}c_{R,k_2k_3}^{}c_{R,k_2}+c_{L,k_1+k_3}^{}c_{L,k_1}c_{L,k_2k_3}^{}c_{L,k_2})].`$
From this expression we see that $`g_2`$ corresponds to a two-particle scattering involving both chiralities; in this model, we can call it either forward scattering or backward scattering since there is no way to distinguish between the two processes in the absence of some other quantum number such as spin. The $`g_4`$ term corresponds to a scattering between two fermions with the same chirality, and therefore describes a forward scattering process.
The quartic interaction in Eq. (64) seems very difficult to analyze. However we will now see that it is easily solvable in the bosonic language; indeed this is one of the main motivations behind bosonization. The bosonic expression for the total Hamiltonian $`H=H_0+V`$ is found to be
$`H=`$ $`{\displaystyle \underset{q>0}{}}q[v_F(b_{R,q}^{}b_{R,q}+b_{L,q}^{}b_{L,q})+{\displaystyle \frac{g_2}{2\pi }}(b_{R,q}^{}b_{L,q}^{}+b_{R,q}b_{L,q})+{\displaystyle \frac{g_4}{2\pi }}(b_{R,q}^{}b_{R,q}+b_{L,q}^{}b_{L,q})]`$ (65)
$`+{\displaystyle \frac{\pi v_F}{L}}(\widehat{N}_R^2+\widehat{N}_L^2)+{\displaystyle \frac{g_2}{L}}\widehat{N}_R\widehat{N}_L+{\displaystyle \frac{g_4}{2L}}(\widehat{N}_R^2+\widehat{N}_L^2).`$
The $`g_4`$ term again renormalizes the velocity. The $`g_2`$ term can then be rediagonalized by a Bogoliubov transformation. We first define two parameters
$`v`$ $`=`$ $`\left[(v_F+{\displaystyle \frac{g_4}{2\pi }}{\displaystyle \frac{g_2}{2\pi }})(v_F+{\displaystyle \frac{g_4}{2\pi }}+{\displaystyle \frac{g_2}{2\pi }})\right]^{1/2},`$
$`K`$ $`=`$ $`\left[(v_F+{\displaystyle \frac{g_4}{2\pi }}{\displaystyle \frac{g_2}{2\pi }})/(v_F+{\displaystyle \frac{g_4}{2\pi }}+{\displaystyle \frac{g_2}{2\pi }})\right]^{1/2}.`$ (66)
Note that $`K<1`$ if $`g_2`$ is positive (repulsive interaction), and $`>1`$ if $`g_2`$ is negative (attractive interaction). (If $`g_2`$ is so large that $`v_F+g_4/(2\pi )g_2/(2\pi )<0`$, then our analysis breaks down. The system does not remain a Luttinger liquid in that case, and is likely to go into a different phase such as a state with charge density order). The Bogoliubov transformation now takes the form
$`\stackrel{~}{b}_{R,q}`$ $`=`$ $`{\displaystyle \frac{b_{R,q}+\gamma b_{L,q}^{}}{\sqrt{1\gamma ^2}}},`$
$`\stackrel{~}{b}_{L,q}`$ $`=`$ $`{\displaystyle \frac{b_{L,q}+\gamma b_{R,q}^{}}{\sqrt{1\gamma ^2}}},`$
$`\gamma `$ $`=`$ $`{\displaystyle \frac{1K}{1+K}},`$ (67)
for each value of the momentum $`q`$. The Hamiltonian is then given by the quadratic expression
$`H=`$ $`{\displaystyle \underset{q>0}{}}vq[\stackrel{~}{b}_{R,q}^{}\stackrel{~}{b}_{R,q}+\stackrel{~}{b}_{L,q}^{}\stackrel{~}{b}_{L,q}]`$ (68)
$`+{\displaystyle \frac{\pi v}{2L}}[{\displaystyle \frac{1}{K}}(\widehat{N}_R+\widehat{N}_L)^2+K(\widehat{N}_R\widehat{N}_L)^2].`$
Equivalently,
$$H=\frac{1}{2}_0^L𝑑x[vK\mathrm{\Pi }^2+\frac{v}{K}(_x\varphi )^2].$$
(69)
The old and new fields are related as
$`\varphi _R`$ $`=`$ $`{\displaystyle \frac{(1+K)\stackrel{~}{\varphi }_R(1K)\stackrel{~}{\varphi }_L}{2\sqrt{K}}},`$
$`\varphi _L`$ $`=`$ $`{\displaystyle \frac{(1+K)\stackrel{~}{\varphi }_L(1K)\stackrel{~}{\varphi }_R}{2\sqrt{K}}},`$
$`\varphi `$ $`=`$ $`\sqrt{K}\stackrel{~}{\varphi }\mathrm{and}\theta ={\displaystyle \frac{1}{\sqrt{K}}}\stackrel{~}{\theta }.`$ (70)
Note the important fact that the vacuum changes as a result of the interaction; the new vacuum $`|\stackrel{~}{0}`$ is the state annihilated by the operators $`\stackrel{~}{b}_{\nu ,q}`$. Since the various correlation functions must be calculated in this new vacuum, they will depend on the interaction through the parameters $`v`$ and $`K`$. In particular, we will see in the next section that the power-laws of the correlation functions are governed by $`K`$.
Given the various Hamiltonians, it is easy to guess the forms of the corresponding Lagrangians. For the non-interacting theory ($`g_2=g_4=0`$), the Lagrangian density describes a massless Dirac fermion,
$$=i\psi _R^{}(_t+v_F_x)\psi _R+i\psi _L^{}(_tv_F_x)\psi _L$$
(71)
in the fermionic language, and a massless real scalar field,
$$=\frac{1}{2v_F}(_t\varphi )^2\frac{v_F}{2}(_x\varphi )^2$$
(72)
in the bosonic language. For the interacting theory in Eq. (69), we find from Eq. (70) that
$$=\frac{1}{2vK}(_t\varphi )^2\frac{v}{2K}(_x\varphi )^2=\frac{1}{2v}(_t\stackrel{~}{\varphi })^2\frac{v}{2}(_x\stackrel{~}{\varphi })^2.$$
(73)
The momentum operator in Eq. (4) has the same expression in terms of the old and new fields, namely,
$$P=k_F(\widehat{N}_R\widehat{N}_L)+_0^L𝑑x_x\varphi _x\theta .$$
(74)
We can check that $`[P,\varphi ]=i_x\varphi `$ and $`[P,\theta ]=i_x\theta `$.
Let us now write down the fields $`\stackrel{~}{\varphi }`$ and $`\stackrel{~}{\theta }`$ in the Heisenberg representation. This is simple to do once we realize that the right- and left-moving fields must be functions of $`x_R=xvt`$ and $`x_L=x+vt`$ respectively. We find that
$`\stackrel{~}{\varphi }(x,t)={\displaystyle \frac{i}{2\sqrt{\pi }}}{\displaystyle \underset{q>0}{}}{\displaystyle \frac{1}{\sqrt{n}_q}}[`$ $`\stackrel{~}{b}_{R,q}e^{iq(x_R+i\alpha /2)}\stackrel{~}{b}_{R,q}^{}e^{iq(x_Ri\alpha /2)}\stackrel{~}{b}_{L,q}e^{iq(x_Li\alpha /2)}`$
$`+\stackrel{~}{b}_{L,q}^{}e^{iq(x_L+i\alpha /2)}]`$
$`{\displaystyle \frac{\sqrt{\pi }}{L}}[`$ $`{\displaystyle \frac{x}{\sqrt{K}}}(\widehat{N}_R+\widehat{N}_L)\sqrt{K}vt(\widehat{N}_R\widehat{N}_L)],`$
$`\stackrel{~}{\theta }(x,t)={\displaystyle \frac{i}{2\sqrt{\pi }}}{\displaystyle \underset{q>0}{}}{\displaystyle \frac{1}{\sqrt{n}_q}}[`$ $`\stackrel{~}{b}_{R,q}e^{iq(x_R+i\alpha /2)}+\stackrel{~}{b}_{R,q}^{}e^{iq(x_Ri\alpha /2)}\stackrel{~}{b}_{L,q}e^{iq(x_Li\alpha /2)}`$
$`+\stackrel{~}{b}_{L,q}^{}e^{iq(x_L+i\alpha /2)}]`$
$`+{\displaystyle \frac{\sqrt{\pi }}{L}}[`$ $`\sqrt{K}x(\widehat{N}_R\widehat{N}_L){\displaystyle \frac{vt}{\sqrt{K}}}(\widehat{N}_R+\widehat{N}_L)].`$ (75)
We observe that the coefficients of $`\widehat{N}_R`$ and $`\widehat{N}_L`$ have terms which are linear in $`t`$. This is necessary because we want the conjugate momentum field to satisfy
$$\stackrel{~}{\mathrm{\Pi }}=\frac{1}{v}_t\stackrel{~}{\varphi }=_x\stackrel{~}{\theta }.$$
(76)
We note that a dual equation holds, namely,
$$\frac{1}{v}_t\stackrel{~}{\theta }=_x\stackrel{~}{\varphi }.$$
(77)
(One can check from Eq. (75) that $`\stackrel{~}{\varphi }_R`$ and $`\stackrel{~}{\varphi }_L`$ are functions of $`x_R`$ and $`x_L`$ alone). In terms of $`\theta `$, the Lagrangian density is
$$=\frac{K}{2v}(_t\theta )^2\frac{Kv}{2}(_x\theta )^2=\frac{1}{2v}(_t\stackrel{~}{\theta })^2\frac{v}{2}(_x\stackrel{~}{\theta })^2.$$
(78)
Although the Lagrangians in Eq. (73) and Eq. (78) have opposite signs, the Hamiltonians derived from the two are identical.
Before ending this section, let us comment on a global symmetry of all these models. It is known that fermionic systems with a conserved charge are invariant under a global phase rotation
$$\psi _Re^{i\lambda }\psi _R,\mathrm{and}\psi _Le^{i\lambda }\psi _L,$$
(79)
where $`\lambda `$ is independent of $`(x,t)`$. Eq. (57) then implies that the corresponding bosonic theories must remain invariant under
$$\varphi \varphi ,\mathrm{and}\theta \theta +\frac{\lambda }{\sqrt{\pi }}.$$
(80)
This provides a constraint on the kinds of terms which can appear in the Lagrangians of such theories.
### 2.3 Field theory of modes near the Fermi momenta
In the last section, we discussed bosonization for a model of fermions which has the following properties.
(i) There are an infinite number of right- and left-moving modes with the momenta going from $`\mathrm{}`$ to $`\mathrm{}`$, and
(ii) the relation between energy and momentum is linear for all values of the momentum.
Neither of these properties is true in condensed matter systems which typically are non-relativistic and have a finite (though possibly very large) number of states. The question is the following: can bosonization give useful results even if these two properties do not hold? We will see that the answer is yes, provided that we are only interested in the long-wavelength, low-frequency and low-temperature properties of such systems.
In an experimental system, the fermions may be able to move either on a discrete lattice of points such as in a crystal, or in a continuum such as the conduction electrons in a metal. For instance, non-interacting fermions moving in a continuum have a dispersion $`ϵ_k=k^2/2m`$, while fermions hopping on a lattice have a dispersion such as $`ϵ_k=t\mathrm{cos}(ka)`$ if $`a`$ is the lattice spacing and $`t`$ is the nearest neighbor hopping amplitude. In either case, a non-interacting system in one-dimension will, at zero temperature, have a Fermi surface consisting of two points in momentum space given by $`k=\pm k_F`$ (see Fig. 2). As stated before, we define the one-particle energy to be zero at the Fermi points. At low temperatures $`T`$ or low frequencies $`\omega `$, the only modes which can contribute are the ones lying close to those points, $`i.e.`$, with excitation energies of the order of or smaller than $`k_BT`$ or $`\omega `$. Near the Fermi points, we can approximate the dispersion relation by a linear one, with the velocity being defined to be $`v_F=(dϵ_k/dk)_{k=k_F}`$. We thus restrict our attention to the right-moving modes with momenta lying between $`k_F\mathrm{\Lambda }`$ and $`k_F+\mathrm{\Lambda }`$, and the left-moving modes with momenta lying between $`k_F\mathrm{\Lambda }`$ and $`k_F+\mathrm{\Lambda }`$. Here $`\mathrm{\Lambda }`$ is taken to be much smaller than the full range of the momentum (which is $`2\pi /a`$ on a lattice if the lattice spacing is $`a`$), but $`v_F\mathrm{\Lambda }`$ is much larger than the temperatures or frequencies of interest. If we include only these regions of momenta, then the second quantized Fermi field can be written in the approximate form
$$\psi (x,t)=\psi _R(x,t)e^{ik_Fx}+\psi _L(x,t)e^{ik_Fx},$$
(81)
where $`\psi _R`$ and $`\psi _L`$ vary slowly over spatial regions which are large compared to the distance scale $`1/\mathrm{\Lambda }`$. The momentum components of these slowly varying fields are related to those of $`\psi `$ as
$`\psi _{R,k}(t)`$ $`=`$ $`\psi _{k+k_F}(t),`$
$`\psi _{L,k}(t)`$ $`=`$ $`\psi _{kk_F}(t),`$ (82)
where $`\mathrm{\Lambda }k\mathrm{\Lambda }`$. These long-wavelength fields are the ones to which the technique of bosonization can be applied.
The definitions in Eqs. (81-82) tell us the forms of the various terms in a microscopic model and also tell us which of them survive in the long-wavelength limit. For instance, the density is given by
$`\rho `$ $`=`$ $`:\psi ^{}\psi :=:\psi _R^{}\psi _R+\psi _L^{}\psi _L+e^{i2k_Fx}\psi _R^{}\psi _L+e^{i2k_Fx}\psi _L^{}\psi _R:`$ (83)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{\pi }}}{\displaystyle \frac{\varphi }{x}}+{\displaystyle \frac{1}{2\pi \alpha }}[\eta _R^{}\eta _Le^{i(2\sqrt{\pi }\varphi 2k_Fx)}+\eta _L^{}\eta _Re^{i(2\sqrt{\pi }\varphi 2k_Fx)}].`$
The terms containing $`\mathrm{exp}(\pm i2k_Fx)`$ in Eq. (83) vary on a distance scale $`k_F^1`$ which is typically of the same order as the inverse particle density $`\rho ^1`$. These terms can therefore be ignored if we are only interested in the asymptotic behavior of correlation functions at distances much larger than $`k_F^1`$. In a lattice model, we have to be more careful about this argument since the lattice momentum only needs to be conserved modulo $`2\pi /a`$ in any process. However, since $`0<k_F<\pi /a`$ in general, and $`x/a`$ is an integer, we see that the last two terms in Eq. (83) vary on the scale of the lattice unit $`a`$; we can therefore ignore those terms if we are only interested in phenomena at distance scales which are much larger than $`a`$.
On the other hand, there are situations when a density term like $`\rho \mathrm{cos}(2k_Fx)`$ is generated in the model; for instance, this happens below a Peierls transition if the fermions are coupled to lattice phonons. We then find that the slowly varying terms in the continuum field theory are given by
$`\mathrm{cos}(2k_Fx)\rho `$ $`=`$ $`{\displaystyle \frac{1}{2}}[\psi _R^{}\psi _L+\psi _L^{}\psi _R]\mathrm{in}\mathrm{general},`$ (84)
$`=`$ $`\psi _R^{}\psi _L+\psi _L^{}\psi _R\mathrm{if}e^{i4k_Fx}=1.`$
The second possibility can arise in a lattice model if $`4k_Fa=2\pi `$, $`i.e.`$, at half-filling; we then call it a dimerized system. We will call the term on the right hand sides of Eq. (84) the mass operator. We will see below that for any value of $`K<2`$, this term produces a gap in the low-energy spectrum. This is called the dimerization gap if it occurs in a lattice system.
We should emphasize an important difference between models defined in the continuum and those defined on a lattice. In the continuum, $`\psi _R^2(x)=\psi _L^2(x)=0`$ due to the anticommutation relations. Therefore a term like $`\psi _R^2(x)\psi _L^2(x)`$ is equal to zero in the continuum. However such a term need not vanish on a lattice, if we take the two factors of $`\psi _R^{}`$ (or $`\psi _L`$) as coming from two neighboring sites separated by a distance $`a`$. In fact, this term is allowed by momentum conservation on a lattice if $`4k_Fa=2\pi `$, and it leads to umklapp scattering.
## 3 Correlation functions and dimensions of operators
We will now use bosonization to compute the correlation functions of some fermionic operators in the interacting theory discussed above. The power-law fall-offs of the correlation functions will tell us the dimensions of those operators.
The bosonic correlation function can be found from the commutation relations in Eq. (51), remembering that all normal-orderings have to be done with respect to the new vacuum $`|\stackrel{~}{0}`$. (Henceforth we will omit the tilde denoting the new vacuum, but we will continue to use the tilde for the new $`\varphi `$ fields). For instance,
$$0|T\stackrel{~}{\varphi }(x,t)\stackrel{~}{\varphi }^{}(0,0)|0=\frac{1}{4\pi }\mathrm{ln}\left[\left(\frac{2\pi }{L}\right)^2\left(x^2(vti\alpha \mathrm{sign}(t))^2\right)\right].$$
(85)
We can use the expressions in Eq. (57) and identities like Eq. (30) to obtain the correlation functions of various operators. For instance,
$`0|Te^{i2\sqrt{\pi }\beta \stackrel{~}{\varphi }_R(x,t)}e^{i2\sqrt{\pi }\beta \stackrel{~}{\varphi }_R(0,0)}|0`$ $``$ $`\left({\displaystyle \frac{\alpha }{vtxi\alpha \mathrm{sign}(t)}}\right)^{\beta ^2},`$
$`0|Te^{i2\sqrt{\pi }\beta \stackrel{~}{\varphi }_L(x,t)}e^{i2\sqrt{\pi }\beta \stackrel{~}{\varphi }_L(0,0)}|0`$ $``$ $`\left({\displaystyle \frac{\alpha }{vt+xi\alpha \mathrm{sign}(t)}}\right)^{\beta ^2},`$ (86)
in the limit $`L\mathrm{}`$; we will assume henceforth that this limit is taken in the calculation of all correlation functions. Consider now the positive-chirality fermion field; according to Eq. (57),
$$\psi _R=\frac{1}{\sqrt{2\pi \alpha }}\eta _Re^{i2\sqrt{\pi }\varphi _R},$$
(87)
where $`\varphi _R`$ is given in Eq. (70) in terms of $`\stackrel{~}{\varphi }_R`$ and $`\stackrel{~}{\varphi }_L`$. Hence its time-ordered correlation function takes the form
$$0|T\psi _R(x,t)\psi _R^{}(0,0)|0\frac{\alpha ^{(1K)^2/2K}}{2\pi (vtxi\alpha \mathrm{sign}(t))^{(1+K)^2/4K}(vt+xi\alpha \mathrm{sign}(t))^{(1K)^2/4K}}.$$
(88)
We see that the correlation function falls off at large space-time distances ($`i.e.`$, large compared to $`\alpha `$) with the power $`(1+K^2)/(2K)`$. This means that the scaling dimension of the operator $`\psi _R`$ or $`\psi _R^{}`$ is $`(1+K^2)/4K`$; this agrees with the familiar value of $`1/2`$ for non-interacting fermions.
If we set $`x=0`$ in Eq. (88), and Fourier transform over time, we find that the one-particle density of states (DOS) has a power-law form near zero frequency,
$$\stackrel{~}{n}(\omega )|\omega |^\beta ,$$
(89)
where
$$\beta =\frac{(1K)^2}{2K}.$$
(90)
The same result holds for the DOS of the negative-chirality fermions. We therefore see that for any non-zero interaction, either repulsive or attractive, the one-particle DOS vanishes as a power. (This result is not to be confused with the bosonic DOS which, from Eq. (68), is a constant near zero energy since the energy is linearly related to the momentum which has a constant density. That leads to a specific heat which is linear in the temperature at low temperatures). Alternatively, we may set $`t=0`$ in Eq. (88) and Fourier transform over space, with a factor of $`\mathrm{exp}(ik_Fx)`$ since the momentum of the right-chirality fermions is measured with respect to the Fermi momentum $`k_F`$. We then see that the momentum distribution function is continuous at $`k_F`$ with a power-law form,
$$n(k)=n(k_F)+\mathrm{constant}\mathrm{sign}(kk_F)|kk_F|^\beta ,$$
(91)
as we have sketched in Fig. 1 (b). These expressions for $`\stackrel{~}{n}(\omega )`$ and $`n(k)`$ are characteristic features of a Luttinger liquid.
Next let us compute the correlation function of an operator which is bilinear in the fermion fields, namely, the mass operator
$$M=\psi _R^{}\psi _L+\psi _L^{}\psi _R=\frac{1}{2\pi \alpha }\left[\eta _R^{}\eta _Le^{i2\sqrt{\pi }\varphi }+\eta _L^{}\eta _Re^{i2\sqrt{\pi }\varphi }\right].$$
(92)
Using the same technique as before, we find that
$$0|TM(x,t)M(0,0)|0\frac{\alpha ^{2(K1)}}{4\pi ^2((vti\alpha \mathrm{sign}(t))^2x^2)^K}.$$
(93)
This shows that the scaling dimension of the mass operator is $`K`$. For the non-interacting case $`K=1`$, we see that the addition of such a term to the Lagrangian density in Eq. (71) makes the Dirac fermion massive; this is why we have called it the mass operator. (For convenience, we will sometimes omit the Klein factors when writing fermionic operators in the bosonic language. We will of course need to restore those factors when calculating the correlation functions; clearly, correlation functions will vanish if the numbers of $`\eta _R`$ and $`\eta _R^{}`$ (or $`\eta _L`$ and $`\eta _L^{}`$) are not equal).
An important operator to consider is the density $`\rho `$. From Eqs. (83), (85) and (93), we see that the density-density equal-time correlation function is asymptotically given by
$$0|\rho (x,0)\rho (0,0)|0=\frac{K}{2\pi ^2x^2}+\mathrm{const}\frac{\mathrm{cos}(4k_Fx)}{x^{2K}}.$$
(94)
We should emphasize that this is only the asymptotic expression; the complete expression generally contains oscillatory terms like $`\mathrm{cos}(4nk_Fx)/x^{2n^2K}`$ for all positive integers $`n`$. However the form of the denominator shows that these terms decay rapidly with $`x`$ as $`n`$ increases.
In general, we can consider an operator of the form
$$O_{m,n}=e^{i2\sqrt{\pi }(m\varphi +n\theta )}.$$
(95)
(Such an operator can arise from a product of several $`\psi `$’s and $`\psi ^{}`$’s if we ignore the Klein factors; then Eq. (57) implies that $`m\pm n`$ must take integer values). We then find the following result for the two-point correlation function
$`0|TO_{m,n}(x,t)O_{m^{},n^{}}^{}(0,0)|0`$
$`\delta _{mm^{}}\delta _{nn^{}}{\displaystyle \frac{\alpha ^{2(m^2K+n^2/K)}}{(vtxi\alpha \mathrm{sign}(t))^{(m\sqrt{K}n/\sqrt{K})^2}(vt+xi\alpha \mathrm{sign}(t))^{(m\sqrt{K}+n/\sqrt{K})^2}}},`$ (96)
where we have taken the limit $`L\mathrm{}`$ as usual. (If $`L`$ had been kept finite, the correlation function in Eq. (96) would have been non-zero even if $`mm^{}`$ or $`nn^{}`$. This may seem surprising since the global phase invariance in Eqs. (79 \- 80) should lead to the Kronecker $`\delta `$’s in Eq. (96) even for finite values of $`L`$. The resolution of this puzzle is that we need to include the appropriate Klein factors in the definition in (95) to show that the correlation function of a product of fermionic operators is zero if it is not phase invariant). We conclude that the scaling dimension of $`O_{m,n}`$ is given by
$$d_O=m^2K+\frac{n^2}{K}.$$
(97)
The appearance of the cut-off $`\alpha `$ in the expressions for the various correlation functions may seem bothersome. This may be eliminated by redefining the operators $`O_{m,n}`$ in Eq. (95) by multiplying them with appropriate $`K`$-dependent powers of $`\alpha `$; then the two-point correlation function has a well-defined limit as $`\alpha 0`$. The important point to note is that all the correlation functions fall off as power-laws asymptotically, and that the exponents give the scaling dimensions of those operators. The significance of the scaling dimension will be discussed in the next section.
For certain applications of bosonization, it is useful to know the forms of the correlation functions in imaginary time. From the various expressions above, it is clear that if $`x`$ is held fixed at some non-zero value, then the poles in the complex $`t`$ plane are either in the first or in the third quadrant. We may therefore rotate $`t`$ by $`\pi /2`$ without crossing any poles. After doing this, we write $`t=i\tau `$ where $`\tau `$ is a real variable. Eq. (96) then takes the form
$$O_{m,n}(x,t)O_{m^{},n^{}}^{}(0,0)\delta _{mm^{}}\delta _{nn^{}}e^{i4mn\zeta }(\frac{\alpha ^2}{x^2+v^2\tau ^2})^{m^2K+n^2/K},$$
(98)
where $`\zeta =\mathrm{tan}^1(v\tau /x)`$, and we have dropped the $`\alpha `$-dependent terms in the denominator since there are no longer any poles for non-zero values of $`x`$.
## 4 Renormalization group analysis of perturbed models
We will now study the effects of some perturbations on the low-energy properties of Luttinger liquids. A standard way to do this is to use the renormalization group (RG) idea. Suppose that we are given an action at a microscopic length scale which may be a lattice spacing $`a`$; the action contains some small perturbations proportional to certain dimensionless parameters $`\lambda _i`$, such that for $`\lambda _i=0`$, we have a gapless system with an infinite correlation length $`\xi `$, $`i.e.`$, all correlations fall off as power laws. Then the RG procedure typically consists of the following steps.
(i) First, a small range of high momentum modes of the various fields are integrated out. Specifically, we will assume that the momenta lie in the range $`[\mathrm{\Lambda },\mathrm{\Lambda }]`$ while the frequencies go all the way from $`\mathrm{}`$ to $`\mathrm{}`$. Then we will integrate out the modes with momenta lying in the two intervals $`[\mathrm{\Lambda },\mathrm{\Lambda }/s]`$ and $`[\mathrm{\Lambda }/s,\mathrm{\Lambda }]`$ and with all frequencies from $`\mathrm{}`$ to $`\mathrm{}`$. Here $`s=e^{dl}`$ where $`dl`$ is a small positive number. The asymmetry between the momentum and frequency integrations is necessary to ensure that the action remains local in time at all stages. (Note that we are using sharp momentum cut-offs in this section, whereas we used a smooth momentum cut-off with the parameter $`\alpha `$ in the previous sections).
(ii) Secondly, the space-time coordinates, the fields and the various parameters are rescaled by appropriate powers of $`s`$ so that the new action looks exactly like the old action. This new action is effectively at a larger length scale equal to $`ae^{dl}`$. Clearly, the changes in the parameters $`\lambda _i`$ must be proportional to the small number $`dl`$. Since we are going to repeat the process of integrating out high momentum modes, we introduce the idea of an effective length scale $`a(l)=ae^l`$; we also define length scale dependent parameters $`\lambda _i(l)`$, where $`\lambda _i(0)`$ denote the values of $`\lambda _i`$ in the original action. We then define the $`\beta `$-functions
$$\beta (\lambda _i)=\frac{d\lambda _i(l)}{dl}.$$
(99)
These are functions of all the $`\lambda _i(l)`$’s so that we get a set of coupled non-linear equations in general. In principle, the $`\beta `$-functions are given by infinite power series in the $`\lambda _i`$, but in practice, we can only obtain the first few terms depending on the number of loops of the various Feynman diagrams that we can compute. The RG analysis is therefore usually limited to small values of $`\lambda _i`$.
(iii) Finally, we integrate the RG equations, $`i.e.`$, the differential equations described by the $`\beta `$-functions, in order to obtain the functions $`\lambda _i(l)`$. For simplicity, let us consider the case of a single perturbation with a coefficient $`\lambda `$. Then one of three things can happen as $`l`$ increases from $`0`$. Either $`\lambda (l)`$ goes to zero in which case we recover the unperturbed theory at long distances; or $`\lambda `$ does not change with $`l`$; or $`\lambda (l)`$ grows with $`l`$ till its value becomes of order $`1`$. In the last case, the RG equation cannot be trusted beyond that length scale since the $`\beta `$-functions are generally only known up to some low order in the $`\lambda `$’s. All that we can say is that beyond the length scale $`ae^l`$ where $`\lambda (l)`$ becomes of order $`1`$, a completely new kind of action is likely to be required to describe the system. Large perturbations of a gapless system often (but not always) correspond to a gapped system whose correlation length $`\xi `$ (which governs the exponential decay of various correlation functions) is of the same order as that length scale $`ae^l`$. Thus, although the blowing up of a parameter $`\lambda `$ at some scale does not tell us what the new action must be beyond that scale, it can give us an idea of the correlation length of that new theory. This is the main use that is made of RG equations. To complete the picture and find the new theory beyond the scale $`\xi `$, one usually has to do some other kind of analysis.
Let us now examine in a little more detail the various kinds of RG equations which can arise at low orders. Suppose that to first order, the RG equation for a single perturbing term is given by
$$\frac{d\lambda }{dl}=b_1\lambda ,$$
(100)
where $`b_1`$ is some constant. If $`b_1<0`$, any non-zero value of $`\lambda `$ at $`l=0`$ flows to $`0`$ as $`l`$ increases. Such a perturbation is called irrelevant. If $`b_1>0`$, it is called a relevant perturbation. A small perturbation then grows exponentially with $`l`$ and reaches a number of order $`1`$ at a distance scale given by $`e^{b_1l}1/\lambda (0)`$. In many models, this gives an estimate of the correlation length $`\xi `$ and of the energy gap $`\mathrm{\Delta }E`$ of the system, namely,
$`\xi `$ $`=`$ $`ae^l={\displaystyle \frac{a}{\lambda (0)^{1/b_1}}},`$
$`\mathrm{so}\mathrm{\Delta }E`$ $`=`$ $`{\displaystyle \frac{v}{\xi }}={\displaystyle \frac{v\lambda (0)^{1/b_1}}{a}}.`$ (101)
Finally, if $`b_1=0`$, the perturbation is called marginal. One then has to go to second order in $`\lambda `$. If the RG equation takes the form
$$\frac{d\lambda }{dl}=b_2\lambda ^2,$$
(102)
then a small perturbation of one particular sign flows to zero and is called marginally irrelevant, while a small perturbation of the opposite sign grows and is called marginally relevant. For instance, suppose that $`b_2>0`$. Then the above equation gives
$$\lambda (l)=\frac{\lambda (0)}{1b_2\lambda (0)l}.$$
(103)
If we start with a negative value of $`\lambda (0)`$, $`\lambda (l)`$ flows to $`0`$. For large $`l`$, $`\lambda (l)`$ goes to zero logarithmically in the distance scale as $`1/(b_2l)`$ independently of the starting value. (It turns out that this produces logarithmic corrections to the power-law fall-offs of the correlation functions at large distances and the various excitation energies ). On the other hand, if we start with a small positive value of $`\lambda (0)`$, then $`\lambda (l)`$ grows and becomes of order $`1`$ at a distance scale which we identify with a correlation length
$$\xi =ae^{1/(b_2\lambda (0))}.$$
(104)
The corresponding energy gap $`\mathrm{\Delta }E=v/\xi `$ is extremely small for small values of $`\lambda (0)`$; it may be very hard to distinguish this kind of a system from a gapless system by numerical studies. This is in sharp contrast to the situation with a relevant perturbation where the gap scales as a power of $`\lambda (0)`$.
There is a simple relation between the scaling dimension of an operator $`O`$ (assumed to be hermitian for simplicity), and the first-order coefficient $`b_1`$ in its $`\beta `$-function. We recall that the scaling dimension $`d_O`$ is defined as half of the exponent appearing in the two-point correlation function at large distances, namely,
$$O(x,0)O(0,0)=|x|^{2d_𝒪}.$$
(105)
It is convenient to define the normalization of $`O`$ in such a way that the right hand side of Eq. (105) has a prefactor equal to $`1`$. Consequently, $`O`$ has the engineering dimensions of $`a^{d_𝒪}`$. Let us now add a perturbation to the Hamiltonian (or to the Lagrangian with a negative sign) of the form
$$\delta H=\lambda a^{d_𝒪2}v𝑑xO,$$
(106)
where the factors of $`a`$ and $`v`$ (the velocity of the unperturbed Luttinger liquid) are put in to make $`\lambda `$ dimensionless; note that $`v/a`$ has the dimensions of energy. Then the first-order RG equation for $`\lambda `$ must take the form given in Eq. (100) with
$$b_1=2d_O.$$
(107)
This important statement will be proved below for the class of operators $`O_{m,n}`$ introduced in Eq. (95). If $`d_O=2`$, the perturbation is marginal and we have to proceed to Eq. (102). It turns out that $`b_2`$ can be obtained from a three-point correlation function, but we will not pursue that here .
It will not come as a surprise that the RG equations for interacting quantum systems in one dimension can often be derived in two different ways, namely, using the fermionic theory or the bosonic one. Although both the derivations are limited in practice to small values of the perturbations $`\lambda _i`$, we will see that the bosonic derivation is superior because it can handle the interactions in Eq. (63) exactly. In the fermionic derivation, we have to assume that not only the $`\lambda `$’s but also the interaction parameters $`g_2`$ and $`g_4`$ are small. We will now discuss some simple examples of $`\beta `$-function calculations to first order in the two kinds of theories.
As a particularly simple exercise, consider a non-interacting massive Dirac theory, where the mass term is to be treated as a perturbation. We define the Fourier components of $`\psi _\nu `$ as
$`\psi _R(x,t)`$ $`=`$ $`{\displaystyle _\mathrm{\Lambda }^\mathrm{\Lambda }}{\displaystyle \frac{dk}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\omega }{2\pi }}e^{i(kx\omega t)}\psi _R(k,\omega ),`$
$`\psi _L(x,t)`$ $`=`$ $`{\displaystyle _\mathrm{\Lambda }^\mathrm{\Lambda }}{\displaystyle \frac{dk}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\omega }{2\pi }}e^{i(kx+\omega t)}\psi _L(k,\omega ).`$ (108)
Then the action takes the form
$`S[\psi _\nu ,\psi _\nu ^{}]={\displaystyle _\mathrm{\Lambda }^\mathrm{\Lambda }}{\displaystyle \frac{dk}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\omega }{2\pi }}[`$ $`\psi _R^{}(k,\omega )(\omega vk)\psi _R(k,\omega )+\psi _L^{}(k,\omega )(\omega vk)\psi _L(k,\omega )`$ (109)
$`\mu (\psi _R^{}(k,\omega )\psi _L(k,\omega )+\psi _L^{}(k,\omega )\psi _R(k,\omega ))].`$
Since $`\mu `$ has the dimensions of energy, the dimensionless parameter must be taken to be
$$\lambda =\frac{a\mu }{v}.$$
(110)
(The value of $`a`$ is completely arbitrary here and it will not appear in any physical quantity as we will see). We consider the partition function in the functional integral representation,
$$Z=𝒟\psi _\nu 𝒟\psi _\nu ^{}e^{iS}.$$
(111)
We integrate out the modes in the momentum and frequency ranges specified in step (i) of the RG procedure outlined above. Since Eq. (109) describes non-interacting fermions, the mode integration produces an action which looks exactly the same, except that the momentum integrations go from $`\mathrm{\Lambda }/s`$ to $`\mathrm{\Lambda }/s`$. To restore this to the original range of $`[\mathrm{\Lambda },\mathrm{\Lambda }]`$, we define the new (primed) quantities
$`k^{}`$ $`=`$ $`sk,`$
$`\omega ^{}`$ $`=`$ $`s\omega ,`$
$`\psi _\nu ^{}(k^{},\omega ^{})`$ $`=`$ $`s^{3/2}\psi _\nu (k,\omega ),`$
$`\lambda ^{}`$ $`=`$ $`s\lambda .`$ (112)
The resultant action in terms of the new variables and fields looks exactly the same as the original action in terms of the old variables. Note that we had to rescale the mass parameter also in order to achieve this. Since $`s=e^{dl}`$, we obtain the RG equation
$$\frac{d\lambda (l)}{dl}=\lambda (l).$$
(113)
Clearly, this describes a relevant perturbation, and $`\lambda (l)`$ grows to $`1`$ at a length scale
$$\xi =ae^l=\frac{v}{\mu }.$$
(114)
The energy gap is $`\mathrm{\Delta }E=v/\xi =\mu `$ as expected.
Now let us add density-density interactions as in Eq. (63) to the above massive theory. The question is the following: do $`\xi `$ and $`\mathrm{\Delta }E`$ scale in the same way with $`\mu `$ as they do in the non-interacting theory? Clearly, it is not easy to answer this in the fermionic language since the interactions themselves are not easy to handle in that language, and the mass perturbation is an additional complication. But bosonization comes to our rescue here since the bosonic theory remains quadratic even after including the four-fermi interactions; hence the mass perturbation is the only thing that needs to be studied.
Let us consider a more general perturbing operator of the form
$$O=O_{m,0}+O_{m,0}^{},$$
(115)
where $`O_{m,n}`$ is defined in Eq. (95); the reason for setting $`n=0`$ will be explained later. From Eq. (97), the scaling dimension of $`O`$ is given by $`d_O=m^2K`$; note that this contains the effects of the four-fermion interaction in a non-trivial way through the parameter $`K`$. In the bosonic language, the perturbed action has the sine-Gordon form,
$$S[\stackrel{~}{\varphi }]=𝑑x𝑑t\left[\frac{1}{2v}(_t\stackrel{~}{\varphi })^2\frac{v}{2}(_x\stackrel{~}{\varphi })^2\frac{v\lambda }{a^2}\mathrm{cos}(2m\sqrt{\pi K}\stackrel{~}{\varphi })\right],$$
(116)
where we have changed variables from $`\varphi `$ to $`\stackrel{~}{\varphi }`$ using Eq. (70). We now have to apply the RG procedure to this action. We introduce the Fourier components of $`\stackrel{~}{\varphi }`$ as
$$\stackrel{~}{\varphi }(x,t)=_\mathrm{\Lambda }^\mathrm{\Lambda }\frac{dk}{2\pi }_{\mathrm{}}^{\mathrm{}}\frac{d\omega }{2\pi }e^{i(kx\omega t)}\stackrel{~}{\varphi }(k,\omega ).$$
(117)
(In principle, the momentum cut-offs for fermion and boson fields need not be equal, but we will use the same symbol $`\mathrm{\Lambda }`$ for convenience). Next we consider the partition function
$$Z=𝒟\stackrel{~}{\varphi }e^{iS},$$
(118)
and expand $`e^{iS}`$ in powers of $`\lambda `$ to obtain an infinite series. Let us write the field $`\stackrel{~}{\varphi }`$ as the sum
$$\stackrel{~}{\varphi }=\stackrel{~}{\varphi }_<+\stackrel{~}{\varphi }_>,$$
(119)
where both $`\stackrel{~}{\varphi }_<`$ and $`\stackrel{~}{\varphi }_>`$ contain all frequencies, but $`\stackrel{~}{\varphi }_<`$ only contains momenta lying in the range $`[\mathrm{\Lambda }/s,\mathrm{\Lambda }/s]`$, whereas $`\stackrel{~}{\varphi }_>`$ only contains momenta lying in the ranges $`[\mathrm{\Lambda },\mathrm{\Lambda }/s]`$ and $`[\mathrm{\Lambda }/s,\mathrm{\Lambda }]`$. Following step (i) of the RG procedure, we have to perform the functional integration over $`\stackrel{~}{\varphi }_>`$, and then re-exponentiate the infinite series to obtain the new action in terms of $`\stackrel{~}{\varphi }_<`$. We will do this calculation only to first order in $`\lambda `$. This is not difficult since $`e^{\pm i2m\sqrt{\pi K}\stackrel{~}{\varphi }}`$ can be written as the product of exponentials of $`\stackrel{~}{\varphi }_<`$ and $`\stackrel{~}{\varphi }_>`$, while the quadratic part of the action decouples as $`S_0[\stackrel{~}{\varphi }]=S_0[\stackrel{~}{\varphi }_<]+S_0[\stackrel{~}{\varphi }_>]`$. Let us denote the expectation value of a functional $`F[\stackrel{~}{\varphi }_>]`$ as
$$F[\stackrel{~}{\varphi }_>]=𝒟\stackrel{~}{\varphi }_>e^{iS_0[\stackrel{~}{\varphi }_>]}F[\stackrel{~}{\varphi }_>].$$
(120)
Now we have to compute
$$e^{\pm i2m\sqrt{\pi K}\stackrel{~}{\varphi }_>(x,t)}.$$
(121)
By translation invariance, the value of this is independent of the coordinates $`(x,t)`$, so we can evaluate it at the point $`(0,0)`$. We then use the fact that $`\stackrel{~}{\varphi }_>^n(0,0)=0`$ if $`n`$ is odd, while
$$\stackrel{~}{\varphi }_>^n(0,0)=(n1)(n3)\mathrm{}1\stackrel{~}{\varphi }_>^2(0,0)^{n/2}$$
(122)
if $`n`$ is even. Thus the expectation value in (121) is given by
$`e^{\pm i2m\sqrt{\pi K}\stackrel{~}{\varphi }_>(0,0)}`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}(\pm i2m\sqrt{\pi }K)^n\stackrel{~}{\varphi }_>^n(0,0)`$ (123)
$`=`$ $`e^{2m^2\pi K\stackrel{~}{\varphi }_>^2(0,0)}.`$
Now we use the fact that
$$\stackrel{~}{\varphi }_>^2(0,0)=2_{\mathrm{\Lambda }/s}^\mathrm{\Lambda }\frac{dk}{2\pi }_{\mathrm{}}^{\mathrm{}}\frac{d\omega }{2\pi }\frac{i}{\omega ^2/vvk^2+iϵ}=\frac{\mathrm{ln}s}{2\pi }$$
(124)
to show that the left hand side of Eq. (123) is equal to $`s^{m^2K}`$. Putting everything together, we find the new action to be
$$S[\stackrel{~}{\varphi }_<]=𝑑x𝑑t\left[\frac{1}{2v}(_t\stackrel{~}{\varphi }_<)^2\frac{v}{2}(_x\stackrel{~}{\varphi }_<)^2\frac{v\lambda s^{m^2K}}{a^2}\mathrm{cos}(2m\sqrt{\pi K}\stackrel{~}{\varphi }_<)\right],$$
(125)
where the momentum integrals only go from $`\mathrm{\Lambda }/s`$ to $`\mathrm{\Lambda }/s`$. To restore the range of the momentum to $`[\mathrm{\Lambda },\mathrm{\Lambda }]`$ and to recover the form of the action in Eq. (116), we have to define
$`k^{}`$ $`=`$ $`sk,\mathrm{and}x^{}=s^1x,`$
$`\omega ^{}`$ $`=`$ $`s\omega ,\mathrm{and}t^{}=s^1t,`$
$`\stackrel{~}{\varphi }^{}(k^{},\omega ^{})`$ $`=`$ $`\stackrel{~}{\varphi }_<(k,\omega ),`$
$`\lambda ^{}`$ $`=`$ $`s^{2m^2K}\lambda ,`$ (126)
and write the action in terms of primed variables. Since $`s=e^{dl}`$, we see that $`d\lambda =\lambda ^{}\lambda `$ satisfies the RG equation
$$\frac{d\lambda }{dl}=(2m^2K)\lambda .$$
(127)
This proves the relation between the first-order $`\beta `$-function coefficient $`b_1`$ and the scaling dimension $`d_O`$. Note that the $`\beta `$-functions of the parameters $`v`$ and $`K`$ remain zero up to this order in the perturbation. However they do get a contribution to second order in $`\lambda `$ as shown in Ref. 3.
The mass perturbation corresponds to the special case of Eq. (115) with $`m=1`$. We now see that it is marginal for $`K=2`$ and is relevant if $`K<2`$. In the latter case, $`\lambda (l)`$ grows till we reach a length scale $`\xi =a/\lambda (0)^{1/(2K)}`$ where the length scale of the coefficient of the cosine term in the Lagrangian becomes of the same order as $`a`$; that is the appropriate point to stop the RG flow of $`\lambda `$. The expression for $`\xi `$ implies that the energy gap of the system is given by
$$\mathrm{\Delta }E=\frac{v}{a}\lambda (0)^{1/(2K)}.$$
(128)
Thus the effect of the renormalization is to produce a sine-Gordon theory with the Lagrangian density
$$=\frac{1}{2v}(_t\stackrel{~}{\varphi })^2\frac{v}{2}(_x\stackrel{~}{\varphi })^2\mathrm{const}\frac{(\mathrm{\Delta }E)^2}{v}\mathrm{cos}(2\sqrt{\pi K}\stackrel{~}{\varphi }),$$
(129)
where $`x`$ and $`t`$ in this expression denote the original coordinates, and it is understood that this Lagrangian is not to be renormalized any further. This theory is exactly solvable and its spectrum is known in detail . It has both bosonic and fermionic (soliton) excitations, and both of them have energy gaps of the order of $`\mathrm{\Delta }E`$ given in Eq. (128).
Finally, let us briefly consider some other relevant and marginal perturbations that can appear in a system which, at the microscopic model, involves fermions on a lattice. If the model has the global phase invariance discussed in Eqs. (79 \- 80), then the operators $`O_{m,n}`$ appearing in the bosonized theory must necessarily have $`n=0`$. The scaling dimension is then $`d_O=m^2K`$. Since $`m1`$, there is only a finite number of relevant operators possible depending on the value of $`K`$. For $`K>2`$, there are no relevant operators at all. For $`1/2<K<2`$, the mass operator is the only relevant term, and so on.
Turning to the possible marginal operators, we see that the umklapp operator $`O_{2,0}=\psi _R^2\psi _L^2`$ is marginal for $`K=1/2`$. This is a particularly important case to consider because a Luttinger liquid at $`K=1/2`$ is known to have a global $`SU(2)`$ symmetry; it therefore describes a large number of gapless systems involving spins. From conformal field theory, the value of $`b_2`$ in the RG equation Eq. (102) for the umklapp operator $`O`$ is exactly known to be $`4\pi /\sqrt{3}`$ for the normalization given in Eq. (105). The coefficient of $`O`$ in the Hamiltonian, namely $`\lambda `$, depends on the microscopic parameters of the model. In general, a system will have a non-zero value of $`\lambda `$. As discussed above, for one sign of $`\lambda `$, the system remains gapless but with logarithmic corrections to various physical quantities; for instance, a $`1/\mathrm{ln}T`$ term appears in the magnetic susceptibility of a spin system at low temperatures. For the other sign of $`\lambda `$, the system spontaneously dimerizes producing a finite correlation length and an energy gap; this leads to an exponentially vanishing susceptibility at low temperatures.
## 5 Applications of Bosonization
We will now study various applications of the method of bosonization. The method, as you have learned, can only be applied in one dimension, so we restrict ourselves to one-dimensional models. As you have also seen, the main advantage of the method of bosonization is that many interacting fermion theories can often be recast (within some approximations) as non-interacting boson theories. This enables the explicit calculation of correlation functions. This is an advantage, even in Bethe ansatz solvable one-dimensional models, because it is often not possible to compute correlation functions using the Bethe ansatz.
We will concentrate on the applications of the bosonization technique in the following problems - (i) the quantum antiferromagnetic spin 1/2 chain, (ii) the Hubbard model in one dimension, (iii) transport in clean quantum wires and (iv) transport through isolated impurities. Since the physics of each of these applications is a huge subject by itself, here we will only concentrate on explaining the model and the quantities that we can obtain through the use of bosonization, rather than go into details of its phenomenology.
## 6 Quantum antiferromagnetic spin 1/2 chain
The model
The first problem that we shall study is the model of a spin 1/2 antiferromagnetic chain. We are picking this model, since you have already learned a lot about the model from the course on quantum spin chains and spin ladders. Here, we will restrict ourselves to just the study of the spin 1/2 anisotropic Heisenberg model with the Hamiltonian given by
$$H=\underset{i=1}{\overset{N}{}}[\frac{J}{2}(S_i^+S_{i+1}^{}+S_i^{}S_{i+1}^+)+J_zS_i^zS_{i+1}^z],$$
(130)
where the interactions are only between nearest neighbor spins, and $`J>0`$. $`S_i^+=S_i^x+iS_i^y`$ and $`S_i^{}=S_i^xiS_i^y`$ are the spin raising and lowering operators. Although this model can be exactly solved using the Bethe ansatz and one has the explicit result that the model is gapless for $`JJ_zJ`$ and gapped for $`J_z>J`$, (there is a phase transition exactly at the isotropic point $`J_z=J`$), it is not easy to compute explicit correlation functions in that approach. Hence, it is more profitable to use field theory methods.
Symmetries of the model
Note that this spin model has a global $`U(1)`$ invariance, which is rotations about the $`S^z`$ axis. Precisely when $`J_z=J`$, the $`U(1)`$ invariance is enhanced to an $`SU(2)`$ invariance, because at this point the model can simply be written as $`H=J_i𝐒_𝐢𝐒_{𝐢+\mathrm{𝟏}}`$. The model also has discrete symmetries under $`S^xS^x`$, $`S^yS^y`$, and under $`S^zS^z`$. Note also that one can change the sign of the $`XY`$ part of the Hamiltonian by making a rotation by $`\pi `$ about the $`S^z`$ axis on alternate sites, without affecting the $`J_z`$ term, although this is not an extra symmetry of the model.
Aside on non-linear sigma models
Even using field theory methods, there are two distinct approaches to the problem. In the large-$`S`$ limit, there exists a semiclassical field theory approach to this model, which leads to an $`O(3)`$ non-linear sigma model ($`NL\sigma M`$), with integer and half-integer spins being distinguished by the absence or presence of a Hopf term in the action. In this approach, it is easy to see that integer spin models have a gap in the spectrum. However, it is less easy to study the effect of the Hopf term and show that 1/2-integer spin models are gapless. In fact, in this case, it was the spin model which gave information about field theories with the Hopf term!
Jordan-Wigner transformation
For spin 1/2 models, it is possible to fermionize and then bosonize the spin model and study its spectrum. That is the approach we will follow in the rest of this lecture. First, we will try to convince you that it is possible to rewrite the spin model in terms of spinless fermions. The spin 1/2 model has two states possible at every site - spin $``$ or spin $``$. Hence,it can be mapped to another two state model which we can construct in terms of fermions. We shall assume that an $``$ spin or $``$ can be denoted by the presence or absence of a fermion at that site. Since no more than a single spinless fermion can sit at a site, the degrees of freedom in both the models are the same. This mapping is implemented by introducing a fermion annihilation operator $`\psi _i`$ at each site and writing the spin at the site as
$`S_i^z`$ $`=`$ $`\psi _i^{}\psi _i1/2=n_i1/2`$
$`S_i^{}`$ $`=`$ $`(1)^i\psi _ie^{i\pi _jn_j},`$ (131)
where the sum runs from one boundary of the chain up to the $`(i1)^{\mathrm{th}}`$ site and $`S_i^+`$ is the hermitian conjugate of $`S_i^{}`$. So an $``$-spin is denoted by $`n_i=1`$ and the $``$-spin by $`n_i=0`$ at the site $`i`$. One might have naively guessed that the spin-lowering operator should be expressed by $`\psi _i`$ which denotes annihilation of a fermion (with the spin raising operator being given by the hermitian conjugate). One can explicitly check that this gives the correct commutation relations of the spin operators at a site because $`[S_i^+,S_i^{}]=2S_i^z`$ just reproduces the correct anticommutation relations for the fermions $`\{\psi _i,\psi _i^{}\}=1`$. The extra string factor has to be added in order to correct for different site statistics - the fermions at different sites anticommute, whereas the spin operators commute. In fact, it is instructive to check explicitly that the string operator changes the commutation relation on different sites.
(H.W. Exercise 1. Check the above explicitly).
The Hamiltonian
Now, we rewrite the spin model in terms of the fermions. We find that
$`H`$ $`=`$ $`{\displaystyle \frac{J}{2}}{\displaystyle \underset{i}{}}[\psi _i^{}e^{i\pi n_i}\psi _{i+1}+h.c.]+J_z{\displaystyle \underset{i}{}}[(n_i1/2)(n_{i+1}1/2)]`$ (132)
$`=`$ $`{\displaystyle \frac{J}{2}}{\displaystyle \underset{i}{}}[\psi _i^{}\psi _{i+1}+h.c.]+J_z{\displaystyle \underset{i}{}}[(n_i1/2)(n_{i+1}1/2)].`$
The point to notice is that the string operator has cancelled out in the nearest neighbor interaction, except for a phase term, which also can be explicitly shown to be just $`1`$ because $`e^{i\pi n_i}`$ precedes a creation operator $`\psi _i^{}`$ which can only act if $`n_i=0`$. The spin-flip terms are like the hopping terms in the fermion Hamiltonian and give rise to motion of fermions whereas the $`S^z`$-$`S^z`$ interaction term leads to a four fermion interaction between fermions on adjacent sites (the analog of the on-site Hubbard interaction for spinless fermions). So for non-zero $`J_z`$, the fermionic model is non-trivial. There exists a competition between the hopping term or kinetic energy term, which gains in energy when the electrons are free to hop from site to site, and the potential energy which costs $`J_z`$ if there are electrons present on adjacent sites. So naively, for large $`J_z`$, one expects the potential energy to win and electrons to be localized on non-adjacent sites, and for small $`J_z`$, one expects the kinetic energy to win and to have delocalized fermions. Let us see whether this expectation is true and how it comes about.
Set $`J_z=0`$
To make the problem simpler, we first consider the case where $`J_z=0`$ or where there are no interaction terms. Then this is just the model of free spinless fermions. By Fourier transforming the fermions, - $`\psi _j=_k\psi _ke^{ikja}/\sqrt{N}`$, ( a is the lattice spacing) where the $`k`$ sum is over momentum values in the first Brillouin zone, - we find that the Hamiltonian is given by
$$H=J\underset{k}{}\mathrm{cos}(ka)\psi _k^{}\psi _k.$$
(133)
(H.W. Exercise 2. Obtain the above Hamiltonian explicitly).
The discrete symmetry of the model under $`S_i^{}S_i^+`$ and $`S_i^zS_i^z`$ implies a particle-hole symmetry $`\psi _i\psi _i^{}`$ in the fermion language. Thus, the ground state has to have total spin $`M_iS_i^z=0`$ or equivalently in the fermionic language, the ground state is precisely half-filled. This symmetry can be broken by the addition of a magnetic field term that couples linearly to $`S^z`$. In the fermionic language, this is equivalent to adding a chemical potential term (which couples to $`n_i`$ which is the $`S^z`$ term) in which case, the ground state no longer has $`M=0`$ and the fermion model is no longer half-filled. Thus, for $`M=0`$, the band is precisely half-filled and the Fermi surface ($`E=0`$) occurs exactly at $`ka=\pm \pi /2k_Fa`$ (because the density of states is symmetric about $`E=0`$. ) Low energy excitations are particle-hole excitations about the Fermi surface, which can occur either at a single Fermi point $`i.e.`$, $`k0`$ modes, or across Fermi points, which are the $`k2k_F=\pi /a`$ modes.
Effective field theory
The next step is to write down an effective field theory for the low energy modes. Now comes the approximation. Let us make the assumption that it is only the modes near the Fermi surface (or here, the Fermi points), which are relevant at low energies. Hence, we are only interested in $`ka`$ values near $`ka=\pm \pi /2`$ and we may approximate the dispersion relation around the Fermi points to be linear - $`i.e.`$, $`\mathrm{cos}(ka)=\mathrm{cos}(\pm k_Fa+k^{}a)=\mathrm{cos}(\pm \pi /2+k^{}a)=\mathrm{sin}(k^{}a)=(k^{}a)`$. We introduce the labels left and right to denote fermion modes near $`ka=\pi /2`$ and $`ka=\pi /2`$ respectively and henceforth drop the primes on the momenta and assume that they are always measured from the Fermi points; as before, we take $`k`$ as increasing towards the right near the right Fermi point, and increasing towards the left near the left Fermi point. If we want to solve the problem without any approximations, we have to allow for excitations about the Fermi points with arbitrary $`k`$. The approximation that is made is that we only allow small values of $`k`$ compared to $`k_F`$. This is why the excitations around the left and right Fermi points can be thought of as independent excitations. In this approximation, the Hamiltonian breaks up into
$$H=Ja\underset{k}{}k(\psi _{R,k}^{}\psi _{R,k}+\psi _{L,k}^{}\psi _{L,k}).$$
(134)
(Note that we have incorporated the change in sign mentioned below Eq. (133)). The fermions around the Fermi points are Dirac fermions since we have linearized the dispersion. These fields do not contain any high momentum modes. In real space, the original non-relativistic fermion field, which has high energy modes (rapidly oscillating factors), has been split up as exponential prefactors times smoothly varying fields -
$`\psi _j`$ $``$ $`e^{ik_Fja}{\displaystyle _{k_Fa\mathrm{\Lambda }}^{k_Fa+\mathrm{\Lambda }}}{\displaystyle \frac{d(ka)}{2\pi }}e^{ikja}\psi _k+e^{ik_Fja}{\displaystyle _{k_Fa\mathrm{\Lambda }}^{k_Fa+\mathrm{\Lambda }}}{\displaystyle \frac{d(ka)}{2\pi }}e^{ikja}\psi _k`$ (135)
$``$ $`e^{ik_Fja}\psi _{Lj}+e^{ik_Fja}\psi _{Rj}.`$
We assume that the $`\mathrm{\Lambda }<<k_Fa`$ and that it is sufficient to keep just these modes, if we are interested in physics at length scales much greater than $`1/\mathrm{\Lambda }`$, which is of course much greater than the lattice spacing. (The real physical cutoff is the lattice length or in momentum space, the Fermi momentum. As a low energy approximation, we are introducing the larger length cutoff $`1/\mathrm{\Lambda }`$ or the smaller momentum cutoff $`\mathrm{\Lambda }`$). For both $`R`$ and $`L`$ fermions, states with $`k>0`$ are empty and correspond to electron operators ($`c_k`$), while states with $`k<0`$ are filled and correspond to hole operators ($`d_k^{}`$). (See Fig. 2). In terms of these operators, the Hamiltonian can be rewritten as
$$H=Ja\underset{k>0}{}k(c_{L,k}^{}c_{L,k}+d_{L,k}^{}d_{L,k}+c_{R,k}^{}c_{R,k}+d_{R,k}^{}d_{R,k}).$$
(136)
We now introduce continuum fermion fields made up of particle (electron) and anti-particle (hole) operators at the left and right Fermi points as
$`\psi _R(x,t)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{Na}}}{\displaystyle \underset{k>0}{}}[c_{R,k}e^{ik(vtx)}+d_{R,k}^{}e^{ik(vtx)}]`$
$`\psi _L(x,t)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{Na}}}{\displaystyle \underset{k>0}{}}[c_{L,k}e^{ik(vt+x)}+d_{L,k}^{}e^{ik(vt+x)}],`$ (137)
where $`v=Ja`$ is defined to be the velocity. Note that the factor of $`1/\sqrt{a}`$ is needed to relate continuum fermions to lattice fermions. (The factor of $`\sqrt{a}`$ is needed to get the dimensions right. The lattice fermions satisfy $`\{\psi _i,\psi _j^{}\}=\delta _{ij}`$, whereas continuum fermions satisfy $`\{\psi (x),\psi ^{}(y)\}=\delta (xy)`$ where the Dirac $`\delta `$-function has the dimension of 1/length. Also $`Na=L`$ gives the conventional box normalization of the continuum fermions). Note also that the standard inclusion of the $`e^{ikvt}`$ for the particle fields and the $`e^{ikvt}`$ for the anti-particle fields, show that the right-movers are a function only of $`x_R=xvt`$ and the left-movers are a function only of $`x_L=x+vt`$. This observation will come in useful when we compute correlation functions. In general, we only need to compute equal time correlation functions. The time-dependent correlations are then obtained by replacing $`x`$ by $`x_R`$ for right-movers and by $`x_L`$ for left-movers.
In terms of the continuum fields, the Hamiltonian is obtained as
$$H=iv𝑑x[\psi _R^{}\frac{d}{dx}\psi _R+\psi _L^{}\frac{d}{dx}\psi _L].$$
(138)
(H.W. Exercise 3. Check that this Hamiltonian reduces to the one in Eq. (136) using Eqs. (137)).
We see that the corresponding Lagrangian density is just the standard one for free fermions given by
$$L=i\psi _R^{}(_t+v_x)\psi _R+i\psi _L^{}(_tv_x)\psi _L.$$
(139)
Using the standard rules of bosonization, this Lagrangian can also be rewritten as
$$L=\frac{1}{2v}(_t\varphi )^2\frac{v}{2}(_x\varphi )^2=\frac{1}{2}_\mu \varphi ^\mu \varphi ,$$
(140)
where the last equality requires setting $`v=1`$. Note: It is also worth checking to see that the same Hamiltonian in Eq. (138) is obtained by directly starting with the real space lattice model given in Eq. (132), rewriting the lattice fermions in terms of the continuum fermions remembering the $`\sqrt{a}`$ conversion factor, using $`_ia=𝑑x`$ and using $`\psi _{i+1}=\psi _i+a_x\psi _i`$.
Correlation functions
Thus, we have a Lorentz-invariant massless Dirac fermion field theory in the low energy approximation. All low energy properties can be obtained from the field theory, which in fact are trivially computed, since this is a free massless field theory. As far as fermionic correlation functions are concerned, one does not even require bosonization. However, for the spin correlations, it depends on how the spins can be expressed in terms of fermions. For instance, we can explicitly obtain the following spin-spin correlation function
$$G^{zz}(x,t)<S^z(x,t)S^z(0,0)>$$
(141)
simply using Wick theorem. We start by writing $`S_j^z`$ in term of the fermions as $`S_j^z=n_j1/2=\psi _j^{}\psi _j1/2=:\psi _j^{}\psi _j:`$ , since the expectation value of $`n_j`$ is half. Since the lattice fermion can be written in terms of the continuum fermions as
$$\psi _j=\sqrt{a}[e^{ik_Fj}\psi _R(x=ja)+e^{ik_Fj}\psi _L(x=ja)],$$
(142)
and since $`e^{i2k_Fj}=e^{(i\pi )x/a}=(1)^{x/a}`$, we find that the spin operator can be written as
$$S_j^z/a=S^z(x=ja,t)=:\psi _L^{}\psi _L:+:\psi _R^{}\psi _R:+(1)^{x/a}[\psi _R^{}\psi _L+\psi _L^{}\psi _R].$$
(143)
Directly using the Wick theorem and the fermion correlators
$`<T\psi _L(x,t)\psi _L^{}(0,0)>`$ $`=`$ $`{\displaystyle \frac{i}{2\pi (x_Li\alpha \mathrm{sign}(t))}}`$
$`\mathrm{and}<T\psi _R(x,t)\psi _R^{}(0,0)>`$ $`=`$ $`{\displaystyle \frac{i}{2\pi (x_R+i\alpha \mathrm{sign}(t))}},`$ (144)
we see that
$$G^{zz}(x,t)=\frac{a^2}{4\pi ^2}[(\frac{1}{x_R^2}+\frac{1}{x_L^2})(1)^{x/a}\frac{1}{x_Rx_L}],$$
(145)
where $`x_R=xt`$ and $`x_L=x+t`$.
(H.W. Exercise 4. Obtain this explicitly).
This can also be computed using bosonization. Note that even without doing the calculation, one could have guessed that the four-point correlation of the fermions must go as $`1/l^2`$, where $`l`$ is a distance, because in 1+1 dimensions, the fermion field has a mass dimension of 1/2 or distance dimension of $`1/2`$. So, in the absence of any other scale in the problem (the fermion field is massless and there are no interactions to cause divergences or introduce any anomalous mass scale), as long as the spin correlations can be expressed purely in terms of local fermion fields, no calculations are needed to see that correlations go as $`1/l^2`$. But we do need to calculate to get the explicit coefficients of $`1/x_R^2`$, etc, because they could be multiplied by dimensionless quantities like $`f(x_R/x_L)`$, $`etc`$.
However, to obtain the correlation function $`<S^+(x,t)S^{}(0,0)>`$ in the fermionic language is more difficult because of the non-local string operator. Here, simple dimensional analysis is not sufficient to give the answer and one actually needs bosonization. The correlation function can be written as
$`G^+(x,t)`$ $`=`$ $`<S^+(x,t)S^{}(0,0)>`$ (146)
$`=`$ $`(1)^{x/a}[e^{ik_Fx/a}\psi _R^{}(x,t)+e^{ik_Fx/a}\psi _L^{}(x,t)]\times `$
$`[e^{i\pi _0^x(:\psi ^{}(x^{},t)\psi (x^{},t):+1/2a)dx^{}}+h.c.]\times `$
$`[\psi _R(0,0)+\psi _L(0,0)],`$
where the string operator stretches between the two positions of the spin operator. (The other terms cancel out between $`S^{}`$ and $`S^+`$). Also, we have explicitly made the string operator hermitian, since it is hermitian in the lattice model. The reason bosonization comes in handy here is because the non-local operator when written in terms of bosons, turns out to be perfectly simple. We just use the bosonization identity
$`{\displaystyle _0^x}𝑑x^{}:\psi ^{}(x^{},t)\psi (x^{},t):`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\pi }}}{\displaystyle _0^x}𝑑x^{}_x^{}\varphi ={\displaystyle \frac{1}{\sqrt{\pi }}}[\varphi (x,t)\varphi (0,t)]`$ (147)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{\pi }}}[\varphi _R(x,t)+\varphi _L(x,t)\varphi _R(0,t)\varphi _L(0,t)].`$
Substituting this in Eq. (146), and substituting for the other fermion operators in terms of bosons, we get
$`G^+(x,t)`$ $`=`$ $`(1)^{x/a}{\displaystyle \frac{a}{2\pi \alpha }}[\eta _R^{}e^{ik_Fx/a}e^{i2\sqrt{\pi }\varphi _R(x,t)}+\eta _L^{}e^{ik_Fx/a}e^{i2\sqrt{\pi }\varphi _L(x,t)}]\times `$ (148)
$`[e^{ik_Fx/ai\sqrt{\pi }(\varphi _R(x)+\varphi _L(x)\varphi _R(0)\varphi _L(0))}+e^{ik_Fx/a+i\sqrt{\pi }(\varphi _R(x)+\varphi _L(x)\varphi _R(0)\varphi _L(0))}]\times `$
$`[\eta _Re^{i2\sqrt{\pi }\varphi _R(0,0)}+\eta _Le^{i2\sqrt{\pi }\varphi _L(0,0)}]`$
fully in terms of bosons. Now we use the operator identity $`e^{A+B}=e^Ae^Be^{[A,B]/2}`$ to write each of the 8 terms that appear in the above equation in terms of products of exponential factors. Just for illustration, we explicitly write the first term which appears by multiplying the first term in each of the square brackets in the above equation.
$$G^+(x,t)=\frac{a}{2\pi \alpha }[\eta _R^{}\eta _Re^{i\sqrt{\pi }\varphi _R(x,t)}e^{i\sqrt{\pi }\varphi _L(x,t)}e^{i\sqrt{\pi }\varphi _R(0,0)}e^{i\sqrt{\pi }\varphi _L(0,0)}+7\mathrm{other}\mathrm{terms}].$$
(149)
Now, we use the standard commutators
$`[\varphi _{R/L}(x),\varphi _{R/L}(y)]=(/+)i\mathrm{sign}(xy)/4`$, and $`[\varphi _{R/L}(x),\varphi _{L/R}(y)]=0`$ (since we are using Klein factors), and the standard algorithm for computing the correlation function
$`<e^{i2\sqrt{\pi }m_1\varphi _L(x)}e^{i2\sqrt{\pi }m_2\varphi _L(0)}>`$ $``$ $`Lim_{\alpha 0}\left({\displaystyle \frac{\alpha }{x_Li\alpha \mathrm{sign}(t)}}\right)^{m_1m_2}`$
$`\mathrm{and}<e^{i2\sqrt{\pi }m_1\varphi _R(x)}e^{i2\sqrt{\pi }m_2\varphi _R(0)}>`$ $``$ $`Lim_{\alpha 0}\left({\displaystyle \frac{\alpha }{x_R+i\alpha \mathrm{sign}(t)}}\right)^{m_1m_2},`$ (150)
when $`m_1`$ and $`m_2`$ have the same sign and vanish when they have opposite signs . This implies that of the 8 terms above, four of them give zero contribution. Adding up the contributions of the remaining four, we obtain
$$G^+(x,t)\frac{1}{(x_Rx_L)^{1/4}}[(1)^{x/a}+\mathrm{const}(\frac{1}{x_R^2}+\frac{1}{x_L^2})].$$
(151)
Note that the Klein factors always come as $`\eta _i^{}\eta _i=1`$ in this correlation function. Also note that one cannot fix the arbitrary constant that can appear between the uniform and the alternating parts of the correlation function because of the normal ordering ambiguities. It is only the exponents which can be found.
(H.W. Exercise 5. Obtain the above explicitly).
Thus even for the non-interacting theory or purely the $`XY`$ model, bosonization comes in handy to compute the correlation functions. As we have already said, the reason the correlation functions are not obtainable just by naive scaling arguments is because the expression for the ‘off-diagonal’ spin correlations in terms of the fermion operators is non-trivial, because of the presence of the string term. These are the only non-zero correlators in the theory. The other correlators such as $`G^{z+}`$ or $`G^{++}`$ are zero by symmetry - $`i.e`$ because of $`U(1)`$ invariance in the spin model or because of charge conservation in the fermion model.
Case when $`J^z0`$
We now consider the Hamiltonian in Eq. (132). In the fermionic language, the last term is given by
$$\delta H=J_z\underset{j}{}:\psi _j^{}\psi _j::\psi _{j+1}^{}\psi _{j+1}:.$$
(152)
At very large $`J_z`$, we would expect electrons to be localized on every alternate site so that adjacent sites are not occupied. However, this will not be true for small $`J_z`$, so the point of the exercise is to see when this happens and what the ground state looks like, for both small $`J_z`$ and large $`J_z`$. In the low energy limit, we can rewrite this term in terms of the continuum Dirac fermions at the Fermi points (use Eq. (142)) as
$`\delta H=aJ_z{\displaystyle 𝑑x}`$ $`[`$ $`:\psi _R^{}(x)\psi _R(x)+\psi _L^{}(x)\psi _L(x)+(1)^{x/a}M(x):]\times `$
$`[`$ $`:\psi _R^{}(x+a)\psi _R(x+a)+\psi _L^{}(x+a)\psi _L(x+a)+(1)^{x/a+1}M(x+a):],`$
$`\mathrm{where}M(x)`$ $`=`$ $`\psi _R^{}(x)\psi _L(x)+\psi _L^{}(x)\psi _R(x).`$ (153)
Using the notation $`\rho _L(x)=\psi _L^{}(x)\psi _L(x)`$ and $`\rho _R(x)=\psi _R^{}(x)\psi _R(x)`$, ($`\rho _R+\rho _L`$ is the charge density, and $`\rho _R\rho _L`$ is the current density; $`\rho _L`$ and $`\rho _R`$ are also called the left and right moving currents respectively), we can rewrite Eq. (153) as
$`\delta H=aJ_z{\displaystyle 𝑑x}`$ $`[`$ $`\rho _R(x)\rho _R(x+a)+\rho _L(x)\rho _L(x+a)+\rho _R(x)\rho _L(x+a)+\rho _L(x)\rho _R(x+a)`$ (154)
$`M(x)M(x+a)].`$
Here we have used the fact that oscillatory factors integrate to zero. (More precisely, they give rise to higher dimension operators, which, however, are irrelevant and ignored in this analysis). In the current-current terms, we can use the expansions $`\rho _L(x+a)=\rho _L(x)+a_x\rho _L(x),\psi _L(x+a)=\psi _L(x)+a_x\psi _L(x)`$, $`etc.`$, and the fact that square of a Fermi field vanishes, $`e.g.`$, $`\psi _L^2(x)=0`$, to deduce that terms of the form $`\rho _L(x)\rho _L(x+a)`$ are higher dimension operators ( they have four fermion operators and at least one derivative term) and renormalize to zero in the $`a0`$ limit. So among the current-current terms, we are only left with $`\rho _L`$-$`\rho _R`$ cross terms of the form $`\rho _L(x)\rho _R(x)`$ as the lowest dimension operators. For the four fermion terms in the second line also, we apply the same expansion. Dropping higher derivative terms, we see that the only term which survives in the product of the curly brackets is of the form $`\rho _R(x)\rho _L(x)\rho _L(x)\rho _R(x)`$. The extra negative sign is because we need to anticommute one of the fields. This adds to the $`\rho _R\rho _L`$ term coming from the first line and we are finally left with
$$\delta H=4J_za𝑑x\rho _L(x)\rho _R(x).$$
(155)
This is a four fermion term which in continuum quantum field theory is called the Thirring term. In the fermionic language, this is an interacting quantum field theory. However, it is easy to solve by bosonization.
By the standard rules of bosonization for non-interacting fermions, we can write
$$\rho _L=\frac{1}{2\sqrt{\pi }}(\frac{1}{v}_t+_x)\varphi ,\mathrm{and}\rho _R=\frac{1}{2\sqrt{\pi }}(\frac{1}{v}_t+_x)\varphi .$$
(156)
In that case, using the units that $`Ja=v=1`$, we get
$$\delta L=\delta H=\frac{J_z}{J\pi }_\mu \varphi ^\mu \varphi ,$$
(157)
where $`\delta L`$ denotes the change in the Lagrangian. This is precisely of the same form as the bosonization of the free fermion Hamiltonian. So the new Lagrangian is given by
$$L=\frac{1}{2K}_\mu \varphi ^\mu \varphi ,$$
(158)
where
$$\frac{1}{K}=1+\frac{2J_z}{J\pi }.$$
(159)
This can be made to look like the free term by redefining the field $`\varphi `$ \- $`i.e.`$, we define a new field $`\stackrel{~}{\varphi }=\varphi /\sqrt{K}`$, so that in terms of $`\stackrel{~}{\varphi }`$, the Lagrangian is just $`\frac{1}{2}(_\mu \stackrel{~}{\varphi }^\mu \stackrel{~}{\varphi })`$. However, the canonical momentum obtained from the rescaled Lagrangian is just $`\stackrel{~}{\mathrm{\Pi }}=_0\stackrel{~}{\varphi }`$ whereas the momentum obtained from the Lagrangian in Eq. (158) $`\mathrm{\Pi }=_0\varphi /K`$. Hence, the momentum gets rescaled compared to the original momentum as $`\stackrel{~}{\mathrm{\Pi }}=\sqrt{K}\mathrm{\Pi }`$. Clearly, the new coordinate and momenta satisfy the canonical commutation relations, since they are rescaled in opposite ways. (Remember that $`\varphi `$ takes values on a compact circle, since the original spin operators are defined in terms of exponential of the boson fields and are invariant under periodic changes of $`\varphi `$). But since the right and left mover fields are defined by taking both the field and the canonical momentum, and they scale in different ways, one can no longer write the right and left moving fields in the tilde representation as just scaled versions of the right and left moving fields of the original theory - in fact, they mix up the left and right moving fields. Explicitly,
$`\stackrel{~}{\varphi }_R(t,x)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[\stackrel{~}{\varphi }(t,x){\displaystyle _{\mathrm{}}^x}𝑑x^{}\stackrel{~}{\mathrm{\Pi }}(t,x^{})]={\displaystyle \frac{1}{2}}[\varphi (t,x)/\sqrt{K}{\displaystyle _{\mathrm{}}^x}𝑑x^{}\sqrt{K}\mathrm{\Pi }(t,x^{})]`$ (160)
$`=`$ $`{\displaystyle \frac{(\varphi _R+\varphi _L)}{2\sqrt{K}}}{\displaystyle \frac{\sqrt{K}(\varphi _L\varphi _R)}{2}}`$
$`=`$ $`{\displaystyle \frac{1}{2}}({\displaystyle \frac{1}{\sqrt{K}}}+\sqrt{K})\varphi _R+{\displaystyle \frac{1}{2}}({\displaystyle \frac{1}{\sqrt{K}}}\sqrt{K})\varphi _L=\mathrm{cosh}\beta \varphi _R+\mathrm{sinh}\beta \varphi _L,`$
where $`e^\beta =\sqrt{K}`$. Similarly
$$\stackrel{~}{\varphi }_L(t,x)=\mathrm{cosh}\beta \varphi _L+\mathrm{sinh}\beta \varphi _R.$$
(161)
One can now express the spin fields in terms of the $`\stackrel{~}{\varphi }`$ fields. They are given by
$`S^z(x,t)`$ $``$ $`\sqrt{{\displaystyle \frac{K}{\pi }}}_x\stackrel{~}{\varphi }+(1)^{x/a}\mathrm{const}e^{i2\sqrt{\pi K}\stackrel{~}{\varphi }}`$
$`S^{}(x,t)`$ $``$ $`(1)^{x/a}e^{i\sqrt{\pi /K}(\stackrel{~}{\varphi }_R\stackrel{~}{\varphi }_L)}+\mathrm{const}\times `$
$`[e^{i(2\sqrt{\pi K}(\stackrel{~}{\varphi }_R+\stackrel{~}{\varphi }_L)+\sqrt{\pi /K}(\stackrel{~}{\varphi }_R\stackrel{~}{\varphi }_L))}+e^{i(2\sqrt{\pi K}(\stackrel{~}{\varphi }_R+\stackrel{~}{\varphi }_L)+\sqrt{\pi /K}(\stackrel{~}{\varphi }_R\stackrel{~}{\varphi }_L))}].`$ (162)
With these substitutions, it is trivial (albeit algebraically more tedious!) to recalculate the spin-spin correlators $`G^{zz}(x,t)`$ and $`G^+(x,t)`$. Since the method is exactly the same as for the free case, we just quote the answers here.
$`G^{zz}(x,t)`$ $``$ $`{\displaystyle \frac{K}{4\pi }}({\displaystyle \frac{1}{x_L^2}}+{\displaystyle \frac{1}{x_R^2}})+(1)^{x/a}\mathrm{const}(x_Rx_L)^K`$
$`G^+(x,t)`$ $``$ $`(1)^{x/a}(x_Rx_L)^{1/4K}+\mathrm{const}(x_Rx_L)^{(\frac{1}{2\sqrt{K}}\sqrt{K})^2}({\displaystyle \frac{1}{x_L^2}}+{\displaystyle \frac{1}{x_R^2}}).`$ (163)
(H.W. Exercise 6: Obtain the above expressions).
Note that at $`K=1/2`$, the two correlations above are the same.
Limitations of this calculation
So the end result is that we have now obtained spin-spin correlation functions even including $`J_z`$. But since we have made a low energy continuum approximation and included only a few low-lying modes around the Fermi point, this derivation of the correlation functions is not true for arbitrary $`J_z`$. For instance, we left out terms that were irrelevant by naive power counting, which only works in the non-interacting case. Once we have interactions, some of those operators could acquire anomalous dimensions and hence become relevant. In other words, we have seen that interactions change the dimensions of operators. However, we have only studied operators of the form $`\rho _L(x)\rho _R(x)`$, which were marginal to start with and seen how they evolved. But we did not keep all the irrelevant operators and see how they evolved. Sometimes, they will also become relevant with sufficiently strong interactions.
A more general effective action approach
However, one can try to understand what can possibly change if we include other corrections that we left out in our approximation. One way of doing this is to look at all possible relevant terms that can appear consistent with the symmetries of the problem. The idea is not to try and derive these terms but to write them down in the effective Lagrangian assuming that if they are not explicitly prohibited by a symmetry, then they will appear. This is the philosophy behind what are called effective field theories.
Aside on how to ‘read off’ dimensions of operators
We know that to see whether an operator is relevant or irrelevant, we have to compute its correlation function and find out its scaling dimension. Then, we have to check whether the scaling dimension is such that the coefficient of the operator grows or becomes smaller as the energy scale is reduced. So given any operator $`O_i`$ in terms of bosons, we first compute the correlation function $`<O_i(x,t)O_i(x^{},0)>`$ which goes as $`1/(xx^{})^{2d_i}`$. For a free fermion theory with no interactions, (equivalently a boson theory with the interaction parameter $`K=1`$) we know that $`<O_i(x,t)O_i(x^{},0)>=1/(xx^{})^{2\stackrel{~}{d}_i}`$ where $`\stackrel{~}{d}_i`$ is just the naive scaling dimension or the engineering dimension of the operator $`O_i`$. The difference between $`\stackrel{~}{d}_i`$ and the $`d_i`$ that appears when we actually compute the correlation function is because of the interactions and is called the anomalous dimension of the operator. As was explained in the other courses in this school , the extra dimensional parameter comes from the cutoff scale. We shall use the term scaling dimension to mean $`d_i`$ itself. For an operator of the form $`O_ie^{i2\sqrt{\pi }\beta (\varphi _L+\varphi _R)}`$, the scaling dimension is given by $`d_i=\beta ^2`$, for the standard (non-interacting) form of the Hamiltonian. Since the space-time dimension is two, it is clear that $`d_i>2`$ implies that the coefficient $`\lambda _i`$ of the operator has to have dimension $`2d_i<0`$. So each time the cutoff is scaled down by a factor $`\mathrm{\Lambda }`$, $`\lambda _i\lambda _i\mathrm{\Lambda }^{2d_i}`$. Hence, after successive rescalings, this term in the action is irrelevant and scales to zero. On the other hand, $`d_i<2`$ denotes relevant operators, whose coefficients grow under scaling downs of the cutoff. $`d_i=2`$ is a marginal operator, whose coefficient remains unchanged under rescalings. (We will come back to this when we study impurity scattering and scaling dimensions of ‘boundary operators’).
Back to the effective action
The only possible Lorentz-invariant relevant terms that can be added to the Lagrangian is either $`\mathrm{cos}2\sqrt{\pi }\beta (\stackrel{~}{\varphi }_L+\stackrel{~}{\varphi }_R)`$ or $`\mathrm{cos}2\sqrt{\pi }\beta (\stackrel{~}{\varphi }_L\stackrel{~}{\varphi }_R)`$; both of these have dimension $`\beta ^2`$ and are thus relevant for $`\beta <\sqrt{2}`$. (The real problem on a lattice, of course, does not have Lorentz invariance. However, in the long distance or low energy limit, all such Lorentz non-invariant interactions will probably be irrelevant). Of these, the $`U(1)`$ symmetry under rotations about the $`z`$-axis in fermion language implies that $`\psi _L`$ and $`\psi _R`$ have to be multiplied by the same phase (because $`S_z`$ which has terms of the form $`\psi _L^{}\psi _R+h.c.`$ should not change). This in turn means that $`\stackrel{~}{\varphi }_L\stackrel{~}{\varphi }_L+c`$ and $`\stackrel{~}{\varphi }_R\stackrel{~}{\varphi }_Rc`$ so that $`\stackrel{~}{\varphi }_L+\stackrel{~}{\varphi }_R\stackrel{~}{\varphi }_L+\stackrel{~}{\varphi }_R`$ and $`\stackrel{~}{\varphi }_L\stackrel{~}{\varphi }_R\stackrel{~}{\varphi }_L\stackrel{~}{\varphi }_R`$+constant. Thus, to be consistent with this symmetry, we can only allow $`\mathrm{cos}2\sqrt{\pi }\beta (\stackrel{~}{\varphi }_L+\stackrel{~}{\varphi }_R)`$.
Furthermore, since the spin operators are all expressed in terms of exponentials of the boson fields, (see Eq. (162)) the boson fields need to be ’compactified on a circle’. This only means that the boson fields are periodic -
$$\stackrel{~}{\varphi }\stackrel{~}{\varphi }+\sqrt{\frac{\pi }{K}}$$
(164)
since the spin fields cannot distinguish between $`\stackrel{~}{\varphi }`$ and $`\stackrel{~}{\varphi }+\sqrt{\pi /K}`$. This restricts $`\beta `$ in $`\mathrm{cos}2\sqrt{\pi }\beta (\stackrel{~}{\varphi }_L+\stackrel{~}{\varphi }_R)`$ to be of the form $`n\sqrt{K}`$ where $`n`$ is an integer.
Finally, we use an unusual feature which occurs in the continuum field theories of many lattice spin models. The translational symmetry of the lattice spin model by one site ( or more sites for more general models) maps to a discrete symmetry in the continuum model, which is distinct from translational symmetry. This can be seen from the continuum definition of the spin in terms of the Dirac fermions - Eq. (143). When we change $`j`$ to $`j+1`$, the oscillatory factor $`(1)^j(1)^j`$. This is a drastic change from site to site. So if in the continuum version, we want to define smooth fields without having this rapid oscillations, we need to define one field for every pair of sites. Thus invariance under translation by $`2a`$ on the lattice becomes translational invariance in the continuum model. But from Eq. (143), we see that translational symmetry through a single site corresponds to the discrete symmetry
$$\psi _Li\psi _L,\psi _Ri\psi _R.$$
(165)
In the bosonic language, this corresponds to
$$\stackrel{~}{\varphi }_L\stackrel{~}{\varphi }_L+\frac{1}{2}\sqrt{\frac{\pi }{K}},\mathrm{and}\stackrel{~}{\varphi }_R\stackrel{~}{\varphi }_R+\frac{1}{2}\sqrt{\frac{\pi }{K}}.$$
(166)
This symmetry implies that the only terms that can be added to the Lagrangian are of the form $`\mathrm{cos}2\sqrt{\pi }2n\sqrt{K}(\stackrel{~}{\varphi }_L+\stackrel{~}{\varphi }_R)`$, so that $`\beta =2n\sqrt{K}`$. This is relevant when $`K<1/2`$ when $`n=1`$. So the system is in a massless phase till $`K`$ reaches $`1/2`$ below which it develops a relevant interaction, and a mass gap.
However, we cannot use our low energy approximate result to estimate the point at which the spin model develops a relevant interaction. Besides adding the possible relevant term mentioned above, the most general thing the other terms that we have neglected can do is to change the relation between $`K`$ and $`J_z`$ in an unpredictable way. In fact, the low energy result relating $`K`$ to the perturbation $`J_z`$ is only true to lowest order in $`J_z/J`$. This particular spin-chain model is, in fact, solvable by Bethe ansatz and the exact answer is
$$\frac{1}{K}=1+\frac{2}{\pi }\mathrm{sin}^1(\frac{J_z}{J})$$
(167)
which, to lowest order in $`J_z/J`$, gives us the relation in Eq. (159). From this, we see that $`1/K2`$ at $`J_z=J`$. This is precisely the $`K`$ value for which the cosine interaction term becomes relevant. To prove that a relevant interaction necessarily leads to a mass gap is non-trivial, but it is certainly plausible. Once, there is a relevant interaction, its coefficient grows under renormalization. It becomes divergent as we make the energy scale lower and lower, so we have to cut it off at some scale, which is the mass scale associated with the theory. However, it could also lead us to a new fixed point, which may not have a mass gap.
To get higher orders in $`J_z/J`$ in this effective field theory approach is not easy, because one needs to go beyond the region of linear dispersion. Also, once one starts including modes with $`kk_F`$, we need to be careful to put in the restriction that $`k_F<k<k_F`$. Also, since the Bethe ansatz already gives the exact answer to all orders, there may not be much point in trying to do this for this problem.
So what does all this formalism gain us? How does the ground state evolve as $`J_z`$ changes? For small values of $`J_z`$ all that happens is that the spin-spin correlations have a slightly different power law fall-off with anomalous non-integer exponents. Does this continue for all values of $`J_z`$? No. Once, $`J_z`$ reaches $`J_z=J`$, the isotropic point, there exists a phase transition to a massive phase where spin-spin correlations fall off exponentially fast at large separations. In this particular problem, of course, one knew this answer from the Bethe ansatz, but the point is that the bosonization method can be used even for other models, which are not exactly solvable by the Bethe ansatz. But without the Bethe ansatz, one cannot analytically find the value of the parameter where the phase transition into a massive phase occurs. The other important gain that we have in this method is that it allows the computation of correlation functions, which is not possible using the Bethe ansatz. Finally, since it is a symmetry analysis, it tells us that for any Hamiltonian of this type, the model is likely to be massless only when the theory has $`U(1)`$ symmetry and the $`Z_2`$ symmetry of translation by a single site and even then, only for some restricted values of the parameter space.
The best reference for this application is Affleck’s lectures on field theories and critical phenomena, which we have followed fairly faithfully.
## 7 Hubbard model
The Hubbard model is one of the simplest realistic models that one can study which has a competition between the kinetic energy and the potential energy. The kinetic energy or the hopping term gains, or rather the energy gets lowered, if the fermions are delocalized - free to move throughout the sample. In this model, the potential energy represents screened Coulomb interactions between electrons and the model is constructed so that it costs energy to put two electrons at the same place. So the potential energy prefers each electron to sit at its own site. The model is given by
$$H=\frac{t}{2}\underset{j\alpha }{}(\psi _{j\alpha }^{}\psi _{j+1\alpha }+h.c.)+U\underset{j}{}n_jn_j+\mu \underset{j,\alpha }{}\psi _{j\alpha }^{}\psi _{j\alpha },$$
(168)
where $`t`$ is the hopping parameter, $`U`$ is a positive constant denoting the repulsion between two electrons at a site, $`\mu `$ is the chemical potential and $`\alpha `$ is the spin index which can be $``$ or $``$. This model is very similar to the fermion model we studied for the spin chain except that the electrons have spin and the chemical potential term allows for arbitrary fillings. The $`U`$ term or Hubbard term is analogous to the nearest neighbor $`J_z`$ interaction term for spinless electrons.
At half-filling (one electron/site, since a filled band implies two electrons/site), for large $`U`$, the model is expected to describe an insulator. One can easily understand this, because at infinite $`U`$, the ground state will have one electron at every site. Any excitation will cost an energy of $`U`$. So there is a gap to excitations and the model behaves as an insulator. It is called a Mott-Hubbard insulator (as opposed to other band insulators) because here the insulating gap is created by interactions.
The question that one would like to ask is, at what value of $`U`$ does the Mott-Hubbard gap open, because naively one may think that at very small values of $`U`$, the model allows free propagation of electrons and describes a metal. Using bosonization, we will show that in one dimension, this expectation is wrong. The Mott-Hubbard gap opens up for any finite $`U`$ if the filling is half and not otherwise. For any other filling, the model at low energies is an example of a Luttinger liquid with separate spin and charge excitations. The spin modes are always gapless whereas the charge modes are gapless at any filling other than half-filling; precisely at half-filling a charge gap opens up. The model for arbitrary filling (other than half-filling) and positive $`U`$ is said to be in the Luttinger liquid phase. Spin and charge correlations fall-off as power laws and we expect power law transport. At half-filling, the model is in a charge-gapped phase called the charge-density wave phase.
Aside: For negative $`U`$, it is found that the spin excitations are always gapped. Here, the model is said to be in the Luther-Emery phase or spin gapped phase.
Bosonization of the model without interactions
How do we go about seeing all that we have described above? In higher dimensions, we would do a mean field theory, but in one space dimension, we know that a mean field analysis is not very useful because of the infrared divergences of the low energy fluctuations. (In other words, if we write down a mean field theory and then try to do systematic corrections about the mean field theory, then order by order in perturbation theory, we find that the integrals which appear in the corrections are divergent). So it seems like a good idea to try and use bosonization. In fact, the way this model is analyzed is very similar to the way we analyzed the spinless fermion model in the previous section. We first switch off the interactions and start with the Fourier decomposition
$$\psi _{j\alpha }=\frac{1}{\sqrt{N}}\underset{k}{}\psi _{k\alpha }e^{ikja}.$$
(169)
We rewrite the Hamiltonian as
$$H_0=\underset{k\alpha }{}(\mu t\mathrm{cos}ka)\psi _{k\alpha }^{}\psi _{k\alpha },$$
(170)
where the $`k`$ values go from $`\pi /a`$ to $`\pi /a`$. In the ground state, all states with $`|k|<k_F`$ are filled, where $`k_F`$ is determined by the chemical potential from the equation $`\mu =\mathrm{cos}k_Fa`$. Just as in the spinless case, we will look only at the low energy modes near the Fermi surface, so that each fermion is written as
$`\psi _{j\alpha }`$ $``$ $`e^{ik_Fja}{\displaystyle _{k_Fa\mathrm{\Lambda }}^{k_Fa+\mathrm{\Lambda }}}{\displaystyle \frac{d(ka)}{2\pi }}e^{ikja}\psi _{k\alpha }+e^{ik_Fja}{\displaystyle _{k_Fa\mathrm{\Lambda }}^{k_Fa+\mathrm{\Lambda }}}{\displaystyle \frac{d(ka)}{2\pi }}e^{ikja}\psi _{k\alpha }`$ (171)
$``$ $`e^{ik_Fja}\psi _{Lj\alpha }+e^{ik_Fja}\psi _{Rj\alpha },`$
so that the $`\psi _{Lj\alpha }`$ and $`\psi _{Rj\alpha }`$ do not contain high energy modes. Substituting this expression in Eq. (168), we get
$`H_0={\displaystyle \frac{t}{2}}{\displaystyle \underset{j\alpha }{}}`$ $`[`$ $`(e^{ik_Fa}\psi _{Lj\alpha }^{}\psi _{Lj+1\alpha }+e^{ik_Fa}\psi _{Rj\alpha }^{}\psi _{Rj+1\alpha }+`$
$`e^{i2k_Fjaik_Fa}\psi _{Lj\alpha }^{}\psi _{Rj+1\alpha }+e^{i2k_Fja+ik_Fa}\psi _{Rj\alpha }^{}\psi _{Lj+1\alpha })+h.c.]`$
$`+\mu {\displaystyle \underset{j\alpha }{}}`$ $`[`$ $`\psi _{Lj\alpha }^{}\psi _{Lj\alpha }+\psi _{Rj\alpha }^{}\psi _{Rj\alpha }+e^{i2k_Fja}\psi _{Lj\alpha }^{}\psi _{Rj\alpha }+e^{i2k_Fja}\psi _{Rj\alpha }^{}\psi _{Lj\alpha }].`$ (172)
The oscillatory terms do not contribute because they have the form
$$\underset{j}{}e^{i2k_Fja}\psi _{Rj\alpha }^{}\psi _{Lj\alpha }\underset{j}{}e^{i2k_Fja}\underset{k,k^{}}{}e^{ikja}e^{ik^{}ja}\psi _{k\alpha }^{}\psi _{k^{}\alpha }\underset{k,k^{}}{}\psi _{k\alpha }^{}\psi _{k^{}\alpha }\underset{j}{}e^{i2k_Fjaikja+ik^{}ja},$$
(173)
and the sum over $`j`$ in the last expression produces $`\delta _{2k_F,kk^{}}`$ which cannot be satisfied for small values of $`k,k^{}`$. (Note that even for the half-filled case where $`k_Fa=\pi /2`$, this cannot be satisfied). Hence, we may drop these terms and we are left with only
$`H_0=`$ $``$ $`{\displaystyle \frac{t}{2}}{\displaystyle \underset{j\alpha }{}}[(e^{ik_Fa}\psi _{Lj\alpha }^{}\psi _{Lj+1\alpha }+e^{ik_Fa}\psi _{Rj\alpha }^{}\psi _{Rj+1\alpha })+h.c.]`$ (174)
$`+`$ $`\mu {\displaystyle \underset{j\alpha }{}}(\psi _{Lj\alpha }^{}\psi _{Lj\alpha }+\psi _{Rj\alpha }^{}\psi _{Rj\alpha }).`$
Now, we expand $`\psi _{j+1\alpha }\psi _\alpha (j+1)=\psi _\alpha (j)+a_x\psi _\alpha (j)`$ \+ higher order irrelevant terms and use the fact that $`\mathrm{cos}k_Fa=\mu `$ (which means that part of the hopping term cancels with the $`\mu `$ term) to get
$`H`$ $`=`$ $`{\displaystyle \frac{at}{2}}{\displaystyle \underset{j\alpha }{}}(e^{ik_Fa}\psi _{Lj\alpha }^{}_x\psi _{Lj\alpha }+e^{ik_Fa}\psi _{Rj\alpha }^{}_x\psi _{Rj\alpha }+h.c.)`$ (175)
$`=`$ $`iat\mathrm{sin}(k_Fa){\displaystyle \underset{j\alpha }{}}(\psi _{Lj\alpha }^{}_x\psi _{Lj\alpha }\psi _{Rj\alpha }^{}_x\psi _{Rj\alpha }),`$
where we have also integrated the hermitian conjugate terms by part to get it in the form of the second equation above. Finally, we can rewrite this Hamiltonian as a continuum Hamiltonian in terms of continuum fields (defined with the usual factor of $`\sqrt{a}`$ as $`\psi _\alpha (j)/\sqrt{a}=\psi _\alpha (x)`$ ) and using $`_ja=𝑑x`$
$$H_0=it\mathrm{sin}(k_Fa)\underset{\alpha }{}𝑑x[\psi _{L\alpha }^{}(x)_x\psi _{L\alpha }(x)\psi _{R\alpha }^{}(x)_x\psi _{R\alpha }(x)],$$
(176)
where we call $`t\mathrm{sin}(k_Fa)=v_Fa`$, the Fermi velocity times $`a`$. The derivation here is very similar to the one for spinless fermions, except that here we have carried it out in real space instead of momentum space. This Hamiltonian can be bosonized using the usual rules of bosonization and we get
$$H_0=\frac{v_Fa}{2}\underset{\alpha }{}𝑑x[\mathrm{\Pi }_\alpha ^2+(_x\varphi _\alpha )^2].$$
(177)
(H.W. Exercise 7. Derive the Hamiltonian in Eq. (176) through a momentum space derivation).
Bosonization of the interaction term
The next step is to figure out the low energy part of the on-site Hubbard interaction. Here, again, the principle is the same. We rewrite the four-fermion term written in terms of the original fermions in terms of the low energy Dirac fermion modes. Just as in the spin model, the $`S^zS^z`$ term or the four fermion term corresponded to a product of normal ordered bilinears, here also the four fermion term in Eq. (168) can be written in terms of the product of normal ordered bilinears if we subtract the average charge densities of the $``$ and $``$ fields. So we may write
$$H_{\mathrm{int}}=U\underset{j}{}n_jn_j=U\underset{i}{}:n_j::n_j:.$$
(178)
In terms of the Dirac fields, this becomes
$$H_{\mathrm{int}}=U\underset{j}{}[(:\psi _{jL}^{}\psi _{jL}:+:\psi _{jR}^{}\psi _{jR}:+\psi _{jR}^{}\psi _{jL}e^{i2k_Fja}+\psi _{jL}^{}\psi _{jR}e^{i2k_Fja})\times ()].$$
(179)
We now expand the products and keep only the terms with no oscillatory factor, to get
$$H_{\mathrm{int}}=U\underset{j}{}(J_{jR}+J_{jL})(J_{jR}+J_{jL})+U\underset{j}{}(\psi _{jR}^{}\psi _{jL}\psi _{jL}^{}\psi _{jR}+h.c.).$$
(180)
The remaining terms have the oscillatory factors of either $`e^{i2k_Fja}`$ or $`e^{i4k_Fja}`$ and can be set to zero for arbitrary filling. Notice however, that $`e^{i4k_Fja}=1`$ and is not oscillatory at half-filling since $`k_Fa=\pi /2`$. We will come back to this point later. Now we first express these fields in terms of the continuum fields and just use the standard bosonization formulae to get
$$H_{\mathrm{int}}=Ua𝑑x[\frac{1}{\pi }_x\varphi _{}_x\varphi _{}+\eta _R^{}\eta _L\eta _L^{}\eta _R\frac{Ua}{2\pi ^2ϵ^2}\mathrm{cos}\sqrt{4\pi }(\varphi _R+\varphi _L\varphi _R\varphi _L)].$$
(181)
(H.W. Exercise 8. Derive the above).
The interesting point to note here is that the cosine term only depends on $`\varphi _{}\varphi _{}`$. (We use the earlier defined notation that $`\varphi =\varphi _L+\varphi _R`$ and $`\theta =\varphi _R+\varphi _L`$). So if we define the charge and spin fields
$$\varphi _c=\frac{\varphi _{}+\varphi _{}}{\sqrt{2}},\mathrm{and}\varphi _s=\frac{\varphi _{}\varphi _{}}{\sqrt{2}},$$
(182)
the Hamiltonian is completely separable in terms of these two fields and we may write $`H=H_0+H_{\mathrm{int}}=H_c+H_s`$ with
$`H_c`$ $`=`$ $`{\displaystyle \frac{v_F}{2}}{\displaystyle 𝑑x[\mathrm{\Pi }_c^2+(1+\frac{U}{\pi v_F})(_x\varphi _c)^2]}`$
$`H_s`$ $`=`$ $`{\displaystyle \frac{v_F}{2}}{\displaystyle 𝑑x[\mathrm{\Pi }_s^2+(1\frac{U}{\pi v_F})(_x\varphi _s)^2+\frac{Ua}{2\pi ϵ^2}\mathrm{cos}\sqrt{8\pi }\varphi _s]},`$ (183)
where the bosonized form of the kinetic energy term given by $`H_0`$ in Eq. (177) along with the first term in Eq. (181) ($`U_x\varphi _{}_x\varphi _{}/\pi `$) can also be written in terms of the charge and spin fields as above. The charge sector is massless, but for the spin sector, one has a cosine term in the Hamiltonian. From our earlier experience of spinless models, we know that a cosine term can lead to a mass gap, when it becomes relevant. So we need to compute the dimension of the operator and see when it becomes relevant. Note that we have chosen the product of the Klein factors to be unity<sup>3</sup><sup>3</sup>3For single chain problems, the Klein factors usually cause no problems and can be set to be unity, in most cases. The only care that we need to take is to remember the negative sign that one gets when two of them are exchanged. But for multi-chain models, when more than four explicit Klein factors exist, one needs to be more careful.. But we only know how to compute correlation functions when the quadratic Hamiltonian is in the standard form. To get that, we need to rescale the $`\varphi `$ fields and their conjugate momenta (in the opposite way so that the commutation relations are preserved) as
$$\overline{\varphi }_c=(1+\frac{U}{\pi v_F})^{1/4}\varphi _c,\mathrm{and}\overline{\mathrm{\Pi }}_c=(1+\frac{U}{\pi v_F})^{1/4}\mathrm{\Pi }_c,$$
(184)
and similarly for the spin fields to get the Hamiltonian in the standard form, from which we can directly read out the dimensions of the operators. In terms of the bar fields, we see that
$`H_c`$ $`=`$ $`(1+{\displaystyle \frac{U}{\pi v_F}})^{1/2}{\displaystyle \frac{v_Fa}{2}}{\displaystyle 𝑑x[\overline{\mathrm{\Pi }}_c^2+(_x\overline{\varphi }_c)^2]}`$
$`H_s`$ $`=`$ $`(1{\displaystyle \frac{U}{\pi v_F}})^{1/2}{\displaystyle \frac{v_Fa}{2}}{\displaystyle 𝑑x[\overline{\mathrm{\Pi }}_s^2+(_x\overline{\varphi }_s)^2+\frac{Ua}{2\pi ϵ^2}\frac{1}{(1\frac{U}{\pi v_F})^{1/2}}\mathrm{cos}\sqrt{\frac{8\pi }{(1\frac{U}{\pi v_F})^{1/4}}}\overline{\varphi }_s]}.`$
The charge sector is purely quadratic (both before and after rescaling!) and remains massless, whereas for the spin sector, the rescaling was necessary to ‘read off’ the dimension of the cosine operator. Since its scaling dimension is given by $`d=2/(1\frac{U}{\pi v_F})^{1/4}`$, it is irrelevant ($`d<2`$) for any weak positive $`U`$ and the spin sector is also massless. On the other hand, for any negative $`U`$, this term has dimension $`d>2`$ and is relevant. As we explained in the spinless case, this means that the spin sector acquires a mass gap for all negative $`U`$.
Also note that the velocities of the charge and the spin modes have got renormalized in different ways. $`v_c=(1+\frac{U}{\pi v_F})^{1/2}v_F`$ is the velocity of the charge mode and $`v_s=(1\frac{U}{\pi v_F})^{1/2}v_F`$ is the velocity of the spin mode. Thus, spin and charge move independently. This is one of the hallmarks of Luttinger liquid behavior in one-dimensional fermion models. It is only for $`U=0`$, that the spin and the charge modes move together.
How does one look for such spin-charge separation in one-dimensional models? Experimentally, one has to look at different susceptibilities and measure the Wilson ratio, which is the ratio of the spin susceptibility to the specific heat coefficient. The specific heat coefficient depends both on spin and charge modes and is given by
$$\frac{\gamma }{\gamma _0}=\frac{1}{2}(\frac{v_F}{v_c}+\frac{v_F}{v_s}),$$
(186)
where $`\gamma _0`$ is the specific heat of non-interacting electrons with velocity $`v_F`$. However, spin susceptibility only depends on the spin mode and is given by
$$\frac{\chi }{\chi _0}=\frac{v_F}{v_s}.$$
(187)
Thus, the Wilson ratio is given by
$$R_W=\frac{\chi /\chi _0}{\gamma /\gamma _0}=\frac{2v_c}{v_c+v_s}.$$
(188)
Clearly, when there is no spin-charge separation, this is given by one. So deviations of the Wilson ratio from unity are a sign of spin-charge separation in real systems.
Finally, let us consider the case exactly at half-filling, $`k_Fa=\pi /2`$. In this case, the $`e^{i4k_Fja}`$ term we neglected in Eq. (179) as oscillatory, is no longer oscillatory, since $`e^{i4k_Fa}=1`$. In this case, there exists a term in the Hamiltonian of the form
$$H_{\mathrm{umklapp}}=U\underset{j}{}(\psi _{jR}^{}\psi _{jL}\psi _{jR}^{}\psi _{jL}+h.c.).$$
(189)
Note that this term destroys two right movers and creates two left movers or vice-versa. So there is an overall change in momentum by $`4k_Fa=2\pi `$, which has to be absorbed by the lattice. It is an umklapp process unlike the earlier interaction term for arbitrary filling which created and destroyed a particle at the left Fermi point and also created and destroyed a particle at the right Fermi point and did not change any momentum. It is easy to see that this term also gives rise to a cosine term by bosonizing, which, after rescaling gives $`\frac{Ua}{2\pi ϵ^2(1+U/\pi v_F)^{1/2}}\mathrm{cos}\sqrt{\frac{8\pi }{(1+U/\pi v_F)^{1/4}}}\overline{\varphi }_c`$ neglecting Klein factors. Thus it appears in the Hamiltonian of the charge sector. This term is irrelevant for any negative $`U`$, but relevant for any positive $`U`$. Thus, precisely at half-filling, the charge sector has a gap. This is similar to the case for spinless fermions where the spin model actually corresponded to a half-filled spinless fermion model. But unlike the case for the spinless fermions where the gap only opens up at $`J=J_z`$, here the gap opens up for any positive $`U`$, however small.
Is there any way one can understand these results in a physical way? For negative $`U`$, we found that the spin sector has a gap. This can be understood by saying that since there is an attractive interaction between the spin $``$ and the spin $``$ densities of fermions, they will like to form singlets and sit on a single site. So to make a spin excitation, one needs to break a pair and this costs energy. But charge excitations can move around as bound spin singlet pairs with no cost in energy. On the other hand, for positive $`U`$, there exists a repulsion between two electrons at a site. So each electron will tend to sit on a different site. At half-filling, hence, there is no way for an electron to move, without trying to sit at a site, at which an electron is already present. And this costs a repulsive energy $`U`$. Hence, there is a gap to charge excitations. But one can flip spins at a site and hence have spin excitations with no cost in energy.
So what results has bosonization given us here? We started with electrons with spin and charge moving together via a hopping term, but with a strong on-site Coulomb repulsion. This term could not be treated perturbatively. However, when we rewrote the theory in terms of bosons, with one boson for the $``$ spin and one for the $``$ spin, we found that the theory decoupled in terms of new spin and charge bosons. For a generic filling, the charge boson was just a massless free boson excitation, whereas the spin boson Hamiltonian had a cosine term, which was relevant when $`U`$ was negative, but irrelevant for positive $`U`$. But at half-filling, the charge excitations develop a gap, for any positive $`U`$. The most important thing to note here is that the charge and spin degrees of freedom have completely decoupled. Since the two fields are scaled differently, they move with different velocities in the system. This is a result that one could never have obtained perturbatively. Thus, at any filling other than half-filling, the low energy limit of the Hubbard model is a Luttinger liquid with massless spin and charge excitations moving with different velocities. A good reference for this part is Shankar’s article which also explains in great detail how to compute correlation functions.
## 8 Transport in a Luttinger liquid - Clean Wire
The last two applications involved the study of correlation functions, with the aim of finding out the different phases possible in a one-dimensional system of interacting fermions. In this part of the course, we will study another application of bosonization, which is to study transport, in particular the DC (or zero frequency) conductivity in one-dimensional wires of interacting fermions.
Firstly, are one-dimensional wires experimentally feasible? The general idea to make narrow wires is to ’gate’ 2D electron gases. In recent times, technology has developed enough to make these wires so narrow, that they contain only one transverse channel. So these are good enough approximations to one-dimensional wires. Another good approximation to coupled chains of one-dimensional models are carbon nanotubes, though those are not the kind of models we will study here.
The next point to note is that even at a qualitative level, transport in low dimensional systems is extremely different from transport in higher dimensions. To understand this point, we will first make qualitative statements about transport and conductivity before we explicitly start computing it using bosonization. The usual aim is to compute the conductance as a function of the voltage, temperature, presence of impurities or disorder and so on. Normally, when currents are measured in wires, one does not worry about quantum effects, because wires are still macroscopic objects, but that is clearly not the case here, since we are interested in one-dimensional wires. In fact, whenever the physical dimensions of the conductor becomes small, (it need not be really one-dimensional), the usual Ohmic picture of conductance where the conductance is given by
$$G=\sigma \frac{W}{L}=\sigma \frac{\mathrm{width}\mathrm{of}\mathrm{conductor}}{\mathrm{length}\mathrm{of}\mathrm{conductor}},$$
(190)
where $`\sigma `$ is a material dependent quantity called conductivity, breaks down. A whole new field called ‘mesoscopic physics’ has now been created to deal with electronic transport in such systems. The term ‘mesoscopic’ in between microscopic and macroscopic is used for systems, where the sizes of the devices are such that it is comparable with a) the de Broglie wavelength ( or kinetic energy) of the electron, b) the mean free path of the electron and c) the phase relaxation length ( the length over which the particle loses memory of its phase) of the electron. Ohmic behavior is guaranteed only when all these length scales are small compared to the size which happens for any macroscopic object. These lengths actually vary greatly depending on the material and also on the temperature. Typically, at low temperatures, they vary between a nanometer for metals to a micrometer for quantum Hall systems.
For mesoscopic wires, in general, quantum effects need to be taken into account. One way of computing these conductances is by using the quantum mechanical formulation of transmission and reflection through impurities and barriers. This formulation is called the Landauer-Buttiker formulation and works for Fermi liquids. However, it does not include interactions. But for one dimensional wires, interactions change the picture dramatically, since the quasi-particles are no longer fermion-like. Hence the Landauer-Buttiker formalism cannot be directly applied and one needs to compute conductances in Luttinger wires taking interactions into account right from the beginning. One way of doing this is by using bosonization and this is the method that we will follow here.
The aim is to compute the conductance of a one-dimensional wire. First, we will compute the conductance through a clean wire ( no impurities or barriers) and argue why the conductance is not renormalized by the interaction. Then we will study the conductance again after introducing a single impurity. Here, we will see that the interactions change the picture dramatically. For a non-interacting one-dimensional wire, from just solving usual one-dimensional quantum mechanics problems, we know that we can get both transmission and reflection depending on the strength of the scattering potential. But for an interacting wire, we shall find that for any scattering potential, however small, for repulsive interactions between the electrons, there is zero transmission and full reflection (implies conductance is zero, or that the wire is ‘cut’) and for attractive interactions between electrons (which is of course possible only for some renormalized ‘effective’ electrons), there is full transmission and zero reflection (implying perfect conductance or ‘healing’ of the wire).
Ballistic conductor
Let us first define the conductance of a mesoscopic ballistic conductor ($`i.e.`$, a conductor with no scattering) without taking interactions into account. We said earlier that the usual definition of conductance as $`G=\sigma \frac{W}{L}`$ breaks down for mesoscopic systems. For instance, it is seen that instead of the conductance smoothly going down as a function of the area or width of the wire $`W`$, it starts going down discretely in steps, each of height $`2e^2/h`$. Also as $`L`$ decreases, instead of increasing indefinitely, $`G`$ saturates at some limiting value $`G_c`$. The general understanding now, is that as the wire becomes thinner and thinner, the current is carried in a very few channels, each of them carrying a current of $`2e^2/h`$ (two for spin degeneracy) until we reach the lowest value which is just a single channel (which we interpret as the lowest eigenstate of the transverse Hamiltonian) carrying this current. Moreover, as the length decreases, the resistance does not decrease indefinitely but instead reaches a limiting value. One way of understanding this is to simply consider this to be a contact resistance, independent of the length of the wire, which arises simply because the conductor and the contacts are different. One cannot make the contacts the same as the conductor, because then our assumption that the voltage drop is across the conductor alone does not make sense. That makes sense only if we assume that the contacts are infinitely more conducting than the conductor. So we are finally left with a non-zero resistance and the wire does not become infinitely conducting. In fact, in this limit, the conductance or resistance of the wire is purely a ’boundary’ property and the ‘conductivity’ of the wire has no real significance. In fact, whether we get a finite conductivity or infinite conductivity depends on how one defines it.
However, for a single channel wire, clearly, the wire is one-dimensional and we know that interactions can change the picture drastically. The question that we want to answer here is precisely that. What is the conductance of a clean one-dimensional interacting wire or Luttinger wire?
Computing conductance of a clean one-dimensional (mesoscopic) wire
(a) Without leads
First, we shall perform a calculation to compute the conductance of a Luttinger liquid without any consideration of contacts or leads. (We shall restrict ourselves to spinless fermions since spin only increases the degrees of freedom and gives an overall multiplicative factor of two in the conductance). The conductance of a wire is calculated by applying an electric field to a finite region $`L`$ of an infinitely long wire and the current $`I`$ is related to the field as
$$I(x)=_0^L𝑑x^{}\frac{d\omega }{2\pi }e^{i\omega t}\sigma _\omega (x,x^{})E_\omega (x^{}),$$
(191)
where $`E_\omega (x^{})`$ is the frequency $`\omega `$ component of the time Fourier transform of the electric field. The conductivity $`\sigma _\omega (x,x^{})`$, in turn, is related to the (imaginary time) current-current correlation function by the usual Kubo formula as
$$\sigma _\omega (x,x^{})=\frac{e^2}{\overline{\omega }}_0^\beta 𝑑\tau <T_\tau j(x,\tau )j(x^{},0)>e^{i\overline{\omega }\tau },$$
(192)
where $`\tau =it`$, $`\omega =i\overline{\omega }+ϵ`$, $`T_\tau `$ is the (imaginary) time ordering operator and $`j(x,\tau )`$ is the current operator. Both these formulae are standard in many books on many body techniques, so here we will confine ourselves to just describing what they mean. The first equation describes the current as a response to an electric field (externally applied plus induced) of frequency $`\omega `$. The proportionality function is the conductivity. To get the usual Ohmic formula, all we need to do is replace $`\sigma =\sigma _0\delta (xx^{})`$ or remember that the $`\sigma (x,x^{})`$ is generally a function which is centered around $`xx^{}`$ and which falls off sufficiently fast elsewhere. The point for mesoscopic systems is that the length of the wire is roughly comparable with the range of $`\sigma (x,x^{})`$. Hence, the current gets contributions from the electric field all over the wire, which is different from what happens in the usual case, where the current at a point gets contributions only from the electric field very near that point. The second equation tells us that the conductivity is related to the current-current correlation function. This is derived by computing the current $`I(x)`$ in a Hamiltonian formulation to first order in the perturbation which is the applied electric field. The Euclidean formulation is used so that the generalization to finite temperature calculations is straightforward, but we shall only work at zero temperature and hence take the $`\beta \mathrm{}`$ limit.
Our aim here will be to compute the current-current correlation function and hence the conductance for a Luttinger wire using bosonization. We shall denote the Euclidean time action of a generic Luttinger liquid as
$$S_E=\frac{1}{2K}𝑑\tau 𝑑x[\frac{1}{v}(_\tau \varphi )^2+v(_x\varphi )^2].$$
(193)
(Note that in the spin model and Hubbard model, $`\tau `$ was replaced by $`it`$). The current can directly be expressed in terms of the boson operators as
$$j(x,\tau )v(\rho _R\rho _L)=\frac{i}{\sqrt{\pi }}_\tau \varphi .$$
(194)
(The extra factor of $`i`$ is because we are now using imaginary time $`\tau `$). Our first step is to obtain the correlation function $`<j(x,\tau )j(x^{},0)>`$ which is similar to the correlation functions for spinless fermions that we computed earlier when we were studying spin models, except that we are now interested in the Euclidean correlation functions. Since we can pull out the $`_\tau `$ outside the correlation function<sup>4</sup><sup>4</sup>4See R. Shankar in for subtleties in pulling out the derivative outside the time ordering operator. One gets an extra term which cancels another term that we have ignored here, a singular $`c`$-number term., all we have to do is compute the propagator given by
$$G(\tau ,x,x^{})=<T_\tau \varphi (x,\tau )\varphi (x^{},0)>,$$
(195)
or equivalently, its Fourier transform
$$G_{\overline{\omega }}(x,x^{})=_0^\beta 𝑑\tau <T_\tau \varphi (x,\tau )\varphi (x^{},0)>e^{i\overline{\omega }\tau }.$$
(196)
The conductivity is then given by
$`\sigma _{\overline{\omega }}(x,x^{})`$ $`=`$ $`{\displaystyle \frac{e^2}{\overline{\omega }\pi }}{\displaystyle _0^\beta }𝑑\tau <T_\tau _\tau \varphi (x,\tau )_\tau \varphi (x^{},0)>e^{i\overline{\omega }\tau }`$ (197)
$`=`$ $`{\displaystyle \frac{e^2\overline{\omega }}{\pi }}G_{\overline{\omega }}(x,x^{}).`$
So now, to compute the conductance, all we have to do is compute the propagator for the boson with a free Euclidean action. The propagator satisfies the equation
$$\frac{1}{K}(v_x^2+\frac{\overline{\omega }^2}{v})G_{\overline{\omega }}(x,x^{})=\delta (xx^{}),$$
(198)
from which upon integrating once, we get
$$\frac{v}{K}_xG(x,x^{})|_{x=x^{}0}^{x=x^{}+0}=1.$$
(199)
The solution to the differential equation is given by
$`G(x,x^{})`$ $`=`$ $`Ae^{|\overline{\omega }|(xx^{})/v},x<x^{}`$ (200)
$`=`$ $`Ae^{|\overline{\omega }|(xx^{})/v},x>x^{}.`$
Using this in Eq. (199), we see that
$`({\displaystyle \frac{2\overline{\omega }}{K}})A`$ $`=`$ $`1,`$ (201)
$`\mathrm{so}\mathrm{that}G_{\overline{\omega }}(x,x^{})`$ $`=`$ $`{\displaystyle \frac{K}{2\overline{\omega }}}e^{|\overline{\omega }|(xx^{})/v},x<x^{}`$ (202)
$`=`$ $`{\displaystyle \frac{K}{2\overline{\omega }}}e^{|\overline{\omega }|(xx^{})/v},x>x^{}`$
$`\mathrm{leading}\mathrm{to}\sigma _{\overline{\omega }}(x,x^{})`$ $`=`$ $`{\displaystyle \frac{Ke^2}{2\pi }}e^{|\overline{\omega }|(xx^{})/v},x<x^{}`$ (203)
$`=`$ $`{\displaystyle \frac{Ke^2}{2\pi }}e^{|\overline{\omega }|(xx^{})/v},x>x^{}.`$
The point to note is that in the $`\overline{\omega }0`$ limit or static limit, the conductivity is finite and does not drop down to zero even for large $`|xx^{}|`$. This is the main difference from macroscopic conductivities which always decay to zero as $`|xx^{}|\mathrm{}`$. Furthermore, for $`x=x^{}`$, even for arbitrary $`\overline{\omega }`$, the Green’s function has a finite value, which is responsible for the saturation value of the conductance. This only happens for a one-dimensional Green’s function. In any other dimension, the Green’s function and hence conductivity will be divergent at $`x=x^{}`$. Using this in the equation for the current, Eq. (191) for a static electric field $`\overline{E}_\omega (x)=2\pi \delta (\omega )E(x)`$, we finally get
$$I(x)=\frac{Ke^2}{2\pi }_0^L𝑑x^{}E(x^{})=\frac{Ke^2}{2\pi }(V_LV_0)$$
(204)
which gives the final result for the conductance as
$$g=\frac{Ke^2}{2\pi }.$$
(205)
There are several subtle points to note in this calculation. One is that we have taken the $`\omega 0`$ before $`|xx^{}|\mathrm{}`$, which is opposite to the usual order of limits in the Kubo formula. The physical justification for the usual order of limits in the Kubo formula comes from the fact that if we first take $`\omega `$ to zero, then we have a static electric field, which is periodic in space. This means that the charge will seek an equilibrium distribution after which there will be no flow of current. Setting $`|xx^{}|\mathrm{}`$, on the other hand, means taking the thermodynamic limit or infinite length limit first, which allows for an unlimited supply of electrons and is probably equivalent to having reservoirs even if we have do not really have infinite length wires. For the mesoscopic systems, however, it is not correct to take the thermodynamic limit first. The physical situation here, is that one applies a static electric field to a finite length of the wire $`L`$, which in fact, is comparable to the range of the conductivity. If we take the $`|xx^{}|\mathrm{}`$ limit first, then it is as if we are looking at a long length of the wire beyond the range of the conductivity. This is the usual limit and we will get the usual Ohm’s law, which however, is wrong in this context. In fact, it is instructive to try out the calculation with the other order of limits -$`i.e.`$ by computing $`\sigma _{\overline{\omega }}(q)`$ and taking the $`q0`$ limit first.
(H.W. Exercise 9: Try the above).
The second point is something we have mentioned earlier - $`i.e.`$, we have not taken contacts or leads into account. This was the initial computation by Kane and Fisher and they obtained the answer in Eq. (205) that the conductance of the clean wire depends on the interaction parameter $`K`$.
(b) Including leads
When any experiment is done, however, one does have explicit contacts or leads. In fact, when a measurement was actually done under conditions where one expected to measure the Luttinger parameter $`K`$, it was found to $`1`$, instead of 0.7, which was expected from other measurements of the $`K`$ value of the wire. (We will see how else $`K`$ can be measured after considering impurity scattering). So we need to understand what happens when we actually try to measure the conductance of a Luttinger wire.
How do we model the leads? The simplest model to consider is that the Luttinger wire is connected to Fermi liquid leads on either side. (See Fig. 6). So the regions A and C can be modelled by the same bosonic model with $`K_L=1`$ and the wire in region B can be modelled as before as a Luttinger wire with $`K=K`$. But now, we have to put appropriate boundary conditions at the points P and P’ between A and B and between B and C respectively. Note that we are making the assumption that one has the same $`\varphi `$ field or same quasiparticle in all the three regions and it is only the LL parameters which are changing. Although, it is interesting to compute the conductance in this case, it is still not clear that this brings the calculation any closer to real experiments, because real experiments will have three dimensional reservoirs.
We start with the action in all the three regions in Euclidean space as
$$S_E=\frac{1}{2}_0^\beta 𝑑\tau _0^L𝑑x[\frac{(_\tau \varphi )^2}{K(x)v(x)}+\frac{v(x)}{K(x)}(_x\varphi )^2],$$
(206)
with $`K(x)=K_L`$, $`v(x)=v_L`$ in regions A and C and $`K(x)=K`$, $`v(x)=v`$ in region B. This is just the free action of a scalar field in all the three regions. Fourier transforming the imaginary time variable with respect to $`\overline{\omega }`$, we obtain
$$S_E=\frac{1}{2}_0^\beta 𝑑\tau _0^L𝑑x[\frac{\overline{\omega }^2\varphi ^2}{K(x)v(x)}+\frac{v(x)}{K(x)}(_x\varphi )^2],$$
(207)
from which we see that the propagator satisfies the equation
$$\{_x(\frac{v(x)}{K(x)}_x)+\frac{\overline{\omega }^2}{K(x)v(x)}\}G_{\overline{\omega }}(x,x^{})=\delta (xx^{}).$$
(208)
Now let us consider the four regions. We assume that the interaction parameter changes abruptly at P and P’, but that the Green’s function is continuous and the derivative of the Green’s function has the correct discontinuity at all the boundaries. So now, we need to solve the Green’s function equation subject to these boundary conditions. Let us choose $`x^{}`$ to lie between $`0`$ and $`L`$. It is then easy to see that the solution is of the form
$`G_{\overline{\omega }}(x,x^{})=`$ $`Ae^{|\overline{\omega }|x/v}\mathrm{for}x0`$
$`=`$ $`Be^{|\overline{\omega }|x/v}+Ce^{|\overline{\omega }|x/v}\mathrm{for}0<xx^{}`$
$`=`$ $`De^{|\overline{\omega }|x/v}+Ee^{|\overline{\omega }|x/v}\mathrm{for}x^{}<xL`$
$`=`$ $`Fe^{|\overline{\omega }|x/v}\mathrm{for}x>L`$ (209)
for semi-infinite leads, because we have assumed that the lengths of the leads are sufficiently long compared to $`L`$ so that we do not need to put any further boundary conditions on them. Note that here the Green’s functions will no longer be functions of $`xx^{}`$ since we have explicitly broken translational invariance. The constants $`A,B,\mathrm{},F`$ are found by matching the boundary conditions. Since we are interested in the DC conductance, we only need the solutions for $`\overline{\omega }0`$ which are easy to obtain and are given by
$`A=F={\displaystyle \frac{K_L}{2\overline{\omega }}},B=E={\displaystyle \frac{K_L+K}{4\overline{\omega }}},C=D={\displaystyle \frac{K_LK}{4\overline{\omega }}}.`$ (210)
From this, we see that $`\sigma _{\overline{\omega }}(x,x^{})`$ is $`x`$ and $`x^{}`$ independent in the $`\overline{\omega }0`$ limit and is equal to $`K_Le^2/2\pi `$ in all regions from which we find the conductance (using Eq. (191)) given by
$$g\frac{I}{V}=\frac{K_Le^2}{2\pi }.$$
(211)
is the same in all the regions. Thus, the conductance is determined by the $`K_L`$ of the leads, which is just $`K_L=1`$ for Fermi liquid leads and does not depend on interactions in the wire. This is a highly counter-intuitive answer! It is telling us that whether we measure the conductance in the leads or in the quantum wire, we get the same answer, so long as we take into account the fact that we are attaching leads, which allow for the fermions to enter and leave the quantum wire. At a very naive level, one may understand this by saying that since the wire itself has no impurities, the only source of resistance is the contact effect between the leads and the wire, which has nothing to to do with the interactions in the wire. However, remember that we have taken semi-infinite leads and abrupt contacts and we are only looking for DC conductance. If any of these assumptions are relaxed, certainly, there are differences in the three regions and one could get more interesting answers.
In fact, using a Landauer-Buttiker scattering approach , there has been some attempt to understand these results more intuitively.
Physically, the difference between this computation and the earlier one is that any real measurement requires Fermi liquid leads. So the end result is that the measurable conductance of an interacting one-dimensional wire is simply given by $`g=e^2/h`$ for spinless fermions and $`g=2e^2/h`$ for fermions with spin .
## 9 Transport in the presence of isolated impurities
Computing conductance with a single impurity
Now let us consider the case when there is a single impurity at the origin. At first, we will model the impurity as a weak barrier and add a term to the action of the form
$$S_{\mathrm{int}}=𝑑x𝑑\tau V(x)\psi ^{}(x)\psi (x).$$
(212)
We assume that $`V(x)`$ is weak and is centred around the origin. For instance, we can choose $`V(x)=\lambda \delta (x)`$, where $`\lambda `$ is much less than the Fermi energy.
First, let us think of what happens when we introduce such a perturbation in a non-interacting wire. In that case, all one has is a one-dimensional quantum mechanics problem with a $`\delta `$-function potential at the origin. We can find the reflection and transmission probabilities for a single particle with momentum $`k`$ as
$$R=\frac{\lambda ^2}{\lambda ^2+k^2}\mathrm{and}T=\frac{k^2}{\lambda ^2+k^2}.$$
(213)
So for any $`\lambda `$, one gets both reflection and transmission. To get the total current, we just have to sum up the contributions of all the electrons close to the Fermi surface. But it is clear that there will be non-zero conductance for any potential, with the amount of current being transmitted depending on the strength of the potential. However, for the Luttinger wire, since there exists interactions between electrons in the wire, and no convenient quasiparticle picture, one cannot solve the problem this way. We have to use the bosonized field theory and include the impurity potential as a perturbation.
Let us first rewrite the impurity potential in terms of the left- and right- moving low energy Dirac modes. We find that
$`\psi ^{}(x)\psi (x)`$ $`=`$ $`(\psi _R^{}e^{ik_Fx}+\psi _L^{}e^{ik_Fx})(\psi _Re^{ik_Fx}+\psi _Le^{ik_Fx})`$
$`=`$ $`\psi _R^{}\psi _R+\psi _L^{}\psi _L+e^{i2k_Fx}\psi _R^{}\psi _L+e^{i2k_Fx}\psi _L^{}\psi _R`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{\pi }}}_x\varphi +{\displaystyle \frac{1}{2\pi \alpha }}(\eta _R^{}\eta _Le^{i2\sqrt{\pi }(\varphi _R+\varphi _L+2k_Fx)}+\eta _L^{}\eta _Re^{i2\sqrt{\pi }(\varphi _R+\varphi _L+2k_Fx)}),`$
where the last line is obtained using standard bosonization. So the full action is given by
$$S=S_E+S_{\mathrm{int}}=S_E\frac{\lambda }{\sqrt{\pi }}_x\varphi (0)+\frac{\lambda }{2\pi \alpha }𝑑\tau \mathrm{cos}2\sqrt{\pi }[\varphi _R(0)+\varphi _L(0)],$$
(215)
where $`S_E`$ is given in Eq. (193) and we have incorporated the fact that the potential only acts at the origin. Moreover, we have simply set both $`\eta _R^{}\eta _L`$ and $`\eta _L^{}\eta _R`$ to be one, with the knowledge that in correlation functions, we will compute $`O(\tau )O^{}(0)`$ so that the Klein factors disappear using $`\eta _{R/L}\eta _{R/L}^{}=\eta _{R/L}^{}\eta _{R/L}=1`$. The first term due to the interaction can be taken care of by a simple redefinition of $`_x\varphi _x\varphi ^{}=_x\varphi +\lambda /2\sqrt{\pi }`$, which makes no difference to the conductance. This could have been seen even from the fermion terms from which it came. The $`\psi _R^{}\psi _R+\psi _L^{}\psi _L`$ term only causes scattering at the same Fermi point with momentum transfers $`q<<2k_F`$. This does not change the direction of propagation of the particles and hence does not affect conductance in any appreciable way. But the cosine term, on the other hand, occurs because of backscattering of fermions from the origin. These represent scattering with $`q|2k_F|`$ \- $`i.e.`$, from the left branch to the right branch and vice-versa and change the direction of propagation of the particles. These scatterings will definitely affect the conductance. The action with this perturbation is no longer quadratic and cannot be exactly solved. However, since $`\lambda `$ is a weak perturbation, one can try to use perturbation theory and the renormalization group approach to see the relevance of this perturbation at low energies.
What is the question that we want to answer? We want to compute the conductance through this barrier at low energies. One way to do that is to see whether this barrier coupling strength grows or becomes smaller as we go to lower energy scales. To check that, we need to perform the usual steps of a renormalization group analysis.
Here since the perturbation term is fixed in space, it is more convenient to first integrate out the variables away from the origin and write down the action purely in terms of the $`\varphi (x=0,\tau )`$ variables. Since integrating out quadratic degrees of freedom is equivalent to using equations of motion for those degrees of freedom, we write down the equations of motion for the action $`S_0`$ as
$$_x^2\varphi \frac{\overline{\omega }^2}{v^2}\varphi =0_x^2\varphi k^2\varphi =0.$$
(216)
The solution to the above equations are given by
$`\varphi `$ $`=`$ $`Ae^{|k|x},x>0`$ (217)
$`=`$ $`Ae^{|k|x},x<0,`$
where $`A\varphi (x=0,\tau )`$. Using this solution in the action, we get the effective action in terms of $`\varphi (\overline{\omega })=\varphi (x=0,\tau )e^{i\overline{\omega }\tau }𝑑\tau `$ as
$`S_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{1}{2K}}{\displaystyle \frac{d\overline{\omega }}{2\pi }_{\mathrm{}}^0𝑑x[v\varphi ^2k^2e^{2kx}+\frac{\overline{\omega }^2}{v}e^{2kx}]}+{\displaystyle \frac{1}{2K}}{\displaystyle \frac{d\overline{\omega }}{2\pi }_0^{\mathrm{}}𝑑x[v\varphi ^2k^2e^{2kx}+\frac{\overline{\omega }^2}{v}e^{2kx}]}`$ (218)
$`=`$ $`{\displaystyle \frac{1}{2K}}{\displaystyle \frac{d\overline{\omega }}{2\pi }\frac{2\varphi ^2\overline{\omega }^2}{v}[\frac{e^{2kx}}{2k}|_{\mathrm{}}^0+\frac{e^{2kx}}{2k}|_0^{\mathrm{}}]}`$
$`=`$ $`{\displaystyle \frac{1}{K}}{\displaystyle \frac{d\overline{\omega }}{2\pi }|\overline{\omega }|\varphi ^2}`$
using $`k=|\overline{\omega }|/v`$. Notice the singular dependence on the Matsubara frequency $`|\overline{\omega }|`$. The reason for its appearance is the following. In real space, even for a quadratic action, all degrees of freedom (dof) are coupled. ( It is only in Fourier space that every mode is decoupled). So when we integrate out all dof except the one at the origin, the dispersion of this degree of freedom can change and has changed. This is why we get the modulus factor in the effective action. When we Fourier transform back to imaginary time, we get
$`{\displaystyle \underset{\overline{\omega }}{}}|\overline{\omega }|\varphi _n^2`$ $``$ $`i{\displaystyle \frac{d\omega }{2\pi }𝑑\tau 𝑑\tau ^{}e^{i\omega (\tau \tau ^{})}|\omega _n|\varphi ^{}(\tau )\varphi (\tau ^{})}`$ (219)
$`=`$ $`{\displaystyle 𝑑\tau 𝑑\tau ^{}\frac{2}{(\tau \tau ^{})^2}\varphi ^{}(\tau )\varphi (\tau ^{})},`$
$`i.e.`$, an explicitly non-local interaction in imaginary time.
So now, we have an action solely in terms of the variables at the origin with the action given by
$$S=\frac{1}{2K}\frac{d\overline{\omega }}{2\pi }|\overline{\omega }|\varphi (\overline{\omega })^2+\lambda \frac{d\overline{\omega }}{2\pi }\mathrm{cos}[2\sqrt{\pi }\varphi (\overline{\omega })].$$
(220)
The RG analysis now involves finding out how the coefficient $`\lambda `$ behaves as we go to lower and lower energies. Before we perform the RG analysis, we may ask why would we want to go to lower energy scales? The general idea is that in spite of the fact that in different physical problems or models, the parameter $`\lambda `$ may be slightly different, qualitatively many such models may have the same behavior. This is because they are all governed by the same fixed point Hamiltonian with the fixed point Hamiltonian being defined as the Hamiltonian one gets when the RG flow stops. So the aim is to keep reducing the energy scale till the RG flow stops so that we can find out the appropriate fixed point Hamiltonian for this model.
In this problem, we want to find out whether the fixed point Hamiltonian has a large barrier or a small barrier. To find that out, let us perform the three steps of the renormalization group transformation. We choose a high frequency cutoff $`\mathrm{\Lambda }`$, which is the real physical cutoff of the theory. Then we rescale $`\mathrm{\Lambda }\mathrm{\Lambda }/s`$ with $`s>1`$ and then divide $`\varphi (\overline{\omega })`$ into $`\varphi _<(\overline{\omega })`$ (slow modes) and $`\varphi _>(\overline{\omega })`$ (fast modes) for the modes with frequencies less than or greater than the cutoff $`\mathrm{\Lambda }/s`$ respectively. Finally, we integrate out the fast modes, which are the modes between $`\mathrm{\Lambda }/s`$ and $`\mathrm{\Lambda }`$ and rescale $`\overline{\omega }\overline{\omega }^{}=s\overline{\omega }`$ or $`\tau \tau ^{}=\tau /s`$ to get back to the original range of integrations. To lowest order, (tree level contribution) we find that
$$\lambda 𝑑\tau \mathrm{cos}2\sqrt{\pi }\varphi _<(x=0,\tau )\lambda s^{1d}𝑑\tau \mathrm{cos}2\sqrt{\pi }\varphi _<(x=0,\tau ),$$
(221)
where $`d`$ is the dimension of the cosine operator. This was explicitly computed earlier and we found that $`d=K`$. The RG equation is now easily obtained by taking $`s=1+dl`$, for infinitesimal $`dl`$. We find that the new $`\lambda ^{}`$ after rescaling is given by
$`\lambda ^{}=\lambda (1+dl)^{1K}`$
$``$ $`\lambda ^{}\lambda =(1K)\lambda dl`$
$``$ $`{\displaystyle \frac{d\lambda }{dl}}=(1K)\lambda .`$ (222)
Normally, one would have had coupled RG equations for $`\lambda `$ and $`K`$. But here since $`1/K`$ is the coefficient of a singular operator, it does not get renormalized to any order.
Notice that in Eq. (221), the coefficient of the operator gets rescaled by a factor $`s^{1d}`$ rather than $`s^{2d_i}`$ as we had mentioned earlier when we computed the dimension of the cosine operator in the spin model. The difference is that the operator in the spin model, in the action, required integration over both space and time. So we rescaled both the space and time ( or equivalently both the momentum and the energy). However, in this case, the operator exists only at a fixed space point. So we only need to integrate over the time coordinate. Hence, the naive scaling dimension or engineering dimension of the operator is 1 and not 2. Such operators are called boundary operators. You will learn more about them in the course on boundary conformal field theory.
The RG equation is now trivial to analyze. For any $`K>1`$, (which corresponds to attractive interactions between the electrons), the $`\lambda `$ renormalizes to zero and for any $`K<1`$, (corresponding to repulsive interactions), it grows stronger and stronger. In other words, for $`K>1`$, the fixed point Hamiltonian is just the free boson Hamiltonian with no barrier and for $`K<1`$, the fixed point Hamiltonian has two disconnected wires to the left and right of the origin. For $`K=1`$, which is the limit of no interactions in the fermionic model, the coupling is marginal. (This was expected, since we know that for free fermions, both transmission and reflection occurs depending on the strength of the barrier potential). Thus, for attractive interactions, the barrier renormalizes to zero and the wire is ‘healed’, whereas for repulsive interactions, the barrier renormalizes to infinity and the wire is ‘cut’. Note that both these answers are completely independent of the strength of the barrier potential .
Strong barrier limit
Since, we are doing perturbation theory, we cannot assume that this result holds for arbitrary $`\lambda `$. It is strictly valid only for $`\lambda 0`$. Once $`\lambda 1`$, the perturbative analysis in $`\lambda `$ breaks down. So it is worthwhile to try and see what happens in the other limit. Supposing we start with two decoupled wires and then allow a small hopping between the two wires. Will this hopping grow at low energies and heal the wire or will it renormalize to zero?
Here, we start with two semi-infinite Luttinger liquid wires and analyze the effect of adding a small hopping term coupling the two wires at $`x=0`$. The models for $`x<0`$ and $`x>0`$ are given by the action
$$S_E=\frac{1}{2}_0^\beta 𝑑\tau _0^L𝑑x[\frac{(_\tau \varphi _i)^2}{K(x)v(x)}+\frac{v(x)}{K(x)}(_x\varphi _i)^2]$$
(223)
for $`i=<`$ and $`i=>`$ respectively. We can also write it in terms of the dual variables as
$$S_E=\frac{1}{2}_0^\beta 𝑑\tau _0^L𝑑x[\frac{K(x)}{v(x)}(_\tau \theta _i)^2+K(x)v(x)(_x\theta _i)^2].$$
(224)
Note that in terms of the dual variables, the action has $`1/K`$ in position of $`K`$. This is because the roles of the fields and the canonically conjugate momenta have interchanged. The fact that the wire is cut implies that at the point $`x=0`$, there is zero density of either $`<`$ or $`>`$ particles - $`\psi _<^{}\psi _<(x=0)=0`$ and $`\psi _>^{}\psi _>(x=0)=0`$. In the bosonic language, this is imposed as $`2\sqrt{\pi }\varphi _<(x=0)=2\sqrt{\pi }\varphi _>(x=0)=\pi /2`$ (and also $`_x\varphi (x=0)=0`$ as can be seen from Eq. (LABEL:imp)). Now a term which hops an electron from one wire to another in the Hamiltonian is just
$`H`$ $`=`$ $`t[\psi _<^{}\psi _>+h.c.]`$ (225)
$`=`$ $`t[\psi _{R<}^{}\psi _{R>}+\psi _{L<}^{}\psi _{L>}+\psi _{R>}^{}\psi _{R<}+\psi _{L>}^{}\psi _{L<}`$
$`+\psi _{L<}^{}\psi _{R>}+\psi _{L>}^{}\psi _{R<}+\psi _{R<}^{}\psi _{L>}+\psi _{R>}^{}\psi _{L<}],`$
where the second equation involves the left and right moving fields and we have already set $`x=0`$. Here, again, the terms that involve fields at one Fermi point are low energy forward scattering terms which do not affect the conductance. In terms of the bosonic fields too, they can be taken care of by trivial redefinitions. But the intra-Fermi point scatterings which will affect the conductance can be bosonized and written in the action as
$`\delta S=t{\displaystyle }d\tau [`$ $`\eta _{L<}^{}\eta _{R>}e^{i(\varphi _{L<}+\varphi _{R>})}+\eta _{L>}^{}\eta _{R<}e^{i(\varphi _{L>}+\varphi _{R<})}`$ (226)
$`+`$ $`\eta _{R<}^{}\eta _{L>}e^{i(\varphi _{R<}+\varphi _{L>})}+\eta _{R>}^{}\eta _{L<}e^{i(\varphi _{R>}+\varphi _{L<})}].`$
Now, we impose the boundary condition on the bosonic fields that we mentioned above, which constrains $`\varphi (0)=\varphi _R(0)+\varphi _L(0)`$ to be equal to $`\pi /2`$. Using this, we can express the above equation solely in terms of the $`\varphi _L(0)\varphi _R(0)=\theta (0)`$ fields and get
$$\delta S=4t𝑑\tau \mathrm{cos}(\theta _>\theta _<),$$
(227)
where once again, we have been able to drop the Klein factors after checking that they do not lead to any extra minus signs in the correlation functions. (Physically, the reason why we only get the $`\theta _i`$ term at the origin is because the constraint has set $`\varphi _i(x=0)=\pi /2`$). Computing the dimension of this operator, we see that to leading order, the RG equations are given by
$$\frac{dt}{dl}=(1\frac{1}{K})t.$$
(228)
($`K`$ has been replaced by $`1/K`$ because we now have to compute the dimensions in the dual action). Thus, for repulsive interactions ($`K<1`$), the hopping term is irrelevant and flows to zero. This confirms the weak barrier calculation that the wire is insulating. On the other hand, for attractive interactions, the hopping strength grows, ultimately healing the wire. This again is in accordance with the weak coupling analysis.
Intermediate fixed points?
We have started from a wire with a weak barrier and shown that under repulsive interactions, the barrier strength grows. We have also started from two decoupled wires and shown that for repulsive interactions, any small hopping term renormalizes to zero. Hence, it seems plausible to conclude that for repulsive interactions in the wire, any barrier will cut the wire and the conductance goes to zero. However, one should keep in mind that our analysis is strictly true only for $`\lambda ,t0`$. Hence, it could happen that for intermediate values of the barrier strength, one could have a pair of non-trivial fixed points (see Fig. 7).
Conductance at finite voltage and temperature
The earlier analysis only tells us how the barrier strength or the tunneling amplitude grows or falls as we go to low energies. But instead of allowing the energy scale to become arbitrarily low, we can cut off the energy scale of renormalization at some finite energy scale, which could be the temperature $`T`$ or the voltage $`V`$. Note that the energy scale at which we want to cutoff the integral is related to the initial high energy scale at which we start the RG as $`E=E_0e^l`$. So for attractive interactions for which weak barriers are irrelevant and for which one would expect perfect transmission at very low energies will have power law corrections when we put the lower energy cutoff as $`E`$. In that case, we have
$$_{\mathrm{\Lambda }_0}^\mathrm{\Lambda }\frac{d\lambda }{\lambda }=_0^{ln(E_0/E)}𝑑l(1K),$$
(229)
which means that the effective barrier strength $`\mathrm{\Lambda }`$ is proportional to $`\mathrm{\Lambda }_0(E/E_0)^{K1}`$. So by choosing $`E=T,V`$, we see that one can get power law corrections to the naive conductance at $`T0,V0`$. In other words, if we measure the conductance at a finite temperature $`T`$, rather than at $`T=0`$, instead of zero conductance for $`K<1`$, we will get conductances which go as $`T^{1K}`$ (roughly the inverse of the barrier strength). Similarly, if instead of measuring conductances as $`V0`$, we measure them at finite voltages, we find that the conductances go as $`V^{1K}`$. On the other hand, for repulsive interactions, we need to start at the strong coupling limit with two decoupled Luttinger wires and allow for a small hopping, which is irrelevant in the RG sense. Here, again, if we cutoff the lower energy scale at $`E`$, we expect instead of zero transmission, power law corrections of the form $`IV^{11/K}`$ and $`IT^{11/K}`$. The only difference in the analysis at the strong coupling fixed point and the weak coupling fixed point is that $`K`$ gets replaced by $`1/K`$ as we saw in the RG equations. This, in fact, is one way in which $`K`$ can be measured in experiments. They could explicitly make a constriction in the quantum wire and measure conductances through it and extract $`K`$.
## 10 Concluding Remarks
Almost any interacting quantum system in one dimension which is gapless and has a linear dispersion for the low-energy excitations can be described as a Luttinger liquid at low energies and long wavelengths. As we have seen, the properties of a Luttinger liquid are determined by the two parameters $`v`$ and $`K`$. These in turn depend on the various parameters which appear in the microscopic Hamiltonian of the system. Some examples of systems where Luttinger liquid theory and bosonization can be applied are quantum spin chains (including some spin ladders), quasi-one-dimensional organic conductors and quantum wires (with or without impurities), edge states in a fractional quantum Hall system, and the Kondo problem. Some of these examples have been discussed above.
Antiferromagnetic spin-$`1/2`$ chains have a long history going back to their exact solution by the Bethe ansatz. In recent years, many experimental systems have been studied which are well-described by quasi-one-dimensional half-odd-integer spin models with isotropic (Heisenberg) interactions. Such systems behave at low energies as a $`K=1/2`$ Luttinger liquid with an $`SU(2)`$ symmetry. It seems to be difficult to vary $`K`$ experimentally in spin systems. In contrast, a single-channel quantum wire (which is basically a system of interacting electrons which are constrained to move along one particular direction) typically has two low-energy sectors, both of which are Luttinger liquids (except at special densities like half-filling). One of these is the spin sector which again has $`K=1/2`$. The other one is the charge sector whose $`K`$ value depends on a smooth way on the different interactions present in the system. Finally, the edge states in a fractional quantum Hall system behave as a chiral Luttinger liquid with $`K`$ taking certain discrete rational values; the value of $`K`$ can be changed by altering the electron density and the magnetic field in the bulk of the system. For all these systems, many properties have been measured such as the response to external electric and magnetic fields (conductivity or susceptibility) and to disorder, scattering of neutrons or photons from these systems, and specific heat; so the two Luttinger parameters can be extracted from the experimental data. The measurements clearly indicate the Luttinger liquid-like behavior of these systems with various critical exponents depending in a non-universal way on the interactions in the system.
On the theoretical side, a large number of exactly solvable models in one dimension have been shown to behave as Luttinger liquids at low energies . These include
(i) models with short range interactions which are solvable by the Bethe ansatz, such as the $`XXZ`$ spin-$`1/2`$ chain (where $`K`$ can take a range of values from $`1/2`$ to $`\mathrm{}`$; this includes the $`XY`$ model with $`K=1`$ and the isotropic antiferromagnet with $`K=1/2`$ as special cases), and the repulsive $`\delta `$-function Bose gas (where $`K`$ can go from $`1`$ in the limit of infinite repulsion to $`\mathrm{}`$ in the limit of zero repulsion), and
(ii) models with inverse-square interactions such as the Calogero-Sutherland model (where $`K`$ can go from $`0`$ to $`\mathrm{}`$) and the Haldane-Shastry spin-$`1/2`$ model (where $`K=1/2`$).
The models of type (ii) are ideal Luttinger liquids in the sense that they are scale invariant; the coefficients of all the marginal operators vanish, and therefore their correlation functions and excitation energies contain no logarithmic corrections. This property makes it particularly easy to study these systems numerically since the asymptotic behaviors are reached even for fairly small system sizes.
What has been left out?
Finally, let us mention the various important things in this field which has been left out. We have only worked with spinless fermions in the transport analysis. When we include spin and do not destroy the $`SU(2)`$ spin symmetry of the system, the results are very similar to the spinless fermion case. For repulsive interactions, the barrier becomes infinite and for attractive interaction, the barrier is healed. However, when the $`SU(2)`$ symmetry is destroyed, there exists possibilities of intermediate (non-trivial) fixed points where either spin or charge can be transmitted and the other reflected. The other thing that has been left out is the phenomenon of resonant tunneling with two impurities. This is an interesting result, because it says that for repulsive interactions, a single impurity cuts the wire, but with two impurities, one can have particular energies, where there can be transmission. The reason, of course, is quantum mechanical tunneling. Here, the energy levels, are the energy levels of the quantum dot that is formed by the two impurities and one can have resonant tunneling at these energy levels. If we include interactions between the electrons on the island, (which is naturally included in the bosonized formalism), we can obtain the physics of the Coulomb blockade. The other important thing that we have left out, from a physical point of view, is what happens if there is a finite density of random impurities. In general, one would expect Anderson localization and no transport. But there are regimes of delocalization as well in the phase diagram. Finally, a very important application where the physics of the Luttinger liquids has actually been experimentally seen is in the edge states of the fractional Quantum Hall fluid. Since here, the edge states are chiral, a lot of the complications of backscattering due to impurities are avoided and it is possible to explicitly construct constrictions and allow tunneling through them. Here, both at the theoretical and experimental level, there are a lot of beautiful results that are worth understanding.
Another important topic not covered here is non-abelian bosonization . This is a powerful technique for studying one-dimensional quantum systems with a continuous global symmetry such as $`SU(2)`$. For instance, isotropic Heisenberg antiferromagnets and Kondo systems are invariant under spin rotations, and they can be studied more efficiently using non-abelian bosonization.
To conclude, let us just say that low dimensional systems and mesoscopic systems have gained in importance in the last few years. Although currently, much of the theoretical work in mesoscopic systems has only involved conventional Fermi liquid theories, it is clear that there are regimes where strong interactions are very important. We expect that bosonization will be one of the important non-perturbative tools to analyze such problems for a few more years to come.
Acknowledgments
DS thanks J. Srivatsava for making Figures 1 to 5.
Figure Captions
1. One-particle momentum distribution function. (a) shows the finite discontinuity at the Fermi momentum $`k_F`$ for a system of interacting fermions in more than one dimension. (b) shows the absence of a discontinuity in an interacting system in one dimension.
2. Picture of the Fermi sea of a lattice model; the momentum lies in the range $`[\pi ,\pi ]`$. The occupied states (filled circles) below the Fermi energy $`E_F=0`$ and the two Fermi points at momenta $`\pm k_F`$ are shown.
3. The one-particle states of a right-moving fermion showing the occupied states (filled circles) below zero energy and the unoccupied states above zero energy.
4. Two possible particle-hole excitations of a right-moving fermionic system showing the occupied states.
5. The one-particle states of a left-moving fermion showing the occupied states below zero energy and unoccupied states above zero energy. Note that the momentum label $`k`$ increases towards the left.
6. The single channel quantum wire with Fermi liquid leads on the left and the right.
7. Renormalization group flow diagram for a quantum wire with repulsive interactions in the presence of an impurity or barrier. In the absence of any non-trivial fixed points, the stable fixed point is the strong coupling fixed point. But perturbative analyses at the strong and weak coupling fixed points cannot rule out a pair of non-trivial fixed points at intermediate strengths of the barrier potential.
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# The Nuclear Dynamics of M32. I.
## 1 Introduction
The presence of a supermassive compact object, presumably a black hole, at the center of the dwarf elliptical galaxy M32 has been suspected for some time (Tonry (1987)). The evidence consists of rapid rotation of the stars near the center of M32 and a central peak in the stellar velocity dispersions (Tonry (1987); Dressler & Richstone (1988); Carter & Jenkins (1993); van der Marel et al. 1994a ; Bender, Kormendy & Dehnen (1996)). The most recent study (van der Marel et al. (1997), 1998) used data from the Faint Object Spectrograph (FOS) on the Hubble Space Telescope (HST) to infer the rotation and dispersion velocities with a spatial resolution of $`0.1^{\prime \prime }`$ in the inner $`0.5^{\prime \prime }`$ of M32. The FOS data revealed a sharper rise in the stellar velocity dispersions than had been observed from the ground; however the velocity resolution of the FOS is limited, making that instrument only marginally useful for the study of a low velocity dispersion system like M32. In fact, van der Marel et al. (1998) found large point-to-point variations in their velocity dispersion measurements, making the dynamical interpretation uncertain. Nevertheless, the case for a supermassive black hole in M32, of mass $`M_h3\times 10^6_{}`$, was considerably strengthened.
Here we present observations of M32 carried out using STIS, the Space Telescope Imaging Spectrograph, on HST. Our data were obtained as part of the STIS Investigation Definition Team’s (IDT) key program to observe the nuclei of a sample of $`15`$ nearby galaxies in the spectral region centered on the Calcium triplet, $`\lambda 8600`$Å. This paper is the first in a series that will present stellar-kinematical evidence from STIS for the presence (or absence) of supermassive black holes in galactic nuclei.
The STIS data improve on earlier ground-based and FOS data from M32 in several ways. The spatial resolution of STIS is $`0.1^{\prime \prime }`$, or $`0.3`$ pc at the distance of M32, similar to that of the single FOS aperture; however STIS provides continuous spatial sampling along a slit. The spectral resolution of STIS in the G750M mode ($`38`$ km/s) is much greater than that of the FOS making STIS a more suitable instrument for observing M32, whose velocity dispersion outside of the nucleus is only $`60`$ km s<sup>-1</sup>. We were able to obtain from the STIS spectra not only the lowest-order moments of the stellar velocity distribution – the rotation velocity and velocity dispersion – but also the full line-of-sight velocity distribution (LOSVD) as a function of position along the major axis.
A number of modelling studies (van der Marel et al. 1994b ; Dehnen (1995); Qian et al. (1995)) have made predictions about the observable signatures at HST resolution of a supermassive black hole in M32. The black hole is expected to be associated with three kinematical features. (1) The stellar rotation velocity should remain flat, or rise slightly, into the central resolution element. (2) The stellar velocity dispersion should exhibit a sudden upturn at a distance of $`0.2^{\prime \prime }0.5^{\prime \prime }`$ from the center, reaching a central value of $`120`$ km s<sup>-1</sup> or greater depending on the mass of the black hole. (3) The distribution of line-of-sight stellar velocities in the resolution element centered on the black hole should be strongly non-Gaussian, with extended, high-velocity wings. Part of the predicted rise in the velocity dispersion near the center would be due to these wings; part to blending of the rotation curve from the two sides of the galaxy; and part to an intrinsic rise in the random velocities.
We see all three signatures of the black hole in the STIS data. The velocity dispersion spike is most impressive; in terms of the usual parametrization $`\sigma _0`$ (the dispersion of the Gaussian core of the LOSVD), the central measured value is $`130`$ km s<sup>-1</sup>, while the value of $`\sigma `$ corrected for the non-Gaussian wings of the LOSVD is considerably greater, at least $`175`$ km s<sup>-1</sup>, or $`3`$ times the value in the main body of M32.
A detailed description of the observations and the data reduction is given in §2. In §3 we describe the methods used to extract the stellar LOSVDs from the STIS spectra. We carried out independent analyses based on two spectral deconvolution routines: the FCQ algorithm of R. Bender (1990), which is a Fourier method; and the MPL routine of D. Merritt (1997), which is based on nonparametric function-estimation techniques. The two algorithms gave consistent results for the low-order velocity moments of M32 (the mean and dispersion) but there were systematic differences in the recovered values of $`h_4`$, the parameter that measures symmetric deviations of the LOSVD from a Gaussian. We argue that the $`h_4`$ values are consistent after taking into account the biases in the two methods; if our analysis is correct, the true values of $`h_4`$ in M32 are significantly greater than zero both within the region where the black hole’s force is dominant and at larger radii.
A preliminary interpretation of the M32 kinematical data is presented in §4, where we address two basic questions: Have we resolved the central spike in the velocity dispersions? and: Is the shape of the central LOSVD consistent with what is expected for a nucleus containing a black hole? We answer “yes” to both questions, although the resolution of the central spike may be marginal if the black hole mass is near the lower limit of the allowed range. In any case, M32 is the first galaxy for which the imprint of the black hole’s gravitation on the stellar velocities has been observed with a resolution comparable to that of the gas-dynamical studies (e.g. Bower et al. (1998); Ferrarese & Ford (1999)).
We show that both the velocity dispersion spike and the wings of the central LOSVD are consistent with the predictions of simple dynamical models containing black holes, with masses $`M_h`$ in the range $`2\times 10^6_{}\begin{array}{c}<\hfill \\ \hfill \end{array}M_h\begin{array}{c}<\hfill \\ \hfill \end{array}5\times 10^6_{}`$. More accurate estimates of $`M_h`$ will be presented in Paper II where the full kinematical data set will be compared to three-integral axisymmetric models.
We adopt a distance to M32 of $`0.7`$ Mpc; thus $`1^{\prime \prime }`$ corresponds to 3.25 pc.
## 2 Observations and Data Reduction
M32 was observed on 1998 September 04 with STIS in the long-slit mode with wavelength centered on the Ca II triplet feature near 8561Å. STIS is described by Woodgate et al. (1998) and its on-orbit performance by Kimble et al. (1998). The present data are part of a survey of the nuclei of nearby galaxies being conducted by the STIS IDT (HST Program ID: 7566). The goal of the survey is to place stellar kinematical constraints on the masses of nuclear black holes. Two orbits of data with a total integration time of 4,898 seconds were obtained. The HST spacecraft tracking was operated in fine lock with a reported jitter of no more than 3 mas rms or 11 mas peak-to-peak. The aperture was $`52^{\prime \prime }\times 0.1^{\prime \prime }`$ with a position angle of 163, coincident with the M32 isophotal major axis. The CCD data were read out in the unbinned mode. Spatial sampling at the focal plane was at every $`0.05071^{\prime \prime }`$, corresponding to a 2-pixel optical resolution of about $`0.115^{\prime \prime }`$ FWHM. Outside of $`\pm 0.7^{\prime \prime }`$, the data were binned spatially to enhance the signal. The spectral resolution was approximately 38 km s<sup>-1</sup>. M32 spectra were obtained in the CR-split mode to assist with cosmic ray (CR) identification and rejection. The location of the galaxy center in the image was moved by approximately 4.5 rows along the aperture between the two orbits to ensure that residual detector sensitivity variations, that may not be completely removed from the data during reduction, are not mistaken for weak features. This form of dithering also assists with the identification of hot pixels in the CCD which do not rectify well.
A spectrum of the star HR 7615 was obtained with the same STIS aperture earlier in the program (Figure The Nuclear Dynamics of M32. I.). This bright, K0 III giant was used as a template when inferring the LOSVD of the stars in M32 from the STIS spectra. Template spectra of a handful of other bright, cool stars of various spectral types were also obtained and used to test the sensitivity of the spectral deconvolution to template mismatch. A set of spectral images of HR 7615 were taken centered at $`0.05^{\prime \prime },0.00^{\prime \prime }`$ and $`+0.05^{\prime \prime }`$ with respect to the centerline of the aperture (i.e. offset along the dispersion direction). These data were added together using appropriate weights to match the aperture illumination profile of M32. Each of these spectral images has a slightly different shift in its velocity centroid and the combined image provides a more accurate template for determining the kinematics of an extended object such as a galaxy. Internal wavelength calibration images (“Wavecals”) as well as an internal continuum lamp image (“Flat field”) were taken for calibration purposes. An image that has had all of the instrumental response removed is said be a rectified image. The contemporaneous flat field spectrum was obtained in the portion of the orbit where HR 7615 was behind the earth. The flat spectrum was taken through the $`0.2^{\prime \prime }\times 0.09^{\prime \prime }`$ aperture rather than through the long aperture since the former is superior for removing the pixel-to-pixel detector response of a stellar point source. A contemporaneous flat field was used to remove the internal fringing which is significant for wavelengths greater than 7500 Å and which changes over time (Goudfrooij, Baum & Walsh (1997)). We note that the fringing is far more serious for low dispersion spectra especially for high S/N than it is for our data taken in medium dispersion.
Cosmic rays (CRs) account for approximately 20% of the total signal and contaminate approximately 5% of the pixels in a typical exposure. CRs were identified and removed using the following procedure. The centroids in the cross-dispersion direction were determined for each rectified image and the images shifted so that the galaxy core appeared on the same row. CRs were identified and removed by comparing the flux in a given pixel to the flux in the corresponding pixel in subsequent images. For each pixel, outlying values were rejected and excluded when the frames were averaged together. Our data set for M32 included 4 raw images. Most pixels were found to have 4 frames contributing to their average values, while fewer pixels had 3 or less frames. Only 2 pixels within the central $`2^{\prime \prime }`$ had contributions from no frames. Those pixels were assigned values representing the average of their adjacent pixels.
The M32 data were reduced using two separate approaches: 1) “Shift and Add,” and 2) “Frame by Frame.” The latter relies heavily on a standard software package called CALSTIS at Goddard Space Flight Center. The “Shift and Add” method starts by removing the detector response using contemporaneous flat, bias, and dark calibration files. CR hits are then removed using the procedure described previously. The frames are averaged together and the resulting frame is remapped to place the spectra from a single location along the aperture onto a single row. As with most spectrographs, STIS produces spectra with S-shaped and pincushion distortions as well as spectra that are not aligned exactly with a row. A cubic interpolation was used to remap the spectra for later analysis. The remapping is not perfect, with centroid errors of approximately 0.1 pixel rms. This level of accuracy was deemed adequate, although we are working to improve it. In addition, the remapping produces very minor ($`\begin{array}{c}<\hfill \\ \hfill \end{array}1\%`$) residual Moiré ripples in the data, which can be minimized but not completely eliminated. Further work is also underway to measure and correct for this residual affect. The strength of the “Shift and Add” technique is that it preserves photometric accuracy. However it has the disadvantage of introducing a subpixel image smearing since each frame is registered to the nearest integer pixel.
The philosophy of the “Frame by Frame” method is to apply all calibration corrections (including the remapping described above) to a single frame before the resulting frames are added together. This approach has the advantage of preserving the highest spatial resolution. However, a substantial fraction of the pixels in each frame in this case are interpolated values, potentially sacrificing some photometric precision. As noted above, CRs contaminate approximately 5% of the pixels in a 20 minute exposure. After remapping, as many as 20% of the pixels must be assigned a reduced weight for the final averaging since a single pixel often gets remapped partially into several pixels in the new image.
Fortunately, both approaches to the data reduction were found to give very similar results. While the ”Frame by Frame” method produced a higher apparent S/N ratio than the ”Shift and Add” method, the former method is more prone to the introduction of systematic error. We adopted the more conservative approach of accepting a somewhat higher variance rather than risk the introduction of a bias. We therefore adopted the ”Shift and Add” spectra for the analysis of the LOSVDs in the present study. Figure 1 shows spectra at several positions along the aperture. As noted above, we are continuing in our efforts to refine the data reduction techniques still further.
## 3 Recovery of the Stellar Velocity Distribution
### 3.1 Method
An observed spectrum $`I(\lambda )`$ is the convolution of the line-of-sight velocity distribution $`N(V)`$ of the stars within the aperture with the spectrum of a single star $`T(\lambda )`$:
$$I(\mathrm{ln}\lambda )=N(V)T(\mathrm{ln}\lambda V/c)𝑑V.$$
(1)
The goal is to extract an estimate $`\widehat{N}(V)`$ of the true stellar broadening function $`N(V)`$ given $`I(\lambda )`$ and $`T(\lambda )`$, both observed with the same instrument. For $`T(\lambda )`$ we adopt the template spectrum of Figure The Nuclear Dynamics of M32. I..
Two independent deconvolution routines were used. The first algorithm, the “Fourier Correlation Quotient” (FCQ) method (Bender (1990); Bender, Saglia & Gerhard (1994)), constructs an estimate of the broadening function using Fourier techniques. The FCQ routine differs from earlier Fourier algorithms (e.g. Sargent et al. (1977)) in that the deconvolution is based on the template-galaxy correlation function rather than on the spectra themselves. This approach is less sensitive to template mismatch (Bender (1990)). The second algorithm, “Maximum Penalized Likelihood” (MPL), finds $`\widehat{N}(V)`$ as the solution to a penalized likelihood problem. The MPL estimate of $`N(V)`$ is computed on a grid in $`V`$ in such a way as to optimize the fit of the convolved template to the galaxy spectrum, subject to a “penalty” that measures the lack of smoothness of $`\widehat{N}(V)`$ (Merritt (1997)).
Both algorithms are nonparametric in the sense that no explicit constraints are placed on the functional form of $`\widehat{N}(V)`$. However they differ in two ways that are important for the current study. The FCQ algorithm requires that the absorption lines in the template spectrum be narrow compared to the broadened lines of the galaxy spectrum, i.e. that the galaxy velocity dispersion be large compared to the instrumental resolution. The MPL routine works well even when the galaxy velocity dispersion is small, as long as both template and galaxy spectra are observed at the same spectral resolution, at least in the case that the template star and galaxy have the same intrinsic absorption line properties (an assumption that will not be tested here).
The two algorithms differ also in the way they deal with the amplification of noise that accompanies the deconvolution. The FCQ routine uses a Wiener filter to suppress high-frequency components of the template-galaxy correlation function $`\stackrel{~}{K}_{T,G}`$. The degree of smoothing is determined by a factor, called here $`W`$, which fixes the width of the Gaussian function used to model the low-frequency, or signal, component of $`\stackrel{~}{K}_{T,G}`$. The choice $`W=1`$ corresponds to “optimal” filtering and larger values produce less smoothing; the FCQ algorithm adopts a default value of $`W=1`$ but automatically increases $`W`$ (to a maximum of 1.3) if the recovered LOSVD shows evidence of significantly non-Gaussian wings. In the MPL algorithm, the level of smoothing is determined by a factor $`\alpha `$ that multiplies the smoothness penalty function. This penalty function defines any $`N(V)`$ that is Gaussian as “smooth,” regardless of its mean or dispersion, via Silverman’s (1982) prescription; in the limit $`\alpha \mathrm{}`$, the MPL estimate of $`N(V)`$ is the Gaussian function which is most consistent, in a maximum-likelihood sense, with the galaxy spectrum. There is no a priori way of computing the optimum value of $`\alpha `$ in the MPL algorithm, a point that we return to below.
The different effects of smoothing on the form of $`\widehat{N}(V)`$ are illustrated in Figure The Nuclear Dynamics of M32. I., which shows estimates of the LOSVD in the central resolution element of M32 as computed by the two routines. Both algorithms produce rapidly-fluctuating solutions when undersmoothed, a consequence of the amplification of noise that always accompanies deconvolution. The only significant difference in this regime is the non-negativity of the MPL estimates, a consequence of the logarithmic form of the penalty function (Silverman (1982)). As the smoothing is increased, systematic differences begin to appear which are related to the different smoothing algorithms in the two codes. Solutions obtained via MPL tend to be more robust with respect to the degree of smoothing, producing in the limit of large $`\alpha `$ a Gaussian fit. However Figure The Nuclear Dynamics of M32. I. suggests that estimates of certain quantities, e.g. the wings of the LOSVD, might depend sensitively on the choice of smoothing level in either algorithm.
Once an estimate of $`N(V)`$ has been obtained, various quantites related to the line-of-sight velocity distribution can be derived. The simplest of these are the mean and rms velocities, which we denote by $`\overline{V}`$ and $`\sigma `$ respectively. As is well known, both quantities are difficult to estimate for numerically-recovered LOSVD’s since they are sensitively dependent on the form of $`\widehat{N}(V)`$ at large velocities where this function is most poorly determined. A standard alternative is to describe $`\widehat{N}(V)`$ by a Gram-Charlier or Gauss-Hermite (GH) series, the product of a normalizing Gaussian with a sum of Hermite polynomials $`H_i`$, both expressed in terms of $`(VV_0)/\sigma _0`$ (Thompson & Tapia (1990)). The parameters $`V_0`$ and $`\sigma _0`$ take the place of $`\overline{V}`$ and $`\sigma `$; while their definitions are to an extent arbitrary, these parameters are typically determined by requiring the coefficients of $`H_1`$ and $`H_2`$, called $`h_1`$ and $`h_2`$, to be zero (Gerhard (1993); van der Marel & Franx (1993)). Because $`V_0`$ and $`\sigma _0`$ describe the Gaussian core of the LOSVD, they are relatively insensitive to deviations of $`\widehat{N}(V)`$ from Gaussianity at high velocities. Information about these deviations is contained in the higher-order coefficients $`h_3`$, $`h_4`$ etc.; $`h_3`$ measures asymmetries in $`N(V)`$ and $`h_4`$ measures the strength of symmetric, non-Gaussian wings.
The FCQ and MPL algorithms derive the GH parameters from $`\widehat{N}(V)`$ in slightly different ways; details are given in Appendix A.
When applied to the STIS M32 spectra, the two algorithms were found to give consistent results for the lowest moments of $`N(V)`$, i.e. $`V_0`$, $`\sigma _0`$ and $`h_3`$. However the estimates of $`h_4`$ differed significantly at positions outside of the central $`0.1^{\prime \prime }`$. The FCQ algorithm gave $`0.15\begin{array}{c}<\hfill \\ \hfill \end{array}\widehat{h}_4\begin{array}{c}<\hfill \\ \hfill \end{array}0`$ at almost all positions; negative values of $`h_4`$ imply an $`N(V)`$ that falls off more sharply than a Gaussian at large velocities. The MPL algorithm gave values for $`h_4`$ in the range $`0\begin{array}{c}<\hfill \\ \hfill \end{array}h_4\begin{array}{c}<\hfill \\ \hfill \end{array}0.1`$, almost all positive, corresponding to LOSVD’s with super-Gaussian wings. Positive values of $`h_4`$ are expected near a black hole (Bahcall & Wolf (1976); van der Marel (1994)) and are also characteristic of models with radially-anisotropic velocity distributions.
We discuss the origin of this discrepancy in Appendix B. We believe that the primary reason for the systematic difference in $`\widehat{h}_4`$ values is the low velocity dispersion of M32. When a galaxy’s velocity dispersion is comparable to the dispersion of the template star spectrum ($`50`$ km s<sup>-1</sup> in the case of HR7615), the FCQ algorithm has difficulty recovering the true LOSVD (Figures B2,3; Bender, Paquet & Nieto (1991)). The $`\widehat{N}(V)`$’s recovered by FCQ in this regime are more sharply truncated than the true $`N(V)`$’s, leading to systematically low estimates of $`\widehat{h}_4`$. For values of $`\sigma _0`$ and S/N comparable to those of M32 at $`1^{\prime \prime }`$, Figure B2 shows that the estimates of $`h_4`$ generated by FCQ depend only weakly on the true $`h_4`$, with a bias that approaches $`0.1`$ for a true $`h_4`$ of 0.1. The MPL algorithm suffers much less from this bias (Figures B3-5).
Bias in nonparametric function estimates can always be reduced by reducing the degree of smoothing (e.g. Silverman (1986)), which in the case of the FCQ algorithm means increasing $`W`$. Figure B5c suggests that increasing $`W`$ from its default value of $`1`$ to values of $`2`$ can reduce the bias in FCQ estimates of $`h_4`$ by factors of $`2`$ or greater, even when $`\sigma _0`$ is as large as $`100`$ km s<sup>-1</sup>. We carried out this experiment with the STIS data; the results are shown in Figure The Nuclear Dynamics of M32. I.. The average $`\widehat{h}_4`$ values recovered by FCQ in M32 are indeed dependent on $`W`$; a change in $`W`$ from $`1`$ to $`1.5`$ has the effect of increasing $`\widehat{h}_4`$ from $`0.08`$ to $`+0.08`$. The latter value is essentially identical to the mean value of $`\widehat{h}_4`$ recovered via MPL.
We note that the dependence of $`\widehat{h}_4`$ on $`W`$ could be due either to the suggested explanation, i.e. the need to include more frequency channels for dispersions close to the instrumental resolution, or alternatively to a mismatch between galaxy and template spectral properties in the wings of the lines which results in incorrect continuum subtraction. We will not explore the second possibility here but note again that the Monte Carlo experiments in Appendix B suggest that values of $`W`$ of order 2 are appropriate even when there is no template mismatch.
The values of $`\widehat{h}_4`$ recovered by MPL are also dependent on the value of the smoothing parameter $`\alpha `$ but much less so (cf. Figure The Nuclear Dynamics of M32. I.), until $`\alpha `$ is made so large that the LOSVD is forced into a Gaussian shape. The Monte Carlo tests summarized in Figure B5b suggest that the bias in $`\widehat{h}_4`$ as recovered by MPL is likely to be only of order $`0.02`$, several times smaller than with FCQ.
We conclude that the values of $`\widehat{h}_4`$ recovered by the two algorithms are consistent once their different biases are taken into account and that the values returned by MPL are likely to be more accurate. Henceforth we adopt the MPL estimates.
The full set of GH parameters derived from the STIS spectra and their $`1\sigma `$ confidence intervals are given in Table 2 (FCQ) and Table 3 (MPL). For radii $`\begin{array}{c}<\hfill \\ \hfill \end{array}0.7^{\prime \prime }`$ from the center the data were sampled at full resolution($`0.05^{\prime \prime }`$) while at larger radii they were binned spatially. The sampling at small radii is fine enough that the data points are somewhat correlated; this was done to ensure that no information concerning the steep radial gradients of the profiles was lost.
### 3.2 Results for M32
Figure The Nuclear Dynamics of M32. I. presents LOSVDs computed via the MPL algorithm at positions separated by about $`0.1^{\prime \prime }`$ along the slit. One expects these broadening functions to obey $`N(V;R)=N(V;R)`$, since for a point-symmetric galaxy, the velocity distributions should reverse after passing from one side of the galaxy to the other. The LOSVDs of Fig. The Nuclear Dynamics of M32. I. show approximately the expected symmetry. The right-hand column of Figure The Nuclear Dynamics of M32. I. plots mean broadening functions averaged over the two sides of the galaxy, $`\overline{N}(V)=\frac{1}{2}[N(V,R)+N(V,R)]`$; the central LOSVD has been symmetrized about $`V=0`$.
These broadening functions show clear and consistent deviations from Gaussian form, in two respects. First, the central LOSVD exhibits strong super-Gaussian “wings” at high velocities. These wings are possibly present also in some of the off-center LOSVDs although with lower amplitude. Second, the off-center LOSVDs are asymmetric, with tails extending toward velocities opposite in sign to the mean velocity at each radius. These asymmetric tails are similar to those exhibited by a rotating system superimposed on a slowly-rotating bulge.
Systematic problems in the spectral deconvolution, e.g. template mismatch or incorrect continuum subtraction, can easily produce features like the wings and tails seen in the broadening functions of Figure The Nuclear Dynamics of M32. I.. Such errors in most cases would be expected to produce features located at the same velocity on both sides of the galaxy (e.g. Bender, Saglia & Gerhard (1994)) and are therefore an unlikely explanation for the asymmetric tails seen in the off-center LOSVDs. The strong wings seen in the central LOSVD might more plausibly be attributed to systematic errors. However we found that the wings in the central LOSVD were robust; they appeared in both the MPL and FCQ estimates of $`N(V)`$ (though less clearly in the latter – see Figure The Nuclear Dynamics of M32. I.) and were relatively unaffected by changes in the assumed continuum level or slope. We carried out MPL deconvolutions where the fit to the galaxy spectrum was restricted to the region around only one, or two, of the three calcium-triplet lines; these LOSVDs also exhibited strong wings. We also tried using one of the other available STIS stellar templates; again the wings were only slightly affected. (The adopted template, Figure 1, produced the best overall fit to the galaxy spectrum.) Finally, we show in §4.2 that the wings are consistent with those predicted by stars in the gravitational field of a supermassive black hole.
Figure The Nuclear Dynamics of M32. I. shows the Gauss-Hermite parameter $`V_0`$, a measure of the stellar rotation, in the inner arc second of M32. Also plotted are $`h_3`$, the lowest, odd GH moment of the LOSVD, and the “corrected” rotation velocity, $`V_{0,c}=V_0+\sqrt{3}\sigma _0h_3`$. $`V_{0,c}`$ is a closer approximation than $`V_0`$ to the true mean line-of-sight velocity $`\overline{V}`$ (van der Marel & Franx (1993)). The corrected rotation velocity is lower in absolute magnitude than $`|V_0|`$ due to the asymmetric wings of the LOSVD noted above.
The STIS rotation curve is consistent with earlier ground-based measurements (Figure The Nuclear Dynamics of M32. I.) at radii $`\begin{array}{c}>\hfill \\ \hfill \end{array}1^{\prime \prime }`$ but with a larger peak value, $`60`$ km s<sup>-1</sup>. Furthermore the rotation curve remains flat or slightly rising into smaller radii than seen heretofore, before falling at $`R\begin{array}{c}<\hfill \\ \hfill \end{array}0.25^{\prime \prime }`$ due to the blending of light from the two sides of the galaxy. There is a suggestion of an east-west asymmetry in the rotation curve though the effect is probably not significant.
The $`h_3`$ profile is approximately antisymmetric about the center of M32, as expected in a relaxed galaxy. $`|h_3|`$ reaches a maximum value of $`0.05`$ at $`|R|0.3^{\prime \prime }`$ and appears to gradually decline at larger radii. This behavior is similar to that predicted in axisymmetric models (e.g. Figure 8 of Dehnen (1995)) where $`h_3`$ remains essentially constant at radii outside the seeing disk.
The Gauss-Hermite parameter $`\sigma _0`$ is shown in Figure The Nuclear Dynamics of M32. I.. Also plotted is $`h_4`$, the lowest, even moment of the LOSVD, and the “corrected” velocity dispersion, $`\sigma _{0,c}=\sigma _0(1+\sqrt{6}h_4)`$; $`\sigma _{0,c}`$ is a closer approximation than $`\sigma _0`$ to the true rms velocity $`\sigma `$. The velocity dispersion rises suddenly inside of $`0.3^{\prime \prime }`$, approximately the same radius at which the rotation curve begins to fall. This coincidence suggests that at least part of the rise in $`\sigma _0`$ is due to averaging of the rotation velocity over the two sides of the galaxy near the center, which has the effect of converting a rotation into an apparent dispersion (Tonry (1987)).
The corrected velocity dispersion $`\sigma _{0,c}`$ rises well above $`\sigma _0`$ near the center due to the strong non-Gaussian wings of the LOSVD. The central value of $`\sigma _{0,c}`$ is $`175`$ km s<sup>-1</sup>; this should probably be interpreted as a lower limit since $`h_4`$ is only sensitive to the inner parts of the wings. (We argue below, based on model fits, that the rms velocity in the central resolution element may be as high as $`200`$ km s<sup>-1</sup>.) The ground-based data (Figure The Nuclear Dynamics of M32. I.) are consistent with the STIS dispersions at radii $`\begin{array}{c}>\hfill \\ \hfill \end{array}1^{\prime \prime }`$ but fail to resolve the continued rise in $`\sigma _0`$ inside of $`0.5^{\prime \prime }`$.
Dynamical models (e.g. Dehnen (1995); Qian et al. (1995)) predict $`h_4(R)`$ profiles similar to that in Figure The Nuclear Dynamics of M32. I. when observed with HST resolution: a central maximum; a rapid drop, to small or negative values, at $`R0.1^{\prime \prime }`$; and a nearly constant value at larger radii. The predicted drop at $`0.1^{\prime \prime }`$ is due to blending of the light from the two sides of the galaxy, which broadens the low-velocity part of the LOSVD and lowers the observed $`h_4`$. The predicted central value of $`h_4`$ depends strongly on the black hole mass and on the PSF; our value, $`\widehat{h}_40.14`$, is larger than in the two studies just cited, but these studies were based on rather low assumed black hole masses, $`M_h=12\times 10^6_{}`$. The true black hole mass is probably greater (van der Marel et al. (1998)).
The behavior of $`\widehat{h}_4`$ at larger radii is surprising. Previous observational studies (e.g. van der Marel et al. 1994a ; Bender, Kormendy & Dehnen (1996)) have returned smaller estimates for $`h_4`$ in M32. However we believe that these earlier results are not inconsistent with ours giving the difficulties involved with estimating this parameter. The van der Marel (1994a) study was based on WHT observations with a much lower spatial resolution than the STIS data. At radii $`\begin{array}{c}<\hfill \\ \hfill \end{array}1^{\prime \prime }`$, the value of $`\widehat{h}_4`$ recovered by those authors was strongly affected by the PSF blending of the rotation curve discussed above, yielding negative values in the central aperture. Outside of $`2^{\prime \prime }`$, van der Marel et al. found $`\widehat{h}_4`$ to increase sharply to $`0.03`$ on both sides of the galaxy (their Figure 12). Van der Marel (private communication) notes that the values of $`\widehat{h}_4`$ derived from the WHT data depended sensitively on the choice of template spectrum and on the algorithm for continuum subtraction. Using a single, best-fit template, $`\widehat{h}_4`$ was found to lie between $`0.3`$ and $`0.5`$ throughout the inner $`2^{\prime \prime }`$; the lower values of $`\widehat{h}_4`$ in the published paper were derived using a spectral deconvolution routine that constructs an “optimal” template by linear superposition of a set of stellar spectra. In another ground-based study, Bender, Kormendy & Dehnen (1996) applied the FCQ algorithm to CFHT data of higher spatial resolution and found $`\widehat{h}_40.05`$ inside of $`0.2^{\prime \prime }`$, gradually falling to $`0`$ at $`1.0^{\prime \prime }`$. However the spectral resolution in this study was only $`80`$ km s<sup>-1</sup>and the derivation of $`h_3`$ and $`h_4`$ correspondingly difficult; as noted above, we also found $`\widehat{h}_40`$ from the STIS data using the FCQ algorithm and argued that these values were significantly negatively biased.
Although we believe that all of these studies are consistent with our conclusion that $`h_4`$ is significantly positive throughout the nucleus of M32, we are less willing to strongly endorse the precise $`\widehat{h}_4`$ values shown in Figure The Nuclear Dynamics of M32. I., due to the sensitive dependence of this parameter on the details of the spectral deconvolution algorithm, continuum subtraction, smoothing level, etc. In Paper II we will construct dynamical models based on a range of assumed $`h_4`$ profiles in order to test the dependence of the inferred black hole mass on this parameter.
We may also compare our results to the van der Marel et al. (1997, 1998) HST/FOS measurements of $`V_0`$ and $`\sigma _0`$ (Figure The Nuclear Dynamics of M32. I.). The FOS measurements were taken through square apertures as small as $`0.1^{\prime \prime }`$ on a side, hence their spatial resolution is comparable to that of the STIS data. However the FOS is a low spectral resolution instrument and not well suited to objects like M32 with a relatively low velocity dispersion; furthermore there are difficulties in positioning the FOS and these were probably the cause of the large point-to-point variations seen by van der Marel et al. One advantage of STIS over FOS is the continuous spatial sampling which avoids potential errors in aperture placement. We find a hint in the STIS data of the asymmetry seen in the FOS $`\sigma _0(R)`$ profile (a more rapid falloff on the west side). The central FOS value of $`\sigma _0`$ seems significantly bigger than found here, and the FOS rotation velocities are systematically lower.
## 4 Analysis
The STIS data show what appear to be clear signatures of the gravitational influence of a massive compact object on the stellar velocity distribution within the central parsec of M32. Here we use simple dynamical models to address the question of whether these observed features are consistent with the presence of a supermassive black hole. Our aim is not to derive the best possible estimate of the black hole mass or the stellar velocity distribution – those are the goals of Paper II – but rather to address two, more basic issues concerning the interpretation of the data.
1. The velocity dispersion profile (Figs. The Nuclear Dynamics of M32. I., The Nuclear Dynamics of M32. I.) exhibits a sudden upturn at a distance of $`0.3^{\prime \prime }`$ from the center, presumably due in part to the gravitational force from a massive compact object. At roughly the same radius, the rotation curve falls (Figs. The Nuclear Dynamics of M32. I., The Nuclear Dynamics of M32. I.), presumably due to blending of light from opposite sides of the galaxy which are rotating in opposite directions. The blending would also be expected to contribute to the rise in the observed velocity dispersion (Tonry (1987)), consistent with the fact that the upturn in $`\sigma `$ and the drop in $`\overline{V}`$ occur at roughly the same radius. We would like to estimate the degree to which the velocity dispersion spike is a product of this blending, and the degree to which it is due to a real upturn in the stellar random velocities. In other words: Do the STIS data resolve the black hole’s sphere of influence?
2. The central LOSVD in M32 exhibits strong, super-Gaussian wings (Fig. The Nuclear Dynamics of M32. I., The Nuclear Dynamics of M32. I.). Such wings are a generic prediction of the black hole model (Bahcall & Wolf (1976); van der Marel (1994); Dehnen (1995)); they result from stars on high-velocity orbits within the black hole’s sphere of influence. However systematic errors in the spectral deconvolution can also produce spurious features in the LOSVD’s, particularly at large velocities where the form of the broadening function is only weakly constrained by the spectra. We would like to verify that the inferred, non-Gaussian shape of the central LOSVD is consistent with that expected from the black hole model.
We note that M32 is the only galaxy so far to exhibit either a resolved central spike in the stellar velocity dispersions, or strong non-Gaussian wings in the nuclear LOSVD. Either feature, if observed with sufficient spatial resolution and S/N, could independently place strong constraints on the mass of a central black hole. We will in fact generate estimates of $`M_h`$ from our analyses of both features but we stress that the best estimates of $`M_h`$ can only come from more complete modelling based on the entire kinematical data set.
### 4.1 The Velocity Dispersion Spike
One expects to see a rise in the stellar velocities in a hot stellar system at a distance $`r_h`$ from the black hole, where
$$r_h=\frac{GM_h}{\sigma _{}^2},$$
(2)
the black hole’s “radius of influence” (Peebles (1972)). Here $`M_h`$ is the black hole mass and $`\sigma _{}^2`$ is the stellar velocity dispersion at $`r>r_h`$. Setting $`M_h=10^6_{}`$ and $`\sigma _{}=60`$ km s<sup>-1</sup> (Figures The Nuclear Dynamics of M32. I., The Nuclear Dynamics of M32. I.) gives $`r_h1.2\mathrm{pc}0.3^{\prime \prime }`$. This value seems comfortably larger than the HST/STIS resolution but it is based on an assumed value of $`M_h`$ and furthermore it refers to the true radius whereas we observe the galaxy in projection, which tends to hide otherwise sharp features. Modelling the spike therefore requires us to predict the true, projected velocity field of the galaxy in two dimensions on the plane of the sky, including both random and rotational velocities, and then to convolve it with the STIS PSF.
We begin by constructing solutions to the stellar hydrodynamical equations. We assume the simplest possible axisymmetric model for the stars, in which the velocity dispersions are isotropic in the meridional plane $`(\varpi ,z)`$, i.e. $`\sigma _\varpi (\varpi ,z)=\sigma _z(\varpi ,z)\sigma (\varpi ,z)`$; the model is flattened by an inequality between $`\sigma ^2`$ and the mean square azimuthal velocity $`\overline{v_\varphi ^2}`$. The second moments of the internal velocity distribution are given by
$$\nu \sigma ^2=_z^{\mathrm{}}\nu \frac{\mathrm{\Phi }}{z}𝑑z,\nu \overline{v_\varphi ^2}=\nu \sigma ^2+\nu \varpi \frac{\mathrm{\Phi }}{\varpi }+\varpi \frac{(\nu \sigma ^2)}{\varpi }$$
(3)
where $`\nu `$ is the stellar number density, $`v_\varphi `$ is the azimuthal velocity and $`\mathrm{\Phi }`$ is the combined gravitational potential from the stars and the central black hole, $`\mathrm{\Phi }(\varpi ,z)=\mathrm{\Phi }_{}(\varpi ,z)GM_h/r`$ (Merritt (1999)).
We evaluated these expressions assuming a stellar density
$$\nu (\varpi ,z)=\nu _0(m/b)^\alpha \left[1+(m/b)^2\right]^\beta \left[1+(m/c)^2\right]^\gamma ,m^2=\varpi ^2+(z/q)^2,$$
(4)
a parametrized form proposed by van der Marel et al. (1998); those authors found a good match between their model (observed edge-on) and M32 with the parameters $`\alpha =1.435,\beta =0.423,\gamma =1.298,b=0.55^{\prime \prime },c=102.0^{\prime \prime },q=0.73`$ and $`\nu _0=0.463\times 10^5L_{,V}`$pc<sup>-3</sup>. The luminosity density $`\nu _0`$ is converted into a mass density $`\rho _0`$ by the factor $`(M/L)(_{}/L_{})`$, with $`M/L`$ the mass-to-light ratio of the stars in solar units.
Given $`\sigma (\varpi ,z)`$ and $`\overline{v_\varphi ^2}(\varpi ,z)`$ as obtained from equations (3) and (4), the projected, line-of-sight, mean square velocity $`\overline{V^2}`$ is obtained by a density-weighted integration through the galaxy, assumed here to be edge-on:
$$\mathrm{\Sigma }(R,Z)\overline{V^2}(R,Z)=2_R^{\mathrm{}}\nu (\varpi ,z)\left[\left(1\frac{R^2}{\varpi ^2}\right)\sigma ^2(\varpi ,z)+\frac{R^2}{\varpi ^2}\overline{v_\varphi ^2}(\varpi ,z)\right]\frac{\varpi d\varpi }{\sqrt{\varpi ^2R^2}}$$
(5)
(Fillmore (1986)), where $`(R,Z)`$ are coordinates on the plane of the sky and $`R`$ is measured parallel to the long axis of the galaxy’s figure; $`\mathrm{\Sigma }(R,Z)`$ is the stellar surface density.
Values of $`\overline{V^2}`$ were computed on a rectangular grid of $`180\times 25`$ locations with separations of $`0.015^{\prime \prime }`$ in $`R`$ and $`0.02^{\prime \prime }`$ in $`z`$. These values were then convolved with the STIS PSF and averaged over the pixel area and the aperture after weighting by the model surface brightness. The STIS PSF at $`8500`$Å has a FWHM of $`0.115^{\prime \prime }`$. The PSF is also slightly asymmetric (Bower et al. (2000)). We ignored this slight asymmetry when carrying out the convolutions with our models.
The second velocity moments of models constructed in this way are uniquely determined by the two parameters $`(M_h,M/L)`$ that specify the potential. Figure The Nuclear Dynamics of M32. I. shows the goodness of fit of the models to the observed, mean square velocities; only data points within the inner $`1.0^{\prime \prime }`$ were used in evaluating $`\chi ^2`$. As estimates of $`\overline{V^2}`$, we took $`V_{0,c}^2+\sigma _{0,c}^2`$, where $`V_{0,c}`$ and $`\sigma _{0,c}`$ are the Gauss-Hermite parameters corrected by $`h_3`$ and $`h_4`$ respectively (Figures The Nuclear Dynamics of M32. I., The Nuclear Dynamics of M32. I.). The best-fit model has $`M_h3.2\times 10^6_{}`$ and $`M/L3.3`$ with $`\stackrel{~}{\chi ^2}=0.64`$; a $`\stackrel{~}{\chi ^2}`$ of unity includes models with $`M_h`$ as small as $`2.2\times 10^6_{}`$ and as large as $`4.3\times 10^6_{}`$.
The degree of net rotation in our models may be adjusted by partitioning the azimuthal motions between streaming, $`\overline{v_\varphi }`$, and dispersion, $`\sigma _\varphi `$. We followed the standard practice (Satoh (1980)) of making $`\sigma _\varphi ^2`$ a weighted average of $`\sigma ^2`$ and $`\overline{v_\varphi ^2}`$, i.e.
$$\sigma _\varphi ^2=k^2\sigma ^2+(1k^2)\overline{v_\varphi ^2}.$$
(6)
The parameter $`k`$ (assumed independent of position) may be varied between zero (corresponding to no streaming motions) and a maximum value, of order unity, at which $`\sigma _\varphi ^2`$ is forced below zero at some point in the meridional plane; $`k=1`$ yields an “isotropic oblate rotator.”
When we add $`k`$ as a free parameter and require the models to fit the observed rotation and velocity dispersion profiles separately, the best-fit values of $`(M_h,M/L)`$ were nearly unchanged but $`\stackrel{~}{\chi ^2}`$ increased to $`3.7`$ – since the model is now being asked to fit twice as many data points with only one extra parameter. The best-fit value of $`k`$ was found always to be close to $`1.2`$, implying slightly smaller $`\sigma _\varphi `$ (i.e. greater rotation) than in an isotropic oblate rotator. Figure The Nuclear Dynamics of M32. I. compares the data with the predicted profiles for $`M_h=3.0\times 10^6_{}`$, close to the best-fit value, and for $`M_h=2.0`$ and $`4.0\times 10^6_{}`$.
We draw the following conclusions from these comparisons.
1. The lowest-order moments of the line-of-sight velocity distribution in M32 are reasonably well fit near the center by our simple axisymmetric model, with a black hole mass $`M_h3\times 10^6_{}`$. The rotation curve is best fit by a smaller mass ($`12\times 10^6_{}`$) and the velocity dispersions by a larger mass ($`34\times 10^6_{}`$); if we require the models only to fit the mean square velocities, the fit is essentially perfect within the inner arc second.
2. The STIS observations probably come close to resolving the central upturn in the stellar velocity dispersions, which is predicted to occur at a projected radius of $`0.1^{\prime \prime }`$ for $`M_h2\times 10^6_{}`$ and $`0.2^{\prime \prime }`$ for $`M_h4\times 10^6_{}`$.
3. Smearing of the stellar rotation field probably accounts for only a small part of the observed upturn in the dispersions.
We note that Dehnen (1995) was able to improve the fit of his axisymmetric models to the M32 data then available by varying the ratio of rotational to non-ordered azimuthal motions, which in our models would correspond to varying $`k`$ with position.
If our simple model for the internal dynamics of M32 is approximately correct, Figures The Nuclear Dynamics of M32. I. and The Nuclear Dynamics of M32. I. imply that the STIS data can place strong constraints on the mass of the central black hole. However such a conclusion must await the results of the more complete modelling of Paper II. That study may yield tighter constraints on $`M_h`$, due to the use of the full kinematical data set; or weaker constraints, due to the increased flexibility of general, anisotropic models.
### 4.2 The Central Broadening Function
The LOSVD near the projected center of a galaxy containing a black hole is expected to be very non-Gaussian due to high velocity stars orbiting near the central mass (Bahcall & Wolf (1976)). If the stellar density follows a power law into the center, $`\rho r^\gamma `$, the LOSVD in an aperture containing the black hole has power-law wings, $`N(V)V^{2\gamma 7}`$, $`V\mathrm{}`$ (van der Marel (1994)). The amplitude of these wings depends on the ratio of black hole mass to stellar mass within the aperture, and on the slope $`\gamma `$ of the stellar density profile (Dehnen (1995)), among other factors. The wings are most prominent in a galaxy, like M32, for which the stellar cusp is steep, $`\gamma 2`$, since a large fraction of the light near the projected center comes from the region near the black hole.
Here we ask whether the strong wings seen in the central M32 broadening function (Figures The Nuclear Dynamics of M32. I., The Nuclear Dynamics of M32. I.) are consistent with the black hole model. To answer this question we must compute a stellar distribution function $`f`$ and integrate it over the two velocity components in the plane of the sky (Merritt (1987)). The result must then be smeared by the instrumental resolution. We once again restrict ourselves to the simplest model which permits a meaningful comparison with the data; in this case, a spherical galaxy with an isotropic distribution function, $`f=f(E)`$. Our model will ignore rotation, which acts to broaden the central LOSVD but leaves it symmetric. In any case we do not know what the contribution of rotation to the stellar velocity distribution is very near to the black hole. For the stellar density profile we assume the spherical version of equation (4); $`f(E)`$ then follows from the standard formula (Eddington (1916)) and the projected velocity distribution is also straightforward to compute (Merritt (1993)). We chose to fix $`M/L`$ for the stars at $`3.5(_{}/L_{})`$ as the black hole mass was varied; this $`M/L`$ reproduces approximately the observed mean square velocities for $`rr_h`$.
Figure The Nuclear Dynamics of M32. I. compares the central LOSVD in M32 to the broadening functions predicted by the spherical model, for black hole masses $`M_h=(2.5,5.0,10.)\times 10^6_{}`$. The M32 LOSVD has been symmetrized about $`V=0`$. The fit is quite reasonable for $`M_h=5.0\times 10^6_{}`$, although the high-velocity wings are better fit by still larger masses. The model LOSVD for $`M_h=5.0\times 10^6_{}`$ has the parameters:
$$\sigma =178\mathrm{km}/\mathrm{s}\sigma _0=133\mathrm{km}/\mathrm{s}h_4=0.085.$$
(7)
The first two numbers are essentially identical to the values inferred from the M32 data (Figure The Nuclear Dynamics of M32. I.); $`h_4`$ is lower than, but consistent with, the M32 estimate ($`0.14\pm 0.03`$).
We conclude that the central LOSVD in M32 is consistent with that expected for a stellar nucleus containing a massive compact object, with a mass comparable to that found in the fit to the axisymmetric models.
## 5 Summary
We used HST and STIS to obtain stellar absorption line spectra near the center of M32 in a wavelength region centered on the Calcium triplet. The spectra were analyzed using two independent spectral deconvolution routines; these gave fully consistent results except in the case of the Gauss-Hermite $`h_4`$ parameter, but we argued that the differences could be reconciled after taking into account the different biases of the two algorithms. The stellar rotation velocities in M32 are slightly higher than observed from the ground and remain constant into $`0.25^{\prime \prime }`$ from the center. The velocity dispersions exhibit a clear spike beginning at approximately the same radius. These two kinematical profiles are consistent with those predicted by simple axisymmetric models containing central black holes with masses in the range $`25\times 10^6_{}`$. The stellar LOSVDs show significant deviations from Gaussian form as measured by the Gauss-Hermite parameters $`h_3`$ and $`h_4`$. The central LOSVD is particularly non-Gaussian, exhibiting strong, high-velocity wings. We showed that the amplitude of these wings is consistent with that predicted by simple models containing black holes with masses of order $`3\times 10^6_{}`$.
Detailed dynamical modelling of M32 based on these data and estimates of the black hole mass will be presented in Paper II.
We thank W. Dehnen and R. van der Marel for helpful discussions. This work was supported by NASA grants NAG 5-3158 and NAG 5-6037, by NSF grant AST 96-17088, and by STIS GTO funding. Data presented here were based on observations with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc. (AURA), under NASA contract NAS5-26555.
Appendix A
Gauss-Hermite Moments
The two spectral deconvolution algorithms described above yield nonparametric estimates $`\widehat{N}(V)`$ of the stellar LOSVD. Here we describe the methods used by the two algorithms to derive the GH moments from $`\widehat{N}(V)`$.
Let $`N(X,Y;V)`$ be the distribution of line-of-sight stellar velocities in the aperture centered at $`(X,Y)`$. Define the GH moments of $`N`$ as
$$h_i(X,Y)=2\sqrt{\pi }_{\mathrm{}}^{\mathrm{}}N(X,Y;V)g(w)H_i(w)𝑑V,$$
()
where $`H_i`$ are the Hermite polynomials (as defined by Gerhard (1993)) and the weight function
$$g(w)=\frac{1}{\sqrt{2\pi }\gamma _0}e^{w^2/2},w=(VV_0)/\sigma _0$$
()
has three free parameters $`(\gamma _0,V_0,\sigma _0)`$. Following van der Marel & Franx (1993), we choose these parameters at every point $`(X,Y)`$ such that
$$h_0(X,Y)=1,h_1(X,Y)=h_2(X,Y)=0.$$
()
These definitions impose the following implicit conditions on $`(\gamma _0,V_0,\sigma _0)`$:
$$\gamma _0=\sqrt{2}\sigma _0_{\mathrm{}}^{\mathrm{}}N(V)e^{w^2/2}𝑑w,$$
()
$$0=_{\mathrm{}}^{\mathrm{}}N(V)e^{w^2/2}w𝑑w,$$
()
$$0=_{\mathrm{}}^{\mathrm{}}N(V)e^{w^2/2}(2w^21)𝑑w.$$
()
The relations (A4) define a nonlinear minimization problem with solutions $`(\gamma _0,V_0,\sigma _0)`$ given $`N(V)`$.
The MPL algorithm (Merritt (1997)) derives the three parameters in just this way, using the NAG routine E04FDF to minimize the sum $`(h_01)^2+h_1^2+h_2^2`$ as a function of $`(\gamma _0,V_0,\sigma _0)`$. The higher-order GH moments are then derived using equation (A1), by numerical integration over $`\widehat{N}(V)`$.
Most spectral deconvolution algorithms of which we are aware derive the parameters $`(\gamma _0,V_0,\sigma _0)`$ in a different way. The LOSVD is compared to the trial function
$$𝒩(V)=\frac{\gamma _0}{\sqrt{2\pi }\sigma _0}e^{w^2/2}\left[1+\underset{j=3}{\overset{j_{max}}{}}h_jH_j(w)\right]$$
()
where $`j_{max}`$ is the index of the highest GH moment fitted to $`\widehat{N}(V)`$; typically $`j_{max}=4`$. The integrated square deviation between $`\widehat{N}(V)`$ and $`𝒩(𝒱)`$ is then minimized by varying the $`(j_{max}+1)`$ free parameters $`(\gamma _0,V_0,\sigma _0,h_3,h_4,\mathrm{},h_{j_{max}})`$. This is the technique used by the FCQ algorithm.
A theorem (Myller-Lebedeff (1908)) guarantees the equivalence of the two approaches if $`j_{max}=\mathrm{}`$ in equation (A5) (van der Marel & Franx (1993)). However if $`j_{max}\mathrm{}`$, and if the input $`N(V)`$ can not be precisely represented by a finite GH series with $`jj_{max}`$, the results given by the two algorithms will differ. For instance, in attempting to represent an $`N(V)`$ having $`h_60`$ using $`j_{max}=4`$, the FCQ algorithm will adjust $`\sigma _0`$ and $`h_4`$ to incorrect values in order to better fit the high-velocity wings of the profile with the limited number of terms allowed to it. This is illustrated in Figure A1, which shows the values of $`\sigma _0`$ and $`h_4`$ generated by the second algorithm, $`\widehat{\sigma }_0`$ and $`\widehat{h}_4`$, compared to the true values for an input $`N(V)`$ with $`h_4=0.15`$ and nonzero $`h_6`$:
$$N(V)=\frac{1}{\sqrt{2\pi }}e^{V^2/2}\left\{1+0.15H_4(V)+h_6H_6(V)\right\}.$$
()
For $`|h_6|\begin{array}{c}>\hfill \\ \hfill \end{array}0.1`$, the errors in $`\sigma _0`$ and $`h_4`$ as derived from the second algorithm are $`\begin{array}{c}>\hfill \\ \hfill \end{array}15\%`$ and $`\begin{array}{c}>\hfill \\ \hfill \end{array}20\%`$ respectively.
Appendix B
Performance Evaluation of the FCQ and MPL Algorithms
Here we compare the performance of the FCQ and MPL algorithms given simulated data. Our primary goal is to understand the source of the systematic offset of $`h_4`$ values as derived from the M32 spectra by the two algorithms (§3.1). Two independent sets of tests were carried out, the first by R. Bender and the second by D. Merritt. All tests were based on synthesized galaxy spectra generated from the STIS template spectrum (Figure The Nuclear Dynamics of M32. I.) by convolving it with an assumed $`N(V)`$ and adding noise.
The first set of tests addressed the accuracy of FCQ estimates when the galaxy velocity dispersion is low. It is well known that the accuracy of FCQ begins to fall off when the galaxy velocity dispersion becomes comparable to the dispersion of the template spectrum (e.g. Bender, Paquet & Nieto (1991)). Figure B1 shows values of $`\widehat{\sigma }_0`$ recovered by FCQ given a Gaussian-broadened template spectrum and thirty random noise realizations. The default value ($`W1`$) of the smoothing parameter was used. There is a positive bias in the estimated values beginning at $`\sigma _0100`$ km s<sup>-1</sup>; the bias increases with decreasing $`\sigma _0`$ becoming significant for $`\sigma _050`$ km s<sup>-1</sup>. The bias is only weakly dependent on S/N. This bias in the estimation of $`\sigma _0`$ is unlikely to be important for the nucleus of M32 where $`\sigma _0\begin{array}{c}>\hfill \\ \hfill \end{array}100`$km s<sup>-1</sup>.
Figure B2 shows the performance of FCQ at recovering $`h_4`$. The template spectrum was broadened using an $`N(V)`$ of the form
$$N_1(V)=\frac{1}{\sqrt{2\pi }}e^{V^2/2\sigma _0^2}\left\{1+h_4H_4(V/\sigma _0)\right\}$$
()
with various values of $`\sigma _0`$ and $`h_4`$. Figure B2 reveals significant biases in $`\widehat{h}_4`$ for $`\sigma _0\begin{array}{c}<\hfill \\ \hfill \end{array}100`$ km s<sup>-1</sup>, even when S/N is as great as 100. When $`\sigma _050`$ km s<sup>-1</sup> and S/N $`30`$, characteristic of M32 at $`1^{\prime \prime }`$, the bias in $`h_4`$ is $`0.1`$ for an input $`h_4`$ of $`0.1`$.
The second set of tests compared the performance of the FCQ and MPL algorithms on galaxy spectra generated from the broadening function
$$N_2(V)=\frac{1}{\pi \sigma }\frac{1}{1+(V/\sigma )^2},$$
()
a Lorentzian function, with $`\sigma =100`$ km s<sup>-1</sup>. This LOSVD is qualitatively similar to what is expected in a black-hole cusp, with $`NV^2`$ high-velocity wings. The non-trivial GH parameters are
$$\gamma _0=0.76986\sigma _0=108.07\mathrm{km}\mathrm{s}^1h_4=0.14546h_6=0.01850.$$
()
Figure B3 shows mean estimates of $`N(V)`$ obtained using the two algorithms for 100 random realizations of the noise. The smoothing parameter in both algorithms was adjusted to minimize the mean square error of $`\widehat{N}(V)`$ (as defined below) for each value of S/N. There is a greater bias in the FCQ estimates, as well as a persistent “ringing” at high velocities.
Figure B4 plots the mean integrated square error (MISE) and integrated square bias (ISB) of the recovered broadening functions as functions of S/N; in the case of the FCQ algorithm, the integrated errors are shown both for the optimal choice of smoothing parameter $`W_{opt}`$ that minimizes the MISE, as well as for the value chosen by the algorithm ($`1.3`$). The MISE of an estimate $`\widehat{f}(x)`$ is defined as
$$\mathrm{MISE}\left[\widehat{f}(x)\right]=E\left\{\widehat{f}(x)f(x)\right\}^2𝑑x$$
()
$$=\left\{E\widehat{f}(x)f(x)\right\}^2𝑑x+\left(E\left\{\widehat{f}^2(x)\right\}E\left\{\widehat{f}(x)\right\}^2\right)𝑑x$$
()
$$=\mathrm{ISB}\left[\widehat{f(x)}\right]+\mathrm{IV}\left[\widehat{f(x)}\right],$$
()
the sum of the integrated square bias ISB and the integrated variance IV (Silverman (1986)); here $`E`$ denotes the expectation value, i.e. the average over many random realizations of the noise. The MISE and ISB displayed in Fig. B4 were divided by the normalizing factor $`\left[N_2(V)\right]^2𝑑V`$.
The MISE of the MPL estimates falls roughly as a power law, $`\mathrm{MISE}[\widehat{N}(V)](S/N)^{1.3}`$, close to the asymptotic $`(S/N)^1`$ of parametric estimators. Approximately $`1/2`$ of the total square error comes from the bias and $`1/2`$ from the variance. In the case of the FCQ algorithm, the MISE behaves in a more complicated way with S/N, at first falling with S/N then appearing to level off for S/N $`\begin{array}{c}>\hfill \\ \hfill \end{array}50`$. This levelling off is a consequence of the low-velocity-dispersion bias of FCQ discussed above. For the FCQ estimates, the bulk of the MISE is due to the variance; adjusting the smoothing parameter primarily affects the bias and has little effect on the MISE. For S/N $`20`$, the mean square error of the optimal FCQ estimate is a factor $`3`$ greater than that of the MPL estimates.
The bias in $`\widehat{N}(V)`$ is in the direction of wider and more steeply truncated functions, particularly in the case of the FCQ estimates (Fig. B3). This bias in $`\widehat{N}(V)`$ is consistent with the negative bias found above in estimates of $`h_4`$. Figure B5 compares the ability of the two algorithms to recover $`h_4`$ from the Lorentzian $`N_2(V)`$. Plotted there are the mean square error (MSE) and bias in estimates of $`h_4`$ from 100 random noise realizations; the MSE is defined, for any estimated parameter $`\widehat{P}`$, as
$$\mathrm{MSE}(\widehat{P})=E\left\{\widehat{P}P\right\}^2,$$
()
which can also be decomposed into contributions from the squared bias SB and the variance V:
$$\mathrm{MSE}(\widehat{P})=\left(E\widehat{P}P\right)^2+\left(E\left\{\widehat{P}^2\right\}E\left\{\widehat{P}\right\}^2\right)$$
()
$$=\mathrm{SB}(\widehat{P})+\mathrm{V}(\widehat{P}).$$
()
The MSE of estimates obtained with the MPL algorithm again varies roughly as a power law, $`\mathrm{MSE}(\widehat{h_4})(S/N)^{1.5}`$. The bias in the MPL estimates is always negative, i.e. in the direction of more Gaussian $`N(V)`$’s; for $`S/N20`$, this bias is a modest $`0.03`$, dropping to $`\begin{array}{c}<\hfill \\ \hfill \end{array}0.01`$ for S/N $`=100`$.
The FCQ estimates of $`h_4`$ show a considerably greater error, both in the bias and the variance. Two sets of FCQ estimates were made: first using the default value of the smoothing parameter returned by the code, $`W1.3`$; and second using the optimum value $`W_{opt}`$ that minimized the MSE of the $`h_4`$ estimates at each S/N. For the default value of $`W`$, the algorithm returns mean estimates of $`h_4`$ that lie in the range $`0.030.05`$ for all values of S/N $`10`$, an average error of $`70\%`$. However the optimum smoothing parameter for the recovery of $`h_4`$ was found to vary strongly with S/N, from $`W_{opt}0.5`$ for S/N $`=5`$ to $`W_{opt}2`$ for S/N $`=100`$ (Fig. B5c). Nevertheless a substantial bias remains when $`W_{opt}`$ is used, of order $`0.05`$ even for S/N $`=50100`$. These biases are larger than found above using a more Gaussian $`N(V)`$ with smaller $`h_4`$ and suggest that FCQ estimates of $`h_4`$ may be substantially biased even for $`\sigma _0`$ as large as $`100`$ km s<sup>-1</sup> when the true $`N(V)`$ is sufficiently non-Gaussian.
Top: STIS spectra of M32 at five positions along the slit. Solid curves are convolutions of the MPL-derived broadening functions $`\widehat{N}(V)`$ (Fig. The Nuclear Dynamics of M32. I.) with the stellar template. Bottom: Spectrum of HR7615, a K0III giant, the template spectrum. The vertical scale of the template spectrum is compressed with respect to that of the M32 spectra.
Broadening functions recovered from the central spectrum of M32 using the two spectral deconvolution algorithms discussed in the text. The degree of smoothing increases downward. Left column: FCQ. (a) $`W=2.4`$; (b) $`W=1.6`$; (c) $`W=1.2`$; (d) $`W=0.8`$. Right column: MPL. (e) $`\alpha =10^3`$; (f) $`\alpha =10^5`$; (g) $`\alpha =10^7`$; (h) $`\alpha =10^9`$. The MPL estimates tend toward a Gaussian for large $`\alpha `$ while the FCQ estimates become increasingly distorted as the smoothing is increased. This is the source of the greater bias in the FCQ estimates (although in practice smoothing parameters as small as $`W=0.8`$ would never be used).
The mean value of $`h_4`$ computed by the FCQ algorithm between $`0.5^{\prime \prime }`$ and $`2.0^{\prime \prime }`$, as a function of smoothing parameter $`W`$. Large values of $`W`$ correspond to small degrees of smoothing and hence to less biased estimates. These $`h_4`$ values are based on data that were heavily binned in radius in order to increase the S/N as much as possible.
Line-of-sight velocity distributions derived from the STIS M32 spectra using the MPL deconvolution algorithm. Note the sudden increase in the width of the broadening functions inside of $`0.2^{\prime \prime }`$. The LOSVDs are roughly antisymmetric about the center of M32, as expected for a relaxed system; the right column shows $`\widehat{N}(V)`$ averaged over the left and right sides, $`\frac{1}{2}[\widehat{N}(V,R)+\widehat{N}(V,R)]`$. The central LOSVD exhibits strong non-Gaussian wings, a likely consequence of high-velocity stars near the central black hole. The broadening functions at larger radii exhibit asymmetries suggestive of a second kinematic subcomponent which rotates with a velocity closer to the systemic velocity of M32.
STIS rotation curve for M32, derived from LOSVDs obtained using the MPL spectral deconvolution algorithm. Upper panel: filled circles: $`V_0`$, the parameter that measures the velocity shift of the Gaussian function that multiplies the Gauss-Hermite series. Open circles: $`V_0+\sqrt{3}\sigma _0h_3`$, an estimate of the true mean line-of-sight velocity. Lower panel: the Gauss-Hermite parameter $`h_3`$ that measures asymmetries in the LOSVDs. The mean velocity is smaller than $`|V_0|`$ due to the nonzero value of $`h_3`$, which in turn reflects asymmetries in the LOSVD’s (Fig. The Nuclear Dynamics of M32. I.).
STIS velocity dispersion profile for M32, derived from LOSVDs obtained using the MPL spectral deconvolution algorithm. Upper panel: filled circles: $`\sigma _0`$, the parameter that measures the dispersion of the Gaussian function that multiplies the Gauss-Hermite series. Open circles: $`\sigma _0(1+\sqrt{h_4})`$, an estimate of the true rms line-of-sight velocity. Lower panel: the Gauss-Hermite parameter $`h_4`$ that measures the amplitude of symmetric non-Gaussian distortions in the LOSVD. The velocity dispersion is generally greater than $`\sigma _0`$ due to the nonzero values of $`h_4`$. This difference is substantial in the inner $`0.2^{\prime \prime }`$ due to the strongly non-Gaussian wings of the central LOSVDs (Fig. The Nuclear Dynamics of M32. I.).
Comparison of $`V_0`$ and $`\sigma _0`$ derived from the M32 STIS data (filled circles) with earlier ground-based determinations. Squares: WHT measurements from van der Marel et al. (1994a). Triangles: CFHT measurements from Bender et al. (1996).
Comparison of $`V_0`$ and $`\sigma _0`$ derived from the M32 STIS data (filled circles) with FOS data of van der Marel et al. (1997) (open squares).
Reduced $`\chi ^2`$ contours describing the fit of the axisymmetric models described in the text to the observed, mean square line-of-sight velocity at points within the inner arc second of M32. The contour spacing is 0.5 and the innermost contour is at 0.75. The plus symbol marks the best-fit model.
Predicted kinematical profiles for three axisymmetric models with different black hole masses. (a) $`M_h=2.0\times 10^6_{}`$; (b) $`M_h=3.0\times 10^6_{}`$; (c) $`M_h=4.0\times 10^6_{}`$. Thin curves show the models as observed with infinite resolution; heavy curves are the models after convolution with the STIS PSF; open circles are the data points. For each $`M_h`$, the mass to light ratio $`M/L`$ and rotational parameter $`k`$ have been adjusted to optimize the fit.
The central M32 LOSVD (heavy line), symmetrized about $`V=0`$, compared to the LOSVD’s predicted by spherical nonrotating models with three black holes masses. Thin line: $`M_h=2.5\times 10^6_{}`$; dashed line: $`M_h=5.0\times 10^6_{}`$; thin line: $`M_h=10.\times 10^6_{}`$.
Fig. A1 —
Estimated Gauss-Hermite parameters $`\widehat{\sigma }`$ and $`\widehat{h_4}`$, derived by fitting the $`N(V)`$ of equation (A6) to the assumed form (A5), with $`j_{max}=4`$. Filled circles indicate the input values of $`\sigma _0`$ and $`h_4`$; these values are recovered only when the input $`N(V)`$ has $`h_6=0`$.
Fig. B1 —
Recovery of $`\sigma _0`$ via FCQ.
Fig. B2 —
Recovery of $`h_4`$ via FCQ. Circles: $`\sigma _0=40`$ km s<sup>-1</sup>. Triangles: $`\sigma _0=60`$ km s<sup>-1</sup>. Diamonds: $`\sigma _0=100`$ km s<sup>-1</sup>. There is a significant negative bias in the recovered values of $`h_4`$ when the velocity dispersion is less than about 100 km s<sup>-1</sup>.
Fig. B3 —
Mean estimates of $`N(V)`$ averaged over 100 random realizations of the observed spectrum, for S/N $`=\{5,20,100\}`$. The input $`N(V)`$ (Equation B2) is shown by the heavy curves.
Fig. B4 —
MISE and ISB of estimates of $`N(V)`$ obtained from the two deconvolution algorithms. The input LOSVD was a Lorentzian (Equation B2) with $`\sigma _0=108`$km s<sup>-1</sup> and $`h_4=0.15`$. Both MISE and ISB have been normalized as described in the text. Solid lines: MPL algorithm. Open circles: FCQ algorithm, using a fixed smoothing parameter $`W=1.3`$. Filled circles: FCQ algorithm, using the value $`W_{opt}`$ that minimizes the MISE of the estimated $`N(V)`$.
Fig. B5 —
MSE and bias in estimates of $`h_4`$ obtained from the two deconvolution algorithms. The input $`N(V)`$ was a Lorentzian (Equation B2) with $`\sigma _0=108`$km s<sup>-1</sup> and $`h_4=0.15`$. Solid lines: MPL algorithm. Open circles: FCQ algorithm, with $`W=1.2`$. Filled circles: FCQ algorithm, using the value $`W_{opt}`$ that minimizes the MISE of the estimate $`\widehat{h}_4`$. $`W_{opt}`$ is plotted vs. S/N in the bottom panel.
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# Eigenvector 1: An Optimal Correlation Space for Active Galactic Nuclei
## 1 Introduction
The search for correlations among observational parameters that describe AGN has been developing rapidly in the past few years. At the same time there has been a theoretical countercurrent that views all quasar spectra as remarkably similar. This erroneous view is due in part to the low resolution and s/n of much available spectral data. Blurred quasar spectra do look remarkably similar but data with s/n$`>`$20 in the continuum and resolution $``$5Å show striking differences from which a pattern is beginning to emerge. Some of the most promising correlations arose when high s/n optical spectra were published for the lower redshift (z$`\stackrel{<}{}`$0.5) sources in the Bright Quasar Survey (Boroson & Green 1992: BG92). A principal component analysis of the BG92 correlation matrix showed “eigenvector 1” correlations involving the widths and strengths of \[OIII\]$`\lambda `$5007, as well as broad H$`\beta `$ and Feii<sub>opt</sub> emission lines. More recently measures of X-ray luminosity, particularly the soft X-ray photon index, have emerged as a related part of the eigenvector 1 correlations (Wang, Brinkmann & Bergeron 1996).
We report a study of correlations involving the best available data samples for low redshift (z$`\stackrel{<}{}`$1.0) Active Galactic Nuclei (AGN). We find a correlation space where we are able to discriminate between all of the major forms of AGN phenomenology (e.g. broad and narrow Line Seyfert 1 galaxies, regular and broad absorption line (BAL) quasars, steep and flat spectrum broad line radio galaxies). We refer to it as the Eigenvector 1 (hereafter E1) space reflecting its partial roots in the BG92 study. The three principal “orthogonal” correlates involve measures of: 1) low ionization broad line width - full width at half maximum of H$`\beta `$ (FWHM H$`\beta _{\mathrm{BC}}`$), 2) ratio of line strengths - the continuum normalized ratio of the optical Feii<sub>opt</sub> and broad H$`\beta `$ emission line strengths (R<sub>FeII</sub>= W(Feii $`\lambda `$4570 blend)/W(H$`\beta _{\mathrm{BC}}`$) and 3) X-ray continuum strength - the soft X-ray photon index ($`\mathrm{\Gamma }_{\mathrm{soft}}`$). Parts of this correlation space have been discussed for almost ten years (e.g. Boroson & Green 1992 (PG92); Boller et al. 1996; Marziani et al. 1996 (MS96); Wang et al., 1996; Laor et al. 1997) but the pieces have not previously been united in this way. We present the basic E1 phenomenology in this paper. In a following paper we will consider the principal physical drivers of the E1 correlations.
## 2 The E1 Correlation Space
Figures 1 a,b,c show the 2D projections of the E1 correlation space. AGN plotted in the figures come from three data sources: 1) the BG92 sample of 87 mostly radio-quiet (RQ) AGN (16 radio-loud: RL), 2) the MS96 mixed sample of 21 RQ and 31 RL sources as well as 3) a sample of 24 new sources with matching HST UV archival and optical ground-based spectra (8 RL). The data sources combine the most complete sample of AGN with high quality spectral data (BG92) and two overlapping samples of comparable quality data with matching optical H$`\beta `$ and UV Civ$`\lambda `$1549. Soft X-ray photon indices are available for a large fraction of the above samples (Wang et al. 1996; Brinkmann et al. 1997; Siebert et al. 1998; Yuan et al. 1998). Values determined with hydrogen column density N<sub>H</sub> as a free parameter were preferred when available. The hard X-ray photon index provides an alternative, but less sensitive, E1 correlate (Piccinotti et al. 1982; Brandt et al. 1997). Our total sample includes 128 sources (45 RL) with optical measures, of which, 76 (39 RL) have UV measures. Mean error bars (2$`\sigma `$) are indicated with: i) average errors for sources in the middle of the diagrams for Feii $`\lambda `$4570 and H$`\beta `$ measures as well as ii) a median value for all $`\mathrm{\Gamma }_{\mathrm{soft}}`$ measures. See BG92 and MS96 for detailed discussion of reduction and analysis procedures.
RQ sources are indicated by solid symbols in the figures while open and crossed circles indicate flat (core-dominated) and steep (lobe-dominated) spectrum RL (as defined by Kellermann et al 1989) sources respectively. The majority of sources show a well defined correlation in Figures 1 a,b,c. One can make a first order interpretation of E1 correlations in terms of RL vs. RQ differences. Sources that fall in the large FWHM H$`\beta `$ and large R<sub>FeII</sub> region (like PG 0043+039; Turnshek et al 1994) are apparently very rare and pathological (possibly associated to mixed Starburst/AGN properties). Fig. 1 1 a,b,c indicate that these two populations show a clear separation with RQ having, for example, a mean FWHM(H$`\beta _{\mathrm{BC}}`$) $``$ 2300 km s<sup>-1</sup> less than RL sources (a few RL sources with FWHM(H$`\beta _{\mathrm{BC}}`$) $``$ 11–20000 km s<sup>-1</sup> are not shown in the figures). An alternative interpretation of E1 sees two populations of RQ AGN: 1) population A (filled boxes) which shows little overlap with the RL domain (65% of the BG92 RQ sample) and 2) population B (filled circles) that occupies essentially the same E1 domain as the RL sources ( 25% of the BG92 RQ sources). Fig. 1 1 d shows a plot of W(Feii $`\lambda `$4570) vs. FWHM(H$`\beta _{\mathrm{BC}}`$) for the BG92 sample that provided the motivation for the population A–B concept. Rather than a correlation this plot shows two disjoint populations of sources with no correlation within either population. Correlations appear when R<sub>FeII</sub> is used instead of W(Feii<sub>opt</sub>) or W(H$`\beta _{\mathrm{BC}}`$). We introduce the population A and B distinction for the sake of argument but suggest that it may have more fundamental significance. Table 1 presents mean parameter values (and associated sample standard deviations) in arbitrary bins of FWHM(H$`\beta _{\mathrm{BC}}`$) that allow one to compare RL vs. RQ as well as Pop. A vs. Pop. B sources.
It is not easy at this time to determine what role sample selection biases play in E1 but it is clear that BG92 selection techniques favored RQ sources with narrow Balmer line profiles. Recent work indicates that selection on the basis of soft X-ray excess favors them even more strongly (e.g. Moran et al. 1996; Grupe et al. 1999). At the same time core-dominated/flat spectrum RL sources are no doubt over-represented if their emission is beamed. This study makes no claim to completeness. Beyond BG92 our sample selection is driven by the availability of i) high s/n and resolution ground based data and ii) matching HST UV spectra. Table 2 presents the results of a correlation analysis for the three principal E1 and related parameters (the Pearson’s correlation coefficient r<sub>P</sub> is reported along with the probability P of a chance correlation in a sample of N objects). As the figures and Table 2 also suggest, the strongest E1 correlations are found for the RQ sources. The highest ranked E1 correlation ($`\mathrm{\Gamma }_{\mathrm{soft}}`$ vs. FWHM H$`\beta _{\mathrm{BC}}`$) is present only in the RQ sample which is not surprising since RL sources show no evidence for a soft X-ray excess. The nearly identical r<sub>P</sub> for the ALL RQ and Pop. A RQ samples suggest that the principal correlations are driven by Pop. A. RQ Pop. B and RL sources show the same E1 parameter domain and neither sample shows evidence for significant internal correlations. The latter result is due, at least in part, to: a) the absence of a soft X-ray excess in these sources and b) the weakness of Feii<sub>opt</sub> emission in these relatively broad-lined sources. If the minimum detectable Feii<sub>opt</sub> emission for sources with profile widths FWHM(Feii $`\lambda `$4570) $``$ FWHM(H$`\beta _{\mathrm{BC}}`$) $``$1000 km s<sup>-1</sup>is W(Feii $`\lambda `$4570) $``$10–15Å then we estimate that it will be $``$ 25 and 40 Å for FWHM(H$`\beta _{\mathrm{BC}}`$) $``$ 5000 and 10000 km s<sup>-1</sup>, respectively. The message here is that values of R<sub>FeII</sub>$`\stackrel{<}{}`$0.2 are very uncertain even with the best available data.
There are several lines of evidence that suggest continuum luminosity is uncorrelated with E1 parameter space. 1) Optical luminosity appears in the second orthogonal eigenvector identified in BG92, 2) BAL sources occupy a similar domain to the NLSy1 but they are “displaced” in: a) X-ray luminosity - BALs are X-ray quiet and b) optical luminosity where they are on average about 4$`\times `$ more luminous than the NLSy1 sources in BG92. 3) RL and Pop. B RQ sources occupy the same parameter domain but the RL sources show $`5\times `$ higher optical continuum luminosities which may in part reflect the presence of beamed sources in the RL sample. 4) We also find no evidence for a correlation between E1 coordinates and radio continuum luminosity for RQ sources (see Kellermann et al. 1989; Falcke et al. 1996). This last result tells us that RL sources are a distinct AGN population that shows fundamental differences in BLR structure and kinematics. E1 indicates that the weak radio emission from RQ sources is unrelated to the RL phenomenon.
E1 also shows evidence for a parameter space separation between steep and flat spectrum RL sources. This shows up in Table 2 as a possible RL correlation between FWHM H$`\beta _{\mathrm{BC}}`$ and R<sub>FeII</sub>. Steep spectrum RL sources represent the opposite extremum from NLSy1 while flat spectrum sources are more similar to RQ sources. The tendency for steep spectrum RL sources to show the broadest Balmer profiles and the weakest Feii<sub>opt</sub> emission is confirmed in two large (with some source overlap between themselves and our RL sample) surveys (Brotherton 1996; Corbin 1997). Siebert et al (1998) provide evidence that this separation may also be present in the X-ray spectral index. NLSy1 are a RQ and population A extremum with the narrowest Balmer line profiles, strong Feii<sub>opt</sub> emission and a strong soft X-ray excess. NLSy1 show a clear continuity and correlation with broader line Seyfert 1 galaxies in all measured parameters. This challenges the idea that they represent a unique or disjoint AGN population. BAL quasars (4 RQ BAL in the BG92 sample) occupy an E1 domain that is similar to the NLSy1. Other (BAL) studies (e.g. Boroson & Myers 1992) have also suggested that BAL quasars show Balmer line FWHM$`\stackrel{<}{}`$ 3000 km s<sup>-1</sup>and moderate to strong R<sub>FeII</sub>measures.
Much less clear is whether RQ population B sources represent a disjoint AGN population or show a smooth continuation of the population A correlations. Figure 1 b for example shows a breakdown of the $`\mathrm{\Gamma }_{\mathrm{soft}}`$ – FWHM(H$`\beta `$) correlation at the nominal boundary between population A and B. Whatever the relation between the two RQ populations, their line profiles show striking differences. Figure 2 shows a comparison of the H$`\beta _{\mathrm{BC}}`$ and Civ$`\lambda `$1549<sub>BC</sub> line profiles for prototype NLSy1 source I Zw 1 and NGC 5548 which is a typical broad line Seyfert galaxy. While differences in the Balmer profiles are striking, the most impressive difference is the apparent kinematic decoupling of the Civ$`\lambda `$1549 and H$`\beta `$ profiles in I Zw 1. Our previous work (MS96) suggested that properly NLR corrected Civ$`\lambda `$1549 in RQ sources (see Sulentic & Marziani 1999) is always blueshifted relative to H$`\beta `$ and the AGN rest frame. Fig. 3 shows the Civ$`\lambda `$1549<sub>BC</sub> line shift vs. FWHM(H$`\beta _{\mathrm{BC}}`$) diagram. We have normalized the Civ$`\lambda `$1549 shift by W(Civ$`\lambda `$1549) in this plot because we see a complementary trend in E1 for W(Civ$`\lambda `$1549) to be smallest in the NLSy1 population. Figure 3 confirms that essentially all RQ sources show a Civ$`\lambda `$1549 blueshift while RL sources show equal red and blueshifts with amplitudes generally less than $`\pm `$10<sup>3</sup> km s<sup>-1</sup>. Correlations of Civ$`\lambda `$1549 shift with R<sub>FeII</sub> and $`\mathrm{\Gamma }_{\mathrm{soft}}`$ (see Table 2) indicate that it is likely to be an important E1 correlate. We will explore it further in succeeding papers.
## 3 From Observational to Physical Parameters
E1 correlates AGN spectroscopic data in a way that removes much of the apparent “randomness” of line properties. It also redefines input parameters for photoionization and kinematical models. We have only begun to explore the E1 parameter space and and the physics that drives it. Accumulating evidence from numerous studies suggests that the correlations in E1 involve at least two principal independent parameters: the AGN luminosity to black hole mass ratio and the source orientation. The role of Fe abundance, disk magnetic fields and black hole angular momentum are beyond the scope of this introduction to E1.
Available evidence suggests that the central source luminosity to black hole mass ratio L/M may be systematically higher in Pop. A sources (i.e. a higher accretion rate). A role for L/M is supported by: a) the existence of a soft X-ray excess, particularly in NLSy1, that has been related to a higher accretion rate (Pounds et al. 1995) and b) the thermal signature of the disk itself (Page et al. 1999; Puchnarewicz et al. 1998). If the soft X-ray excess is the high energy end of a thermal signature of the accretion disk then a steep rising blue optical continuum (in low z sources) can be interpreted as the low energy end of the “big blue bump”. This would explain why BG92, with quasars selected on the basis of U–B colors, would favor detection of NLSy1. Civ$`\lambda `$1549 shows the largest blueshifts in NLSy1 sources, and the lowest Civ$`\lambda `$1549 W values, both consistent with the idea that the line emission arises in a disk wind (e.g. Murray & Chiang 1997; Bottorff et al. 1997). The observations are consistent with the idea that L/M is only high enough in RQ sources (especially Pop. A) to trigger a radiation pressure driven outflow. The Balmer profiles in Pop. A sources are most easily interpreted as arising in an optically thick illuminated accretion disk (see e.g. Sulentic et al. 1998). RL sources show Balmer lines with both large red and blue profile shifts.
RL sources show no evidence for a soft X-ray excess which may reflect a real absence of this component or that it is dominated by a relatively hard component related to X-ray emission related to the radio loudness. We favor the former interpretation because most RQ population B sources also show a soft X-ray deficit. The absence of a soft X-ray component in RL+RQ pop B sources coupled with the absence of the CIV blueshift are consistent with weak or absent disk structure. The broad lines may arise in a biconical structure in many or all of these sources (Marziani et al. 1993; Sulentic et al. 1995). Evidence for a bicone origin of the BLR includes: a) sources with double peaked broad lines (e.g. Sulentic et al. 1995; Eracleous & Halpern 1998), b) sources with single peaked broad lines showing both large red and blue displacements (e.g. Marziani et al. 1983; Halpern et al. 1998), c) sources with a transient BLR (e.g. Storchi-Bergmann et al. 1993; Ulrich 2000) and spectropolarimetry of a) and b) sources (Corbett et al. 1998).
The E1 parameter space is measuring aspects of both the geometry and kinematic of the broad line region in AGN. The strength of the correlations and reasonable orthogonality of the parameters suggest that a better diagnostic space for AGN is unlikely to be found. Even in the preliminary presentation, E1 allows us to resolve some AGN conundrums. 1) RL sources are found to be an AGN population with fundamental geometrical and kinematic difference from the RQ majority. 2) NLSy1 are found to be an extremum of the RQ population rather than a pathological or disjoint AGN population. 3) If the width of the Balmer lines bear any signature of source orientation then some NLSy1 (at least, I Zw 1) are seen at or near pole-on (face-on accretion disk) orientation, while some BAL QSOs are likely to be misaligned NLSy1. 4) Population B RQ quasars with Balmer lines broader than $``$ 4000 km s<sup>-1</sup>may represent a distinct class (RL pre/post cursors) from narrower population A RQ sources. 5) Civ$`\lambda `$1549 line shifts indicate the presence in Pop. A (and possibly Pop. B) RQ sources of an ubiquitous disk outflow/wind.
PM acknowledges financial support from the Italian Ministry of University and Scientific and Technological Research (MURST) through grant Cofin 98-02-32 and from the IA-UNAM where part of this work was done.
## 4 References
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# References
## Abstract
Feynman’s path integral is generalized to quantum mechanics on $`p`$-adic space and time. Such $`p`$-adic path integral is analytically evaluated for quadratic Lagrangians. Obtained result has the same form as that one in ordinary quantum mechanics.
On $`p`$-adic path integral
Branko Dragovich
Institute of Physics, P.O.Box 57, 11001 Belgrade, Yugoslavia
Dedicated to the memory
of N.N.Bogolyubov
1. It is well known that dynamical evolution of any one-dimensional quantum-mechanical system, described by a wave function $`\mathrm{\Psi }(x,t)`$, is given by
$$\mathrm{\Psi }(x^{\prime \prime },t^{\prime \prime })=𝒦(x^{\prime \prime },t^{\prime \prime };x^{},t^{})\mathrm{\Psi }(x^{},t^{})𝑑x^{},$$
$`(1)`$
where $`𝒦(x^{\prime \prime },t^{\prime \prime };x^{},t^{})`$ is the kernel of the corresponding unitary operator acting as follows:
$$\mathrm{\Psi }(t^{\prime \prime })=𝒰(t^{\prime \prime },t^{})\mathrm{\Psi }(t^{}).$$
$`(2)`$
$`𝒦(x^{\prime \prime },t^{\prime \prime };x^{},t^{})`$ is also called Green’s function, or the quantum-mechanical propagator, and the probability amplitude to go a particle from a point $`(x^{},t^{})`$ to a point $`(x^{\prime \prime },t^{\prime \prime })`$. One can easily deduce the following three general properties:
$$𝒦(x^{\prime \prime },t^{\prime \prime };x,t)𝒦(x,t;x^{},t^{})𝑑x=𝒦(x^{\prime \prime },t^{\prime \prime };x^{}t^{}),$$
$`(3)`$
$$𝒦^{}(x^{\prime \prime },t^{\prime \prime };x^{},t^{})𝒦(x^{\prime \prime },t^{\prime \prime };x,t^{})𝑑x^{\prime \prime }=\delta (x^{}x),$$
$`(4)`$
$$𝒦(x^{\prime \prime },t^{\prime \prime };x^{},t^{\prime \prime })=lim_{t^{}t^{\prime \prime }}𝒦(x^{\prime \prime },t^{\prime \prime };x^{},t^{})=\delta (x^{\prime \prime }x^{}).$$
$`(5)`$
Since all information on quantum dynamics can be deduced from the propagator $`𝒦(x^{\prime \prime },t^{\prime \prime };x^{},t^{})`$ it can be regarded as the basic ingredient of quantum theory. In Feynman’s formulation of quantum mechanics, $`𝒦(x^{\prime \prime },t^{\prime \prime };x^{},t^{})`$ was postulated to be the path integral
$$𝒦(x^{\prime \prime },t^{\prime \prime };x^{},t^{})=\mathrm{exp}\left(\frac{2\pi i}{h}_t^{}^{t^{\prime \prime }}L(\dot{q},q,t)𝑑t\right)𝒟q,$$
$`(6)`$
where $`x^{\prime \prime }=q(t^{\prime \prime })`$ and $`x^{}=q(t^{})`$, and $`h`$ is the Planck constant.
In its original form, the path integral (6) is the limit of the corresponding multiple integral of $`n1`$ variables $`q_i=q(t_i),(i=1,2,\mathrm{},n1),`$ when $`n\mathrm{}`$. For the half of century of its history, the path integral has been a subject of permanent interest in theoretical and mathematical physics. At present days (see, e.g. ) it is one of the most profound and promising approaches to foundations of quantum theory (in particular, quantum field theory and superstring theory). Feynman’s path integral is inevitable in formulation of $`p`$-adic and adelic quantum mechanics. It is worth noting that just Feynman’s path integral approach enables natural foundation of quantum theory on $`p`$-adic and adelic spaces.
2. Recall that the set of rational numbers $`Q`$ plays an important role in mathematics as well as in physics. From algebraic point of view, $`Q`$ is the simplest number field of characteristic $`0`$. The usual absolute value and $`p`$-adic valuation ($`p`$ is any of prime numbers) exhaust all possible non-trivial norms on $`Q`$ . Completion of $`Q`$ with respect to metrics induced by these norms leads to the field of real numbers $`R`$ and the fields of $`p`$-adic numbers $`Q_p`$, $`(p=2,3,5,\mathrm{})`$. Thus $`Q`$ is dense in $`R`$ and all $`Q_p`$. From physical point of view, all numerical results of measurements are rational numbers. However, theoretical models of physical systems are traditionally constructed using real and complex numbers. One can ask the following question: Why real (and complex) numbers are so good in description of usual physical phenomena, and, is there any aspect of physical reality which has to be described by $`p`$-adic numbers. Construction of $`p`$-adic models and their appropriate interpretation can gradually give answer to this question. Since 1987, there have been many publications (for a review, see, e.g. ) on possible applications of $`p`$-adic numbers in modern theoretical and mathematical physics. For a systematic approach to this subject, $`p`$-adic and adelic quantum mechanics have been formulated.
Recall also that any $`p`$-adic number $`xQ_p`$ can be presented as the following infinite expansion
$$x=p^\nu (x_0+x_1p+x_2p^2+\mathrm{}),\nu Z,$$
where $`x_i=0,1,\mathrm{},p1`$ are digits. We will use the Gauss integral
$$_{Q_p}\chi _p(\alpha x^2+\beta x)𝑑x=\lambda _p(\alpha )2\alpha _p^{\frac{1}{2}}\chi _p\left(\frac{\beta ^2}{4\alpha }\right),\alpha 0,$$
where $`\chi _p(a)=\mathrm{exp}(2\pi i\{a\}_p)`$ is the additive character, and $`\{a\}_p`$ is the fractional part of $`aQ_p`$. $`\lambda _p(x)`$ is a complex-valued arithmetic function (for a definition, see ) with the following properties:
$$\lambda _p(0)=1,\lambda _p(a^2x)=\lambda _p(x),\lambda _p(x)\lambda _p(y)=\lambda _p(x+y)\lambda _p(x^1+y^1),\lambda _p^{}(x)\lambda _p(x)=1.$$
3. $`p`$-Adic quantum mechanics, we are interested in, contains complex-valued functions of $`p`$-adic arguments. There is not the corresponding Schr$`\ddot{o}`$dinger equation, but Feynman’s path integral approach seems to be quite natural. Feynman’s path integral for $`p`$-adic propagator $`𝒦_p(x^{\prime \prime },t^{\prime \prime };x^{},t^{})`$, where $`𝒦_p`$ is complex-valued and $`x^{\prime \prime },x^{},t^{\prime \prime },t^{}`$ are $`p`$-adic variables, is a direct $`p`$-adic generalization of (6), i.e.
$$𝒦_p(x^{\prime \prime },t^{\prime \prime };x^{},t^{})=\chi _p\left(\frac{1}{h}_t^{}^{t^{\prime \prime }}L(\dot{q},q,t)𝑑t\right)𝒟q,$$
$`(7)`$
where $`\chi _p(a)`$ is $`p`$-adic additive character. The Planck constant $`h`$ in (6) and (7) is the same rational number. Integral $`_t^{}^{t^{\prime \prime }}L(\dot{q},q,t)𝑑t`$ we consider as the difference of antiderivative (without pseudoconstants) of $`L(\dot{q},q,t)`$ in final $`(t^{\prime \prime })`$ and initial $`(t^{})`$ times. $`𝒟q=_{i=1}^{n1}dq(t_i)`$, where $`dq(t_i)`$ is the $`p`$-adic Haar measure. Thus, $`p`$-adic path integral is the limit of the multiple Haar integral when $`n\mathrm{}`$. To calculate (7) in this way one has to introduce some order on $`tQ_p`$, and it is successfully done in Ref. . On previous investigations of $`p`$-adic path integral one can see . Our main task here is derivation of the exact result for $`p`$-adic Feynman’s path integral (7) for the general case of Lagrangians $`L(\dot{q},q,t)`$, which are quadratic polynomials in $`\dot{q}`$ and $`q`$, without making time discretization.
A general quadratic Lagrangian can be written as follows:
$$L(\dot{q},q,t)=\frac{1}{2}\frac{^2L_0}{\dot{q}^2}\dot{q}^2+\frac{L_0}{\dot{q}}\dot{q}+\frac{^2L_0}{\dot{q}q}\dot{q}q+L_0+\frac{L_0}{q}q+\frac{1}{2}\frac{^2L_0}{q^2}q^2,$$
$`(8)`$
where index $`0`$ denotes that the Taylor expansion of $`L(\dot{q},q,t)`$ is around $`\dot{q}=q=0`$. The Euler-Lagrange equation of motion is
$$\frac{^2L_0}{\dot{q}^2}\ddot{q}+\frac{d}{dt}\left(\frac{^2L_0}{\dot{q}^2}\right)\dot{q}+\left[\frac{d}{dt}\left(\frac{^2L_0}{\dot{q}q}\right)\frac{^2L_0}{q^2}\right]q=\frac{L_0}{q}\frac{d}{dt}\left(\frac{L_0}{\dot{q}}\right).$$
$`(9)`$
General solution of (9) is
$$qx(t)=C_1x_1(t)+C_2x_2(t)+w(t),$$
$`(10)`$
where $`x_1(t)`$ and $`x_2(t)`$ are two linearly independent solutions of the corresponding homogeneous equation, and $`w(t)`$ is a particular solution of the complete equation (9). Note that $`x(t)`$ denotes the classical trajectory. Imposing the boundary conditions $`x^{}=x(t^{})`$ and $`x^{\prime \prime }=x(t^{\prime \prime })`$, constants of integration $`C_1`$ and $`C_2`$ become:
$$C_1C_1(t^{\prime \prime },t^{})=\frac{(x^{}w^{})x_2^{\prime \prime }(x^{\prime \prime }w^{\prime \prime })x_2^{}}{x_2^{\prime \prime }x_1^{}x_1^{\prime \prime }x_2^{}},$$
$`(11a)`$
$$C_2C_2(t^{\prime \prime },t^{})=\frac{(x^{\prime \prime }w^{\prime \prime })x_1^{}(x^{}w^{})x_1^{\prime \prime }}{x_2^{\prime \prime }x_1^{}x_1^{\prime \prime }x_2^{}}.$$
$`(11b)`$
Since $`C_1(t^{\prime \prime },t^{})`$ and $`C_2(t^{\prime \prime },t^{})`$ are linear in $`x^{\prime \prime }`$ and $`x^{}`$, the corresponding classical action $`\overline{S}(x^{\prime \prime },t^{\prime \prime };x^{},t^{})=_t^{}^{t^{\prime \prime }}L(\dot{x},x,t)𝑑t`$ is quadratic in $`x^{\prime \prime }`$ and $`x^{}`$. Note that the above expressions have the same form in $`R`$ and $`Q_p`$.
Quantum fluctuations lead to deviations of classical trajectory and any quantum path may be presented as $`q(t)=x(t)+y(t)`$, where $`y^{}=y(t^{})=0`$ and $`y^{\prime \prime }=y(t^{\prime \prime })=0`$. The corresponding Taylor expansion of $`S[q]`$ around classical path $`x(t)`$ is
$$S[q]=S[x+y]=S[x]+\frac{1}{2!}\delta ^2S[x]=S[x]+\frac{1}{2}_t^{}^{t^{\prime \prime }}\left(\dot{y}\frac{}{\dot{q}}+y\frac{}{q}\right)^{(2)}L(\dot{q},q,t)𝑑t,$$
$`(12)`$
where we used $`\delta S[x]=0`$. We have now
$$𝒦_p(x^{\prime \prime },t^{\prime \prime };x^{},t^{})=\chi _p\left(\frac{1}{h}S[x]\right)\chi _p\left(\frac{1}{2h}_t^{}^{t^{\prime \prime }}\left(\dot{y}\frac{}{\dot{q}}+y\frac{}{q}\right)^{(2)}L(\dot{q},q,t)𝑑t\right)𝒟y$$
$`(13)`$
with $`y^{\prime \prime }=y^{}=0`$ and $`S[x]=\overline{S}(x^{\prime \prime },t^{\prime \prime };x^{},t^{})`$.
Note that $`𝒦_p(x^{\prime \prime },t^{\prime \prime };x^{},t^{})`$ has the form
$$𝒦_p(x^{\prime \prime },t^{\prime \prime };x^{},t^{})=N_p(t^{\prime \prime },t^{})\chi _p\left(\frac{1}{h}\overline{S}(x^{\prime \prime },t^{\prime \prime };x^{},t^{})\right),$$
$`(14)`$
where $`N_p(t^{\prime \prime },t^{})`$ does not depend on end points $`x^{\prime \prime }`$ and $`x^{}`$. To calculate $`N_p(t^{\prime \prime },t^{})`$ we use conditions (3) and (4). Substituting $`𝒦_p(x^{\prime \prime },t^{\prime \prime };x^{},t^{})`$ into (4) we obtain (for details, see ):
$$N_p(t^{\prime \prime },t^{})=\left|\frac{1}{h}\frac{^2\overline{S}}{x^{\prime \prime }x^{}}(x^{\prime \prime },t^{\prime \prime };x^{},t^{})\right|_p^{\frac{1}{2}}A_p(t^{\prime \prime },t^{}),$$
$`(15)`$
where $`|A_p(t^{\prime \prime },t^{})|_{\mathrm{}}=1`$, ($`||_p`$ and $`||_{\mathrm{}}`$ denote $`p`$-adic and absolute value, respectively). Replacing (15) in equation (3) we get conditions:
$$A_p(t^{\prime \prime },t)A_p(t,t^{})\lambda _p(\alpha )=A_p(t^{\prime \prime },t^{}),$$
$`(16)`$
$$\left|\frac{1}{h}\frac{^2\overline{S}}{x^{\prime \prime }x}(x^{\prime \prime },t^{\prime \prime };x,t)\right|_p^{\frac{1}{2}}\left|\frac{1}{h}\frac{^2\overline{S}}{xx^{}}(x,t;x^{},t^{})\right|_p^{\frac{1}{2}}|2\alpha |_p^{\frac{1}{2}}=\left|\frac{1}{h}\frac{^2\overline{S}}{x^{\prime \prime }x^{}}(x^{\prime \prime },t^{\prime \prime };x^{},t^{})\right|_p^{\frac{1}{2}},$$
$`(17)`$
where
$$\alpha =\frac{1}{2h}\left[\frac{^2\overline{S}}{x^2}(x^{\prime \prime },t^{\prime \prime };x,t)+\frac{^2\overline{S}}{x^2}(x,t;x^{},t^{})\right].$$
$`(18)`$
Analysing the above formulae we obtain
$$A_p(t^{\prime \prime },t^{})=\lambda _p\left(\frac{1}{2h}\frac{^2\overline{S}}{x^{\prime \prime }x^{}}(x^{\prime \prime },t^{\prime \prime };x^{},t^{})\right).$$
$`(19)`$
For details of a quite rigorous derivation of (19), see .
As the final result we have
$$𝒦_p(x^{\prime \prime },t^{\prime \prime };x^{},t^{})=\lambda _p\left(\frac{1}{2h}\frac{^2\overline{S}}{x^{\prime \prime }x^{}}\right)\left|\frac{1}{h}\frac{^2\overline{S}}{x^{\prime \prime }x^{}}\right|_p^{\frac{1}{2}}\chi _p\left(\frac{1}{h}\overline{S}(x^{\prime \prime },t^{\prime \prime };x^{},t^{})\right)$$
$`(20)`$
which is the $`p`$-adic Feynman path integral for quadratic Lagrangians. The corresponding path integral of ordinary quantum mechanics can be transformed into the same form as (20), i.e. in such case index $`p`$ is replaced by index $`\mathrm{}`$. This supports Volovich’s conjecture that fundamental physical laws should be invariant under interchange of number fields $`Q_p`$ and $`R`$.
Acknowledgement. Author wishes to thank the organizers of the Bogolyubov Conference: Problems of Theoretical and Mathematical Physics, for invitation to participate in Moscow and Dubna parts of the Conference.
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# 0^♯ and Elementary End Extensions of 𝑉_𝜅
## 1. Introduction
In this paper we consider the question of existence of elementary end extensions of models of the form $`(V_\kappa ,)`$.
###### Definition 1.1.
* Let $`(𝔼_M,_M)`$ denote the structure of all non-trivial elementary end extensions of M, with $`A_MB`$ iff B is an elementary end extension of A.
* Let $`(𝔼_M^{\text{wf}},_M)`$ denote the structure of all non-trivial well founded elementary end extensions of M, with $`A_MB`$ iff B is an elementary end extension of A.
Several results regarding the existence of elements in $`𝔼_M`$ were proved by Keisler, Silver and Morley.
###### Theorem 1.2 (Keisler, Morley).
Let M be a model of ZFC, $`cof(On^M)=\omega `$. Then $`𝔼_M\mathrm{}`$.
###### Theorem 1.3 (Keisler,Silver).
Let $`M=(V_\kappa ,)`$ be a model of ZFC, where $`\kappa `$ is weakly compact cardinal. Then for every $`SM`$ $`𝔼_{(V_\kappa ,,S)}^{\text{wf}}\mathrm{}`$.
Villaveces , has proved several other results regarding the existence of elementary end extensions of $`V_\kappa `$.
###### Theorem 1.4 (Villaveces).
The theory “ZFC + GCH + $`\lambda (\lambda `$ measurable) + $`\kappa [\kappa `$ inaccessible not weakly compact $``$ transitive $`M_\kappa ZFC`$ such that $`o(M)=\kappa `$ and $`𝔼_M^{\text{wf}}=\mathrm{}`$\]” is consistent relative to the theory “ZFC + $`\lambda (\lambda `$ measurable) + the weakly compact cardinals are cofinal in On”.
He also proved that the property $`𝔼_{V_\kappa }^{\text{wf}}\mathrm{}`$, is not preserved in certain generic extensions by destroying a weakly compact cardinal. In this paper we consider the problem of downwards absoluteness of the existence of well founded elementary end extensions of $`V_\kappa `$. We prove the following :
###### Theorem 1.5.
If $`0^{\mathrm{}}`$ exists then for every cardinal $`\kappa `$ there is an inner model M such that
(1.1)
$$M𝔼_{V_\kappa }=\mathrm{}.$$
In particular weak compactness is never downwards absolute, once we have $`0^{\mathrm{}}`$ in the universe. On the other hand we will prove that any cardinal with uncountable cofinality is Mahlo in any strict inner model of $`L[0^{\mathrm{}}]`$. I would like to thank the referee for pointing out an inaccuracy in the formulation of lemma 3.1 and for asking the question at the end of the paper.
## 2. Main Theorem
In this section we prove theorem 1.5. Let $`\kappa `$ be a cardinal. Since we assume that $`0^{\mathrm{}}`$ exists we can construct our model inside the inner model $`L[0^{\mathrm{}}]`$ . Note that since $`\kappa `$ is a cardinal in $`V`$ it remains a cardinal in $`L[0^{\mathrm{}}]`$, and hence it is weakly compact in $`L`$. Our model will be a generic extension of $`L`$, such that we will be able to construct a generic object inside $`L[0^{\mathrm{}}]`$. The basic idea will be to construct a generic Suslin tree and then to code it. For the construction of the Suslin tree we will follow Kunen’s construction , while the coding will use Levy collapse of certain $`L`$ cardinals. Then we will obtain the generic filter inside $`L[0^{\mathrm{}}]`$.
The following theorem by Kunen gives us the forcing for generating the Suslin tree.
###### Theorem 2.1.
Let $`\kappa `$ be a weakly compact cardinal and $`P_\kappa `$ be the forcing for adding a Cohen subset to $`\kappa `$. Then $`P_\kappa R_\kappa T_\kappa `$, where $`R_\kappa `$ is a forcing that adds a Suslin tree $`T_\kappa `$ to $`\kappa `$, and $`T_\kappa `$ is the forcing defined by the tree.
Let P be the reverse Easton iteration for adding a Cohen subset to each inaccessible, defined by :
###### Definition 2.2.
(2.1)
$$\text{P}=(P_\alpha ,Q_\alpha |\alpha On),$$
where
* $`P_0=\mathrm{}`$.
* If $`\alpha `$ is not inaccessible then $`P_\alpha Q_\alpha =\mathrm{}`$
* If $`\alpha `$ is inaccessible then $`Q_\alpha `$ is a $`P_\alpha `$ name for a partial order adding a Cohen subset to $`\alpha `$ i.e. $`P_\alpha Q_\alpha =(2^{<\alpha },)`$.
Direct limits are taken at inaccessible limits of inaccessibles and inverse limits otherwise.
Solovay (see M. Stanley ) proved that the reverse Easton support iteration for adding Cohen subsets to every $`L`$ inaccessible has a generic filter in $`L[0^{\mathrm{}}]`$, and therefore our iteration up to $`\kappa `$ has a generic filter as well.
Let $`\text{G}=G_\alpha |\alpha \kappa `$ be P generic. By Kunen’s theorem we can interpret $`G_\kappa `$ as a pair $`G_\kappa =T_\kappa ,b_\kappa `$ where $`T_\kappa `$ is a $`\kappa `$ Suslin tree and $`b_\kappa `$ is a branch through $`T_\kappa `$.
Next we define the forcing used to code the tree $`T_\kappa `$. Let S be the Easton supported product of collapsing of $`\alpha ^{+3}`$ to $`\alpha ^{+2}`$ defined inside $`L`$.
(2.2)
$$\text{S}=\{S_\alpha :\alpha \text{ is inaccesible }\}$$
where $`S_\alpha =Coll(\alpha ^{+2},\alpha ^{+3})`$.
###### Proposition 2.3.
There is a $`\text{P}\times \text{S}`$ generic over $`L`$, inside $`L[0^{\mathrm{}}]`$.
###### Proof.
The method of proof of this lemma is almost identical to the proof of M. Stanley of Solovay’s theorem that there exists a P generic filter over $`L`$ inside $`L[0^{\mathrm{}}]`$. We shall build the generic filter by induction on the Silver indiscernibles. The main point will be taking care that at limits the generic filter will be the direct limit of the previously built generic filters.
Let $`i_\alpha :\alpha <\kappa `$ be an increasing enumeration of the indiscernibles below $`\kappa `$. For any indiscernible $`\lambda `$ the forcing can be factored as
(2.3)
$$\text{P}\times \text{S}=\left(\text{P}^{\lambda +1}\text{P}_{\lambda +1}\right)\times \left(\text{S}^\lambda \times \text{S}_\lambda \right)$$
where $`\text{P}^\lambda `$ is the iteration up to $`\lambda `$, and $`\text{P}_\lambda `$ is the iteration from $`\lambda `$ upwards. For each $`\alpha `$ we shall define $`(G^{i_\alpha },H^{i_\alpha })`$, and then define $`(G^{i_\alpha +1},H^{i_\alpha +1})`$ such that $`G^{i_\alpha +1}\times H^{i_\alpha +1}`$ is $`\left(\text{P}^{i_\alpha }Q_{i_\alpha }\right)\times \text{S}^{i_\alpha +1}`$ generic over $`L[G^{i_\alpha }\times H^{i_\alpha }]`$.
$`i_0`$ or $`i_{\alpha +1}`$.
We have that in $`L`$ for every indiscernible $`\lambda `$ both $`\text{P}_{\lambda +1}^{}`$ and $`\text{S}_\lambda `$ are $`\lambda ^{++}`$ closed, where
(2.4)
$$\text{P}_{\lambda +1}^{}=\{\tau :\tau \text{ is a name and }_{\text{P}^{\lambda +1}}\tau \stackrel{~}{𝐏}_{\lambda +1}\}$$
is the term forcing for $`\text{P}_{\lambda +1}`$. Hence $`\text{P}_{\lambda +1}\times \text{S}_\lambda `$ is $`\lambda ^+`$-distributive over $`L^{\text{P}^{\lambda +1}\times \text{S}^\lambda }`$, since $`\text{P}^{\lambda +1}\times \text{S}^\lambda `$ is obviously $`\lambda ^+`$-c.c.
By the same argument $`\text{P}_{i_\alpha +1}^{i_{\alpha +1}}\times \text{S}_{i_\alpha }^{i_{\alpha +1}}`$ is also $`i_\alpha ^+`$ distributive. Let
(2.5)
$$M=L^{\text{P}^{i_\alpha +1}\times \text{S}^{i_\alpha +1}}.$$
Note that each $`L`$ name for dense subset of $`\text{P}_{i_\alpha +1}^{i_{\alpha +1}}\times \text{S}_{i_\alpha }^{i_{\alpha +1}}`$ in $`M`$, belongs to the Skolem hull of the ordinals up to $`i_\alpha `$ and finitely many indiscernibles above $`i_{\alpha +1}`$, say $`\{i_{\alpha +1},\mathrm{},i_{\alpha +n}\}`$. Hence in $`L[0^{\mathrm{}}]`$ we can represent the dense subsets of $`\text{P}_{i_\alpha +1}^{i_{\alpha +1}}\times \text{S}_{i_\alpha +1}^{i_{\alpha +1}}`$ in $`M`$, by a countable union of families of dense subsets each of size $`i_\alpha `$. Now using the $`i_\alpha ^+`$ distributivity we can meet each of these dense subsets. To ensure downwards compatibility we also demand that $`(G^{i_\alpha +1},H^{i_\alpha +1})`$ extends $`(G^{i_\alpha },H^{i_\alpha })`$. Finally use the same distributivity argument to define a generic filter $`G(i_{\alpha +1})`$ for $`Q_{i_{\alpha +1}}`$ over $`L^{\left(\text{P}^{i_{\alpha +1}}\times \text{S}^{i_{\alpha +1}}\right)}`$. Again in order to ensure extension we demand that $`G(i_{\alpha +1})`$ extends $`G(i_\alpha )`$, by putting a condition forcing it into the generic. Since S is not active at these stages and using the fact that P is a reverse Easton iteration this is possible.
$`i_\alpha `$ for $`\alpha `$ limit.
We have built generic objects $`G^i\times H^i:i<\alpha `$ for the product up to $`\alpha `$. Now we would like to build a generic filter for $`\text{P}^{i_\alpha }\times \text{S}^{i^\alpha }`$. Note that since $`i_\alpha `$ is Mahlo in $`L`$ we take direct limit. Moreover $`\text{P}^{i_\alpha }\times \text{S}^{i_\alpha }`$ is $`i_\alpha `$c.c. Define $`G^{i_\alpha },H^{i_\alpha }`$ by
(2.6)
$$pG^{i_\alpha }\text{ iff }\gamma <i_\alpha p\gamma G^{i_\gamma }.$$
(2.7)
$$sH^{i_\alpha }\text{ iff }\gamma <i_\alpha s\gamma H^{i_\gamma }.$$
We prove that $`G^{i_\alpha }\times H^{i_\alpha }`$ is $`\text{P}^{i_\alpha }\times \text{S}^{i_\alpha }`$ generic over $`L`$. Suppose that
$`D\text{P}^{i_\alpha }\times \text{S}^{i_\alpha }`$ is dense open. $`D`$ belongs to the Skolem hull of finitely many ordinals below $`i_\alpha `$ $`\text{a}=\gamma _1,\mathrm{},\gamma _n`$ and finitely many indiscernibles above $`\alpha `$ say $`\text{i}_n=i_{\alpha +1},\mathrm{},i_{\alpha +n}`$. Let $`sup(\text{a})<i_\beta <i_\alpha `$. Define an elementary embedding $`j:LL`$ by
(2.8)
$$j(i_\gamma )=\{\begin{array}{cc}i_\gamma & \text{ if }\gamma <\beta \\ i_{\alpha +\delta }& \text{ if }\gamma =\beta +\delta ,0\delta \end{array}$$
Obviously $`D\text{rng}j`$, and $`j^1(D)`$ is dense open in $`\text{P}^\beta \times \text{S}^\beta `$. Let $`(p^{},q^{})j^1(D)\left(\text{P}^\beta \times \text{S}^\beta \right)`$. Since both $`p^{},q^{}`$ are trivial on an end segment we obtain that
(2.9)
$$j((p^{},q^{}))=(p,q)^{}\mathrm{}^{Q_\gamma \times S_\gamma }:\beta \gamma <\alpha .$$
Hence by our choice of $`(G^{i_\alpha },H^{i_\alpha })`$ we obtain that $`j((p^{},q^{}))(G^{i_\alpha },H^{i_\alpha })`$.
Finally we prove that we can find a generic object $`G(i_\alpha )`$ for $`Q_{i_\alpha }`$ over $`L^{\left(G^{i_\alpha }\times H^{i_\alpha }\right)}`$. Define
(2.10)
$$G(i_\alpha )=_{\beta <\alpha }G(i_\beta ).$$
Let $`D`$ be a dense subset of $`Q_{i_\alpha }`$ in $`L^{G^{i_\alpha }\times H^{i_\alpha }}`$. Let $`\stackrel{~}{D}`$ be a name for $`D`$ in $`\text{P}^{i_\alpha }\times \text{S}^{i_\alpha }`$ Again $`\stackrel{~}{D}`$ is in the Skolem hull of some $`i_\beta <i_\alpha `$ and finitely many indiscernibles $`\text{i}_n=(i_{\alpha +1},\mathrm{},i_{\alpha +n})`$. Define $`j:LL`$ as above. As we have proved if $`(p,q)G^{i_\beta }\times H^{i_\beta }`$ then $`j(p,q)G^{i_\alpha }\times H^{i_\alpha }`$. Hence the embedding $`j`$ has a canonical extension to an embedding $`\widehat{j}:L[G^{i_\beta }\times H^{i_\beta }]L[G^{i_\alpha }\times H^{i_\alpha }]`$ defined by
(2.11)
$$\widehat{j}(\tau (G^{i_\beta }\times H^{i_\beta }))=j(\tau )(G^{i_\alpha }\times H^{i_\alpha }).$$
Since $`\stackrel{~}{D}`$ is in $`\text{rng}j`$ we have $`D\text{rng}\widehat{j}`$. The proof ends as follows:
Let
(2.12)
$$p^{}G(i_\beta )\widehat{j}^1(D).$$
$`p^{}`$ exists since by induction hypothesis. $`G(i_\beta )`$ is $`Q_{i_\beta }`$ generic, and $`\widehat{j}^1(D)`$ is dense in $`Q_{i_\beta }`$ by elementarity, and hence $`\widehat{j}(p^{})D`$. Since $`p^{}L_{i_\beta }[G^{i_\beta }\times H^{i_\beta }]`$ we have $`j(p^{})=p^{}`$. So
(2.13)
$$p^{}G(i_\beta )DG(i_\alpha )D.$$
Let $`\text{G}\times \text{H}`$ be $`\text{P}\times \text{S}`$ generic over $`L`$. Suppose that $`\text{H}=h_\alpha |\alpha <\kappa `$ is the S generic filter. Let $`<,>`$ be a definable pairing function in L, such that for every $`\beta ,\gamma `$, $`<\beta ,\gamma >`$ is an L inaccessible. Since the pairing is definable and $`\kappa `$ is an indiscernible it is closed under the pairing function.
Let $`T`$ be the tree part of $`G(\kappa )`$. Our final model will be $`N=L[T,h_\alpha |\alpha C_T]`$ where
$$C_T=\{\alpha |\beta ,\gamma (\alpha =<\beta ,\gamma >\beta <_T\gamma ).$$
To finish the proof of the theorem we have to prove:
###### Proposition 2.4.
(2.14)
$$N\mathrm{`}\mathrm{`}V_\kappa \text{ has no elementary end extension”.}$$
###### Proof.
The proof will be done by a sequence of claims.
###### Claim 2.5.
$`NT\text{ is Suslin}`$.
###### Proof.
The claims follows from the fact that the forcing S is $`\kappa `$-Knaster in $`L[T]`$. Hence $`\text{S}\times T`$ is $`\kappa `$-c.c. in $`L[T]`$, so especially $`T`$ is $`\kappa `$-c.c. in $`N^{}=L[T,h_\alpha |\alpha <\kappa ]`$. But $`NN^{}`$ and $`\kappa ^N=\kappa ^N^{}`$, thus $`N`$ contains no large anti-chains of $`T`$ as well. ∎
###### Claim 2.6.
For every inaccessible $`\alpha `$
(2.15)
$$N\alpha ^{+++L}<\alpha ^{+++}\alpha C_T.$$
###### Proof.
Since for every $`\alpha C_T`$ the claim obviously holds, it will be enough to prove that other cardinals are not collapsed inside $`L[\text{G},h_\alpha |\alpha C_T]`$. For each $`\mu C_T`$ we can even work inside $`L[\text{G},h_\alpha |\alpha \mu ]`$. However since both forcing notions P and
$$\text{S}^\mu =\{S_\alpha :\alpha \mu \text{ and }\alpha \text{ is inaccesible }\}$$
factors nicely, it is obvious that the only $`L`$-cardinals collapsed are the triple successors of cardinals in $`C_T`$. ∎
Notice that by the inaccessibility of $`\kappa `$ all the collapsing functions are inside $`V_\kappa ^N`$.
Now we finish the proof of proposition 2.4. In $`(V_\kappa ^N,)`$ the tree $`T`$ is definable by the first order formula:
$$\beta <_T\gamma \alpha (\alpha \text{ is inaccessible }\alpha =<\beta ,\gamma >\alpha ^{+++L}<\alpha ^{+++}).$$
$`(V_\kappa ^N,)T\text{ is a }\kappa \text{ tree}`$,i.e., for every ordinal $`\alpha `$ $`\{xT|\text{hight}_T(x)=\alpha \}`$ is a set, and for every ordinal $`\alpha `$ there is an element of $`T`$ of hight $`\alpha `$. Assume that $`(M,E)`$ is an end extension of $`(V_\kappa ^N,)`$. Let $`a`$ be a new ordinal in $`M`$. In $`M`$ there is a tree $`T^{}`$ which end extends the tree $`T`$, since $`T`$ was definable. By elementarity
$$M\text{there is a branch }b\text{ in }T^{}\text{ of length }a.$$
Now it follows that
$$N\{xb|rk(x)<\kappa \}\text{ is a branch through }T.$$
Hence any end extension of $`(V_\kappa ^N,)`$ will provide a branch through $`T`$ in $`N`$. This is a contradiction since $`NT\text{ is Suslin}`$. ∎
## 3. Mahloness in inner models
In view of the previous result it is natural to ask whether we can get an inner model $`ML[0^{\mathrm{}}]`$ such that for every inaccessible cardinal $`\alpha M`$, $`(V_\alpha ,)`$ has no well founded elementary end extension. This turns out to be impossible by the following lemma:
###### Lemma 3.1.
Let $`\kappa >\mathrm{}_1`$ be a cardinal in $`L[0^{\mathrm{}}]`$, $`\text{cf}(\kappa )>\mathrm{}_0`$, then $`\kappa `$ is weakly Mahlo in any strictly inner model $`ML[0^{\mathrm{}}]`$. Moreover if $`\kappa `$ is a limit cardinal then $`\kappa `$ is strongly Mahlo in every $`ML[0^{\mathrm{}}]`$.
###### Proof.
The basic idea is to use the covering theorem to prove that certain cardinals are not collapsed, in any strict inner model of $`L[0^{\mathrm{}}]`$. Then we use the covering theorem again to prove that actually there must be a stationary set of inaccessibles below $`\kappa `$. Let $`ML[0^{\mathrm{}}]`$ be an inner model. Let $`I=\left\{i_\alpha \right|\alpha On\}`$ be an increasing enumeration of Silver’s indiscernibles. Then for every $`\alpha `$ such that $`\omega <\text{cf}(\alpha )`$ we have $`Mi_\alpha ^{+L}`$ is a cardinal. The proof of this uses an idea of Beller . Assume $`Mi_\alpha ^{+L}`$ is not a cardinal. Then $`|i_\alpha ^{+L}|^M=|i_\alpha |^M`$. By the covering theorem also $`M\text{cf}(i_\alpha ^{+L})=\text{cf}(i_\alpha )=|i_\alpha |`$. So in $`M`$ there is an $`f`$ $`f:i_\alpha i_\alpha ^{+L}`$ which maps in an order preserving way a cofinal subset of $`i_\alpha `$ into a cofinal subset of $`i_\alpha ^{+L}`$. Since $`L[0^{\mathrm{}}]\text{cf}(i_\alpha ^{+L})=\omega `$ choose a cofinal sequence (in $`i_\alpha ^{+L}`$) $`\{\beta _n:n<\omega \}`$ inside $`L[0^{\mathrm{}}]`$. Now let $`\gamma _n`$ be the least $`\gamma `$ such that $`f(\gamma )>\beta _n`$. We obtain that $`\{\gamma _n:n<\omega \}`$ is cofinal in $`i_\alpha `$ so $`L[0^{\mathrm{}}]\text{cf}(i_\alpha )=\omega `$. This contradicts the fact that $`i_\alpha `$ has uncountable cofinality. Hence every limit of indiscernibles of uncountable cofinality is a limit cardinal. By the covering theorem it must be a regular cardinal, so it is weakly inaccessible. Especially any uncountable cardinal is weakly inaccessible.
Suppose now that $`i_\alpha `$ is not Mahlo in $`M`$ and $`i_\alpha `$ is a limit of indiscernibles of uncountable cofinality. Then there is a club $`Ci_\alpha `$ consisting of singular cardinals in $`M`$. By the covering theorem (between $`L`$ and $`M`$) each element of $`C`$ is singular in $`L`$. Hence $`CI=\mathrm{}`$. Hence $`L[0^{\mathrm{}}]\text{cf}(i_\alpha )=\omega `$ (since it has two disjoint clubs through $`i_\alpha `$). Therefore if $`\text{cf}(i_\alpha )>\omega `$ and $`i_\alpha `$ is a limit of indiscernibles of uncountable cofinality it must be Mahlo in any strict inner model.
If $`\kappa `$ is also a limit cardinal in $`L[0^{\mathrm{}}]`$ it is strong limit by GCH. Hence it is strong limit in any inner model, so it is strongly Mahlo in $`M`$. ∎
Therefore if $`\kappa `$ is limit in $`L[0^{\mathrm{}}]`$ and $`\text{cf}(\kappa )>\omega `$, then in every inner model there is an inaccessible $`\alpha <\kappa `$ such that $`𝔼_{(V_\alpha ,\epsilon )}^{\text{wf}}\mathrm{}`$.
A natural question is whether one can have no weakly compacts in a strictly inner model of $`L[0^{\mathrm{}}]`$. We comment that if there is a $`\kappa `$ such that $`L[0^{\mathrm{}}]\kappa (\omega )^{<\omega }`$ then by a result of Silver , any inner model $`M`$, $`M\kappa (\omega )^{<\omega }`$, hence there are many ineffable cardinals in $`M`$. Similarly if there is a subtle cardinal $`\kappa `$, in $`L[0^{\mathrm{}}]`$, then obviously $`\kappa `$ is subtle in every inner model (the definition is $`\mathrm{\Pi }_1`$). Hence there are many large cardinals below it in any inner model (e.g., totally indescribables).
However the following question remains open:
Question: $`(ZFC+V=L[0^{\mathrm{}}])`$. Let $`M`$ be an inner model. Is it consistent that $`M`$ has no weakly compact cardinals ? Is it consistent that for no $`\kappa `$ $`M\kappa (\omega )^{<\omega }`$ ?
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# Wess-Zumino-Witten model on elliptic curves at the critical level
## 0. Introduction
The goal of this article is to construct a lattice model which is a variant of the Gaudin model with the help of the Wess-Zumino-Witten (WZW) model on elliptic curves at the critical level and to find its eigenvectors by means of the bosonisation of the WZW model.
As is well known, correlation functions of the WZW model on the Riemann sphere satisfy the Knizhnik-Zamolodchikov equations when the level $`k`$ of the model is not critical, i.e., $`kh^{}`$, where $`h^{}`$ is the dual Coxeter number of the simple Lie algebra $`𝔤`$ which describes symmetry of the model. In constrast to this case, when the level is critical, $`k=h^{}`$, there can be no longer such an equation since the Sugawara construction of the energy-momentum tensor breaks down.
In \[FFR\] Feigin, Frenkel and Reshetikhin found an interpretation of such case as a lattice model called the (rational) Gaudin model \[G1\], \[G2\], \[G3\], a quasi-classical limit of the totally isotropic spin chain model (the XXX model). (See \[S\].) It is also shown in \[FFR\] that the free field realisation of the WZW model provides a new diagonalisation method which recovers the results of the Bethe Ansatz method.
In this paper we apply this story to the WZW model on elliptic curves. The state space of the lattice model thus obtained is a space of functions over the Cartan subalgebra. The transfer matrix (the generating function of the Hamiltonians) is a quasi-classical limit of the IRF type lattice model. Therefore we name the model a “face type Gaudin model”. The bosonisation technique is also applied to find a Bethe Ansatz type eigenvectors.
Felder and Varchenko \[FV1\], \[FV2\] studied this system and its Bethe Ansatz, which arose from the stationary phase method of the integral representation of solutions of the KZB equations. Their results are used in the analysis of a spin chain model with elliptic exchanges in \[I\]. The same kind of system (the Gaudin-Calogero model) has been also studied by Enriquez, Feigin and Rubtsov \[ER\], \[EFR\] who started from the quantisation of the Hitchin system on elliptic curves.
This paper is organised as follows. In Section 1 we state our main results. Section 2 explains how to derive our transfer matrix from the WZW model on elliptic curves at the critical level. In order to construct the Bethe eigenvectors we make use of the Wakimoto modules with the critical level, which we recall in Section 3. The last section, Section 4, shows that the free field theory gives the eigenvector of the transfer matrix in the form of the Bethe vector.
Details shall be published in a forthcoming paper. Mathematically delicate conditions like finiteness of modules are not specified unless they are essential.
## 1. Main results
First we fix notations. Throughout this paper $`\tau `$ is a complex number with positive imaginary part. The elliptic curve with modulus $`\tau `$ is denoted by $`X=X_\tau =/+\tau `$.
Let $`𝔤`$ be a finite dimensional simple Lie algebra of rank $`l`$, $`𝔥`$ be its Cartan subalgebra and
(1.1)
$$𝔤=𝔥\underset{\alpha \mathrm{\Delta }}{}𝔤_\alpha $$
be the root space decomposition, where $`\mathrm{\Delta }`$ is the set of roots. We use the Cartan-Killing form normalised as follows:
(1.2)
$$(A|B):=\frac{1}{2h^{}}\mathrm{Tr}_𝔤(\mathrm{ad}A\mathrm{ad}B)\text{ for }A,B𝔤,$$
We identify $`𝔥`$ and its dual space $`𝔥^{}`$ through this inner product. We fix the simple roots $`\{\alpha _1,\mathrm{},\alpha _l\}`$, Chevalley generators $`\{H_i,E_i,F_i\}_{i=1,\mathrm{},l}`$ and a basis $`e_\alpha `$ of $`𝔤_\alpha `$, such that $`e_{\alpha _i}=E_i`$ for $`i=1,\mathrm{},l`$ and $`(e_\alpha |e_\alpha ^{})=\delta _{\alpha ,\alpha ^{}}`$. The set of positive and negative roots are denoted by $`\mathrm{\Delta }_+=\{\beta _1,\mathrm{},\beta _s\}`$ and $`\mathrm{\Delta }_{}`$ respectively. We fix an orthonormal basis $`\{h_r\}_{r=1,\mathrm{},l}`$ of $`𝔥`$ and the coordinate system of $`𝔥`$, $`(\xi _1,\mathrm{},\xi _l)^l`$, associated to it.
Let $`V_i`$ ($`i=1,\mathrm{},N`$) be finite dimensional irreducible representations of $`𝔤`$ with the highest weight $`\lambda _i`$ and $`V(0)`$ be the $`0`$-weight space of $`V=V_1\mathrm{}V_N`$. We denote the action of $`𝔤`$ on the $`i`$-th factor $`V_i`$ by $`\rho _i`$ as usual. The dual (right) action of $`𝔤`$ on $`\mathrm{\Phi }V^{}=\mathrm{Hom}_{}(V,)`$ is denoted as
(1.3)
$$\rho _i^{}(A)\mathrm{\Phi }(v):=\mathrm{\Phi }(\rho _i(A)v),$$
and the $`0`$-weight space of $`V^{}`$ by $`V^{}(0)`$.
We define a differential operator $`\widehat{\tau }(u)`$ on $`V^{}(0)`$-valued functions on
(1.4)
$$S=\{H𝔥\alpha (H)\text{ for all }\alpha \mathrm{\Delta }\}$$
with fixed complex parameters $`z_i`$ ($`i=1,\mathrm{},N`$) and a spectral parameter $`u`$ as follows:
(1.5)
$$\begin{array}{cc}& \widehat{\tau }(u):=\frac{1}{2}\underset{r=1}{\overset{l}{}}_{\xi _r,u}^2+\hfill \\ & +\frac{1}{2}\underset{i,j=1}{\overset{N}{}}\underset{\alpha \mathrm{\Delta }}{}w_{\alpha (H)}(z_iu)w_{\alpha (H)}(z_ju)\rho _j^{}(e_\alpha )\rho _i^{}(e_\alpha ),\hfill \end{array}$$
where
(1.6)
$$_{\xi _r,u}:=\frac{}{\xi _r}\underset{i=1}{\overset{N}{}}\zeta _{11}(z_iu)\rho _i^{}(h_r).$$
The quasi-periodic functions $`w_c(z)`$ and $`\zeta _{11}(z)`$ are defined by (A.2).
Since this operator can be interpreted as a trace of square of the dynamical (or modified) classical $`r`$-matrix which is a classical limit of the IRF type lattice models (cf. \[F1\], \[F2\], \[FS\]), we call $`\widehat{\tau }(u)`$ the transfer matrix of the face type Gaudin model.
###### Theorem 1.1.
The operator $`\widehat{\tau }(u)`$ commutes with itself:
(1.7)
$$[\widehat{\tau }(u),\widehat{\tau }(u^{})]=0.$$
This can be checked by direct computation, but we shall show in Section 2 that it can be proved with the help of the WZW model at the critical level.
Taking Theorem 1.1 into account, we can pose a question of simultaneous diagonalisation of $`\widehat{\tau }(u)`$. Our main result is the following Bethe Ansatz solution of this problem. Assume that $`V_i`$ is the dual Verma module $`M_{\lambda _i}^{}`$ of $`𝔤`$ and that the sum of the weights $`\lambda _i`$ belongs to the positive root lattice:
(1.8)
$$\underset{i=1}{\overset{N}{}}\lambda _i=\underset{j=1}{\overset{M}{}}\alpha _{i(j)},$$
for a sequence $`\{\alpha _{i(j)}\}_{j=1,\mathrm{},M}`$ of simple roots. The symbol $`ȷ(v)`$ for $`vM_\lambda ^{}`$ is the canonical paring of $`v`$ with the highest weight vector of the Verma module $`M_\lambda `$.
###### Theorem 1.2.
If there are complex numbers $`t_j`$ ($`j=1,\mathrm{},M`$) satisfying the Bethe Ansatz equation,
(1.9)
$$\underset{i=1}{\overset{N}{}}(\alpha _{i(j)}\lambda _i)\zeta _{11}(t_jz_i)=\underset{j^{}j}{}(\alpha _{i(j)}\alpha _{i(j^{})})\zeta _{11}(t_jt_j^{}),$$
for any $`j=1,\mathrm{},M`$, then
(1.10)
$$\mathrm{\Psi }(H;v):=\underset{\{I_j\}}{}\underset{a=1}{\overset{N}{}}I_a;v_a;z_a,t_j(jI_a)$$
is an eigenvector of $`\widehat{\tau }(u)`$. Here $`\{I_j\}`$ is a partition of the set $`\{1,\mathrm{},M\}`$ into $`N`$ sets, $`I_1\mathrm{}I_N=\{1,\mathrm{},M\}`$, and the symbol $`I;v;t_j(jI)`$ is defined as follows:
(1.11)
$$\begin{array}{c}I;v;z,t_j(jI):=\underset{\sigma 𝔖}{}w_{\alpha _{i(\sigma (1))}}(t_{\sigma (1)}t_{\sigma (2)})w_{\alpha _{i(\sigma (1))}+\alpha _{i(\sigma (2))}}(t_{\sigma (2)}t_{\sigma (3)})\times \mathrm{}\hfill \\ \hfill \mathrm{}\times w_{\alpha _{i(\sigma (1))}+\alpha _{i(\sigma (2))}+\mathrm{}+\alpha _{i(\sigma (m))}}(t_{\sigma (m)}z)ȷ(E_{i(\sigma (m))}\mathrm{}E_{i(\sigma (1))}v)\end{array}$$
if $`I=\{1,\mathrm{},m\}`$. (We write $`w_\alpha `$ instead of $`w_{\alpha (H)}`$ for short.) The eigenvalue of $`\mathrm{\Psi }(H;v)`$ is
(1.12)
$$\tau _\mathrm{\Psi }(u):=\frac{1}{2}\underset{r=1}{\overset{l}{}}𝜻(h_r;z,t;u)^2+\frac{}{u}𝜻(\rho ;z,t;u).$$
Here $`𝛇`$ is defined by
(1.13)
$$𝜻(h;z,t;u):=\underset{i=1}{\overset{N}{}}\lambda _i(h)\zeta _{11}(z_iu)\underset{j=1}{\overset{M}{}}\alpha _{i(j)}(h)\zeta _{11}(t_ju),$$
and $`\rho `$ is the half sum of the positive roots of $`𝔤`$.
We shall show how to prove this theorem by bosonisation in Section 4.
## 2. WZW model on elliptic curves at the critical level
The idea behind the definition (1.5) and Theorem 1.1 is that the linear functional $`\mathrm{\Phi }(H;v)`$ and $`(\widehat{\tau }(u)\mathrm{\Phi })(H;v)`$ are analogues of an $`N`$-point function of the WZW model and the correlation funcion of the energy-momentum tensor, respectively.
First we define the geometric data on which the WZW model lives. Let $`𝔛=S\times X_\tau `$ be the trivial family of elliptic curves and $`\pi _{𝔛/S}:𝔛S`$ be the projection. The divisor of $`𝔛`$ corresponding to the parameter $`z_i`$ is denoted by $`P_i`$ and their sum by $`D`$:
(2.1)
$$P_i:=S\times z_imod+\tau 𝔛,D:=P_1+\mathrm{}+P_N$$
The following general definitions of $`N`$-point functions is valid not only for the WZW model but also for the free field theories which we use later.
Assume that the following Lie algebraic data on $`𝔛`$ are given:
* $`𝔞_𝔛`$: a Lie algebra bundle with a fibre isomorphic to a Lie algebra $`𝔞`$.
* $`:𝔞_𝔛\times _{𝒪_S}𝔞_𝔛\mathrm{\Omega }_{𝔛/S}^1`$: an $`𝒪_S`$-bilinear $`\mathrm{\Omega }_{𝔛/S}^1`$-valued pairing which is a 2-cocycle up to exact forms:
(2.2)
$$\begin{array}{c}𝔞_𝔛\times 𝔞_𝔛(A,B)AB\mathrm{\Omega }_{𝔛/S}^1,\\ AB+BAd_{𝔛/S}𝒪_𝔛,\\ [A,B]C+[B,C]A+[C,A]Bd_{𝔛/S}𝒪_𝔛,\end{array}$$
for any $`A,B,C𝔞_𝔛`$. For example, if there is a connection $``$ on $`𝔞_𝔛`$ along the fibre compatible with the Lie algebra structure,
(2.3)
$$\begin{array}{c}:𝔞_𝔛𝔞_𝔛_{𝒪_𝔛}\mathrm{\Omega }_{𝔛/S}^1,\\ [A,B]=[A,B]+[A,B]\text{ for }A,B𝔞_𝔛,\end{array}$$
and an invariant $`𝒪_𝔛`$-inner product $`()`$ of $`𝔞_𝔛`$ which satisfies
(2.4)
$$d_{𝔛/S}(A|B)=(AB)+(AB)\mathrm{\Omega }_{𝔛/S}^1\text{for }A,B𝔞_𝔛,$$
then $`AB:=(AB)`$ has desired properties.
###### Example 2.1.
In this section we use the following data. The Lie algebra bundle $`𝔤_𝔛`$ over $`𝔛`$ is defined as the quotient of $`S\times \times 𝔤`$ by the $`^2`$-action,
(2.5)
$$(m,n)(H;t;A)=(H;t+m\tau +n;e^{2\pi im\mathrm{ad}H}A),$$
for $`(m,n)^2`$ and $`(H;t;A)S\times \times 𝔤`$. The connection $``$ of (2.3) is the trivial differentiation $`d/dtdt`$ along the elliptic curve and the invariant bilinear form $`()`$ is defined by (1.2).
For $`i=1,\mathrm{},N`$, we put
(2.6)
$$\begin{array}{cc}\hfill 𝔞_S^{P_i}& :=(\pi _{𝔛/S})_{}(𝔞_𝔛(P_i))_{P_i}^{}(𝔞𝒪_S)((x_i)),\hfill \\ \hfill 𝔞_{S,+}^{P_i}& :=(\pi _{𝔛/S})_{}(𝔞_𝔛)_{P_i}^{}(𝔞𝒪_S)[[x_i]],\hfill \\ \hfill 𝔞_S^D& :=\underset{i=1}{\overset{N}{}}𝔞_S^{P_i},\widehat{𝔞}_S^D:=𝔞_S^D𝒪_S\widehat{k},\hfill \end{array}$$
where $`x_i`$ is a formal local parameter and the central extension $`\widehat{𝔞}_S^D`$ is defined by the $`𝒪_S`$-valued 2-cocycle
(2.7)
$$c(A,B):=\underset{i=1}{\overset{N}{}}Res_{P_i}A_iB_i,$$
where $`A=(A_i)_{i=1}^N,B=(B_i)_{i=1}^N𝔞_S^D`$
###### Example 2.2.
For Example 2.1, the central extension $`\widehat{𝔤}_S^{P_i}`$ is naturally isomorphic to the space of $`\widehat{𝔤}`$-valued functions on $`S`$, $`\widehat{𝔤}_S=\widehat{𝔤}_{}𝒪_S`$, where $`\widehat{𝔤}`$ is the affine Lie algebra associated to $`𝔤`$.
Let $`𝔞_{\dot{𝔛}}^D`$ be the space of meromorphic sections of $`𝔞_𝔛`$ which are globally defined along the fibre of $`\pi _{𝔛/S}`$ and holomorphic except at $`D`$:
(2.8)
$$𝔞_{\dot{𝔛}}^D:=(\pi _{𝔛/S})_{}(𝔞_𝔛(D)),$$
which is naturally regarded as a Lie subalgebra of $`\widehat{𝔞}_S^D`$.
###### Example 2.3.
For the Lie algebra bundle $`𝔤_𝔛`$ in Example 2.1, we denote $`𝔞_{\dot{𝔛}}^D`$ by $`𝔤_{\dot{𝔛}}^D`$. The section $`w_{\alpha (H)}(tz_i)e_\alpha `$ (cf. (A.2)) belongs to $`𝔤_{\dot{𝔛}}^D`$.
Let us take $`\widehat{𝔞}_S^{P_i}`$-modules $`_i`$ ($`i=1,\mathrm{},N`$) with the same level $`\widehat{k}=k`$ and define $`=_1\mathrm{}_N`$.
###### Definition 2.4.
(i) The sheaf of conformal blocks $`𝒞(𝔞_𝔛,D,)`$ is defined to be the space of $`𝒪_S`$-linear functionals on $``$ which vanishes on $`𝔞_{\dot{𝔛}}^D`$: $`\mathrm{\Phi }\mathrm{𝑜𝑚}_{𝒪_S}(,𝒪_S)`$ belongs to $`𝒞(𝔞_𝔛,D,)`$ if and only if it satisfies
(2.9)
$$\mathrm{\Phi }(A_{\dot{𝔛}}v)=0\text{for all }A_{\dot{𝔛}}𝔞_{\dot{𝔛}}^D\text{ and }v\text{.}$$
This equation (2.9) is called the Ward identity.
(ii) There is a flat connection on $`𝒞(𝔞_𝔛,D,)`$. A flat section is called the $`N`$-point function. The set of $`N`$-point functions is denoted by $`𝒞^{\mathrm{hor}}(𝔞_𝔛,D,)`$.
The flat connection on $`𝒞(𝔤_𝔛,D,)`$ is defined as follows. We use the notation,
(2.10)
$$\rho _i^{}(h_r\{\theta _i(x_i)\})=\underset{m}{}\theta _{i,m}\rho _i^{}(h_r[m]),$$
for an element $`\theta _i(x_i)=_m\theta _{i,m}x_i^m`$, of $`𝒪_S((x_i))`$. The connection on $`𝒞(𝔤_𝔛,D,)`$ in the direction of $`\xi _r`$ is defined by
(2.11)
$$_{/\xi _r}^{}=\frac{}{\xi _r}\rho ^{}(h_r\{Z(t)\}):=\frac{}{\xi _r}\underset{i=1}{\overset{N}{}}\rho _i^{}(h_r\{Z(x_i+z_i)\}),$$
where $`Z(t)`$ is a meromorphic function with poles in $`D`$ and has quasi-periodicity
(2.12)
$$Z(t+m\tau +n)=Z(t)2\pi im.$$
(In (2.11) $`Z(t)`$ is expanded around $`z_i`$ in the power series of $`x_i=tz_i`$ and substituted into (2.10).) For example, we can take $`Z(t)=\zeta _{11}(tz_1)`$ (cf. (A.2)).
We denote $`𝒞(𝔤_𝔛,D,)`$ and $`𝒞^{\mathrm{hor}}(𝔤_𝔛,D,)`$ by $`𝒞(D,)`$ and $`𝒞^{\mathrm{hor}}(D,)`$ for short. We mean by the WZW model the theory on $`𝒞(D,)`$ or $`𝒞^{\mathrm{hor}}(D,)`$. Hereafter we assume that the $`\widehat{𝔤}`$-module $`M_i`$ is generated over $`\widehat{𝔤}`$ by a $`𝔤`$-submodule $`V_i`$ on which the centre $`\widehat{k}\widehat{𝔤}`$ acts as multiplication by $`k`$ and $`Ax^m`$ ($`A𝔤`$, $`m>0`$) acts by $`0`$. Put $`_i:=M_i𝒪_S`$, $`V=V_1\mathrm{}V_N`$.
###### Lemma 2.5.
An $`N`$-point function $`\mathrm{\Phi }(H;v)`$ of the WZW model is determined by its values on $`vV(0)`$, where $`V(0)`$ is the $`0`$-weight space of $`V`$. In other words, the following restriction map is injective:
(2.13)
$$𝒞^{\mathrm{hor}}(D,)\mathrm{\Phi }\mathrm{\Phi }(H;v)V^{}(0)𝒪_S.$$
This is a consequence of the Ward identity (2.9) and the flatness condition.
Let us return to the discussion for the general Lie algebra bundle $`𝔞_𝔛`$. Let $`Q`$ be a point not contained in $`D`$ and $`\mathrm{V}ac_{Q,k}`$ be the $`\widehat{𝔞}_S^Q`$-module induced from the trivial $`\widehat{𝔞}_{S,+}^Q`$-module $`𝒪_{S,k}`$ of level $`k`$:
(2.14)
$$\mathrm{V}ac_{Q,k}:=\mathrm{Ind}_{\widehat{𝔞}_{S,+}^Q}^{\widehat{𝔞}_S^Q}𝒪_{S,k}$$
where $`𝒪_{S,k}=𝒪_S|0`$ as a linear space, $`𝔞_{S,+}^Q𝒪_{S,k}=0`$ and $`\widehat{k}`$ acts as a multiplication by $`k`$. We call $`\mathrm{V}ac_{Q,k}`$ the vacuum module of level $`k`$ at $`Q`$.
###### Lemma 2.6.
The canonical linear map $`vv|0\mathrm{V}ac_{Q,k}`$ induces isomorphisms:
(2.15)
$$\begin{array}{cc}\hfill 𝒞(𝔞_𝔛,D,)& \stackrel{}{}𝒞(𝔞_𝔛,D+Q,\mathrm{V}ac_{Q,k}),\hfill \\ \hfill 𝒞^{\mathrm{hor}}(𝔞_𝔛,D,)& \stackrel{}{}𝒞^{\mathrm{hor}}(𝔞_𝔛,D+Q,\mathrm{V}ac_{Q,k}).\hfill \end{array}$$
This property is called propagation of vacua in \[TUY\].
The following proposition reveals the real nature of the operator $`\widehat{\tau }(u)`$ in (1.5). The operator obtained directly from the WZW model differs from $`\widehat{\tau }(u)`$ by conjugation. Let us define an operator $`\stackrel{~}{\tau }(u)`$ by
(2.16)
$$\begin{array}{cc}\hfill \stackrel{~}{\tau }(u)& =\mathrm{\Pi }(H,\tau )\widehat{\tau }(u)\mathrm{\Pi }(H,\tau )^1\hfill \\ & =\frac{1}{2}\underset{r=1}{\overset{l}{}}_{\xi _r,u}^2+\underset{r=1}{\overset{l}{}}\frac{}{\xi _r}\mathrm{log}\mathrm{\Pi }(H,\tau )_{\xi _r,u}\hfill \\ & +\frac{1}{2}\underset{i,j=1}{\overset{N}{}}\underset{\alpha \mathrm{\Delta }}{}w_{\alpha (H)}(z_iu)w_{\alpha (H)}(z_ju)\rho _j^{}(e_\alpha )\rho _i^{}(e_\alpha )+2\pi ih^{}\frac{}{\tau }\mathrm{log}\mathrm{\Pi }(H,\tau ),\hfill \end{array}$$
where $`\mathrm{\Pi }(H,\tau )`$ is the normalised Weyl-Kac denominator,
(2.17)
$$\mathrm{\Pi }(H,\tau )=q^{dim𝔤/24}(q;q)_{\mathrm{}}^l\underset{\alpha \mathrm{\Delta }_+}{}(e^{\pi i\alpha (H)}e^{i\pi \alpha (H)})\underset{\alpha \mathrm{\Delta }}{}(qe^{2\pi i\alpha (H)};q)_{\mathrm{}}.$$
Here we use the usual notations, $`q=e^{2\pi i\tau }`$ and $`(x;q)_{\mathrm{}}=_{n=0}^{\mathrm{}}(1xq^n)`$.
###### Proposition 2.7.
According to Lemma 2.6, there is a $`(N+1)`$-point function $`\stackrel{~}{\mathrm{\Phi }}`$ corresponding to $`\mathrm{\Phi }𝒞^{\mathrm{hor}}(D,)`$. We have
(2.18)
$$(\stackrel{~}{\tau }(u)\mathrm{\Phi })(H;v)=\stackrel{~}{\mathrm{\Phi }}(H;vS[2]|0),$$
for $`HS`$, $`vV(0)`$, where $`u`$ is the coordinate of $`Q`$ on the complex plane and $`S[2]`$ is defined as a coefficient of the Sugawara tensor,
(2.19)
$$S(u):=\frac{1}{2}\underset{p=1}{\overset{dim𝔤}{}}\genfrac{}{}{0.0pt}{}{}{}J_p(u)J_p(u)\genfrac{}{}{0.0pt}{}{}{}=\underset{n}{}S[n]z^{n1}.$$
Here $`\{J_p\}_{p=1,\mathrm{},dim𝔤}`$ is an orthonormal basis of $`𝔤`$ and the symbol $`\genfrac{}{}{0.0pt}{}{}{}\genfrac{}{}{0.0pt}{}{}{}`$ is the normal ordering operation.
This means that $`\stackrel{~}{\tau }(u)\mathrm{\Phi }(v_1\mathrm{}v_N)`$ is the correlation function $`S(u)v_1(z_1)\mathrm{}v_N(z_n)`$ of $`S(u)`$ in the context of the conformal field theory.
Hitherto the level is arbitrary. To prove Theorem 1.1, we need to fix the level to the critical value, $`k=h^{}`$, where $`S[2]|0`$ is a singular vector of imaginary weight. Roughly speaking, by virtue of this fact the correlation function of two Sugawara tensors, $`S(u)S(u^{})v_1(z_1)\mathrm{}v_N(z_n)`$, is irrelevant to the order of insertion of $`S(u)`$ and $`S(u^{})`$, from which the commutativity (1.7) follows.
## 3. Wakimoto modules at the critical level
The Bethe vector (1.10) is constructed by means of the Wakimoto realisation of affine Lie algebras from the free field theory. In this section we review basic facts about the Wakimoto representations, following \[K\]. See also \[W\], \[FF1\], \[FF2\], \[FFR\].
The bosonic ghost fields,
(3.1)
$$\beta _\alpha (z)=\underset{m}{}z^{m1}\beta _\alpha [m],\gamma ^\alpha (z)=\underset{m}{}z^m\gamma ^\alpha [m],(\alpha \mathrm{\Delta }_+)$$
satisfy the following operator product expansions:
(3.2)
$$\beta _\alpha (z)\gamma ^\alpha ^{}(w)\frac{\delta _\alpha ^\alpha ^{}}{zw}.$$
We denote the Heisenberg algebra generated by $`\beta _\alpha [m]`$ and $`\gamma ^\alpha [m]`$ ($`\alpha \mathrm{\Delta }_+`$, $`m`$) by $`\widehat{\mathrm{Gh}}(𝔤)`$.
The ghost Fock space $`^{\mathrm{gh}}`$ is defined as a left $`\widehat{\mathrm{Gh}}(𝔤)`$-module generated by the vacuum vector $`|0^{\mathrm{gh}}`$, satisfying
(3.3)
$$\beta _\alpha [m]|0^{\mathrm{gh}}=0,\gamma ^\alpha [n]|0^{\mathrm{gh}}=0$$
for any $`\alpha \mathrm{\Delta }_+`$, $`m0`$, $`n>0`$. The normal ordered product $`:P:`$ of a monomial $`P`$ of $`\beta _\alpha [m]`$’s and $`\gamma ^\alpha [m]`$’s is defined by putting annihilation operators of $`|0^{\mathrm{gh}}`$ appearing in $`P`$ to the right side in the product.
The free boson fields,
(3.4)
$$\begin{array}{c}\varphi _i(z):=\varphi _i[0]\mathrm{log}z+\underset{m\{0\}}{}\frac{z^m}{m}\varphi _i[m],\varphi _i(z):=\underset{m}{}z^{m1}\varphi _i[m],\\ \varphi (H;z)=\underset{i=1}{\overset{l}{}}a_i\varphi _i(z),\varphi [H;m]=\underset{i=1}{\overset{l}{}}a_i\varphi _i[m],\end{array}$$
for $`H=_{i=1}^la_iH_i𝔥`$ have trivial operator product expansions:
(3.5)
$$\varphi (H;z)\varphi (H^{};w)0,\varphi (H;z)\varphi (H^{};w)0,$$
for any $`H,H^{}𝔥`$.
The commutative algebra generated by $`\varphi _i[m]`$ ($`i=1,\mathrm{},l`$, $`m`$) is denoted by $`\widehat{\mathrm{Bos}}(𝔤)`$.
For any one-form $`\lambda (x)dx𝔥^{}((x))dx`$, we define a one-dimensional representation $`\sigma _{\lambda (x)dx}`$ of $`\widehat{\mathrm{Bos}}(𝔤)`$ by
(3.6)
$$\begin{array}{c}\sigma _{\lambda (x)dx}=|\lambda (x)dx,\\ f(x)|\lambda (x)dx=Res_{x=0}(\lambda (x),f(x))dx,\end{array}$$
where $`f(x)𝔥((x))`$ is identified with an element of $`\widehat{\mathrm{Bos}}`$ by the isomorphism defined by $`H_ix^mH_i[m]`$ and $`(,)`$ is the canonical pairing of $`𝔥^{}`$ and $`𝔥`$. In other words,
(3.7)
$$\varphi [H;m]|\lambda (x)dx=\lambda ^{(m1)}(H)|\lambda (x)dx,$$
where $`\lambda (x)dx=_n\lambda ^{(n)}x^ndx`$, $`\lambda ^{(n)}𝔥^{}`$. Hereafter we assume that $`\lambda ^{(n)}=0`$ for $`n2`$.
###### Proposition 3.1.
\[W\], \[FF1\], \[FF2\], \[K\]. For each Chevalley generator $`H_i`$, $`E_i`$ or $`F_i`$ of $`𝔤`$, there exists a differential polynomial of the free fields,
$$X(z)=\underset{m}{}X[m]z^{m1}:=:R(X;\gamma (z),\beta (z),\varphi (z)):,$$
which gives the corresponding Kac-Moody current. Namely, a Lie algebra homomorphism $`\omega `$ from the affine Lie algebra $`\widehat{𝔤}=𝔤[t,t^1]\widehat{k}`$ to a completion of $`\widehat{\mathrm{Gh}}(𝔤)\widehat{\mathrm{Bos}}(𝔤)`$ can be defined by
(3.8)
$$\omega (Xt^m)=X[m],\omega (\widehat{k})=h^{},$$
for all $`X𝔤`$, $`m`$, where $`\widehat{k}`$ is the centre of $`\widehat{𝔤}`$ and $`h^{}`$ is the dual Coxeter number of $`𝔤`$. Moreover, the Sugawara tensor $`S(z)`$ defined by (2.19) is expressed in terms of the free bosons as follows:
(3.9)
$$S(z):=\frac{1}{2}\underset{r=1}{\overset{l}{}}:\varphi (h_r;z)\varphi (h_r;z):\frac{1}{2}^2\varphi (2\rho ;z),$$
###### Definition 3.2.
Denote $`^{\mathrm{gh}}\sigma _{\lambda (x)dx}`$ by $`\mathrm{W}ak_{\lambda (x)dx}`$. We regard this module as a $`\widehat{𝔤}`$-module through the above homomorphism $`\omega :\widehat{𝔤}\widehat{\mathrm{Gh}}(𝔤)\widehat{\mathrm{Bos}}(𝔤)`$ and call it a Wakimoto module of level $`h^{}`$ (or of critical level) with highest weight $`\lambda (x)dx`$.
The Wakimoto module contains a $`𝔤`$-submodule isomorphic to the dual Verma module $`M_\lambda ^{}`$ of $`𝔤`$ with the highest weight $`\lambda =\lambda ^{(1)}`$. (See Proposition 4.4 of \[K\].) We denote it by $`\mathrm{W}ak_{\lambda (x)dx}^0`$. It satisfies for any $`m>0`$ and $`X𝔤`$,
(3.10)
$$X[m]\mathrm{W}ak_{\lambda (x)dx}^0=0.$$
## 4. Free field theories and Bethe vectors
We can “decompose” the WZW model in Section 2 into free field theories, using the Wakimoto realisation, Theorem 3.1. The Bethe vectors in Theorem 1.2 are nothing but the $`N`$-point functions of the free field theories.
Let us define free field theories in the framework introduced in Section 2. We defined the Lie algebra bundle $`𝔤_𝔛`$ as a quotient of $`S\times \times 𝔤`$ by the $`^2`$-action (2.5). Note that this action preserves the triangular decomposition $`𝔤=𝔫_+𝔥𝔫_{}`$. We regard vector bundles,
(4.1) $`𝜷_𝔛`$ $`:=^2\backslash S\times \times 𝔫_+`$
(4.2) $`𝜸_𝔛`$ $`:=(^2\backslash S\times \times 𝔫_{})_{𝒪_𝔛}\mathrm{\Omega }_{𝔛/S}^1,`$
(4.3) $`\mathrm{Bos}_𝔛`$ $`:=^2\backslash S\times \times 𝔥,`$
as abelian Lie algebra bundles and put $`\mathrm{Gh}_𝔛:=𝜷_𝔛𝜸_𝔛`$. We apply Definition 2.4 to $`𝔞_𝔛=\mathrm{Gh}_𝔛`$ and $`𝔞_𝔛=\mathrm{Bos}_𝔛`$. The pairing (2.2) for $`\mathrm{Bos}_𝔛`$ is trivial (i.e., $`=0`$) and the pairing for $`\mathrm{Gh}_𝔛`$ is defined by
(4.4)
$$(A_1,B_1dt)(A_2,B_2dt)^{\mathrm{gh}}:=(A_1B_2)dt(B_1A_2)dt\mathrm{\Omega }_𝔛^1,$$
where $`A_i𝜷_𝔛`$, $`B_idt𝜸_𝔛`$ for $`i=1,2`$ and $`(|)`$ denotes the inner product defined by (1.2). Here we identify $`𝜷_𝔛`$ and $`𝜸_𝔛`$ with a subbundle of $`𝔤_𝔛`$ and a subbundle of $`𝔤_𝔛\mathrm{\Omega }_{𝔛/S}^1`$ respectively.
The algebra defined by (2.6) for $`D=P`$ (a point) in this case is isomorphic to $`\widehat{\mathrm{Gh}}_S(𝔤)=\widehat{\mathrm{Gh}}(𝔤)_{}𝒪_S`$ when $`𝔞_𝔛=\mathrm{Gh}_𝔛`$ and to $`\widehat{\mathrm{Bos}}_S(𝔤)=\widehat{\mathrm{Bos}}(𝔤)_{}𝒪_S`$ when $`𝔞_𝔛=\mathrm{Bos}_𝔛`$.
We denote the sheaf of conformal blocks $`𝒞(𝔞_𝔛,D,)`$ and the space of $`N`$-point functions $`𝒞^{\mathrm{hor}}(𝔞_𝔛,D,)`$ defined in Definition 2.4 for $`𝔞_𝔛=\mathrm{Gh}_𝔛,\mathrm{Bos}_𝔛`$ by
(4.5) $`𝒞^{\mathrm{gh}}(D,):=𝒞(\mathrm{Gh}_𝔛,D,M),`$ $`𝒞^{\mathrm{bos}}(D,):=𝒞(\mathrm{Bos}_𝔛,D,M),`$
(4.6) $`𝒞^{\mathrm{gh},\mathrm{hor}}(D,):=𝒞^{\mathrm{hor}}(\mathrm{Gh}_𝔛,D,M),`$ $`𝒞^{\mathrm{bos},\mathrm{hor}}(D,):=𝒞^{\mathrm{hor}}(\mathrm{Bos}_𝔛,D,M),`$
and define
(4.7)
$$\begin{array}{cc}\hfill 𝒞^{\mathrm{free}}(D,)& :=𝒞(\mathrm{Gh}_𝔛\mathrm{Bos}_𝔛,D,),\hfill \\ \hfill 𝒞^{\mathrm{free},\mathrm{hor}}(D,)& :=𝒞^{\mathrm{hor}}(\mathrm{Gh}_𝔛\mathrm{Bos}_𝔛,D,).\hfill \end{array}$$
Assume that $`_i`$ is a $`\widehat{\mathrm{Gh}}_S^{P_i}`$-module and $`𝒩_i`$ is a $`\widehat{\mathrm{Bos}}_S^{P_i}`$-module. Then $`_i_{𝒪_S}𝒩_i`$ is a $`(\widehat{\mathrm{Gh}}_S^{P_i}\widehat{\mathrm{Bos}}_S^{P_i})`$-module and hence is a $`(\mathrm{Gh}_𝔛\mathrm{Bos}_𝔛)_S^{,P_i}`$-module. The important fact is that the $`N`$-point function of the free field theory naturally gives the $`N`$-point function of the WZW model.
###### Proposition 4.1.
$`𝔤_{\dot{𝔛}}^D((_i_{𝒪_S}𝒩_i))(\mathrm{Gh}_{\dot{𝔛}}^D\mathrm{Bos}_{\dot{𝔛}}^D)((_i_{𝒪_S}𝒩_i))`$. Hence the identity morphism from $`(_i_{𝒪_S}𝒩_i)`$ to itself induces an $`𝒪_S`$-linear map:
(4.8)
$$\iota :𝒞^{\mathrm{free}}(D,(_i_{𝒪_S}𝒩_i))𝒞(D,(_i_{𝒪_S}𝒩_i)).$$
Moreover, this induces a $``$-linear map between spaces of $`N`$-point functions:
(4.9)
$$\iota :𝒞^{\mathrm{free},\mathrm{hor}}(D,(_i_{𝒪_S}𝒩_i))𝒞^{\mathrm{hor}}(D,(_i_{𝒪_S}𝒩_i)),$$
If $`_i=^{\mathrm{gh}}𝒪_S`$ and $`𝒩_i=\sigma _{\mu _i(x)dx}𝒪_S`$ for $`\mu _i(x)dx𝔥^{}((x))dx`$, then an $`N`$-point function of the free field gives an $`N`$-point function of the WZW model with the Wakimoto modules. It is easy to see that
(4.10)
$$𝒞^{\mathrm{free},\mathrm{hor}}(D,(_i_{𝒪_S}𝒩_i))𝒞^{\mathrm{gh},\mathrm{hor}}(D,)_{}𝒞^{\mathrm{bos},\mathrm{hor}}(D,𝒩),$$
hence it is enough to find $`N`$-point functions in $`𝒞^{\mathrm{gh},\mathrm{hor}}(D,(^{\mathrm{gh}})^N)`$ and $`𝒞^{\mathrm{bos},\mathrm{hor}}(D,_{i=1}^N\sigma _{\mu _i(x)dx}𝒪_S)`$ to construct an $`N`$-point function of the WZW model. Let us denote $`𝝈_{\stackrel{}{\mu }dx}:=_{i=1}^N\sigma _{\mu _i(x)dx}𝒪_S`$ and $`|\stackrel{}{\mu }dx:=_{i=1}^N|\mu _i(x)dx`$ for simplicity.
###### Lemma 4.2.
(i) $`𝒞^{\mathrm{gh},\mathrm{hor}}(D,(^{\mathrm{gh}})^N)=\mathrm{\Phi }^{\mathrm{gh}}(H;v),`$ where
$$\mathrm{\Phi }^{\mathrm{gh}}(H;v)=\mathrm{\Pi }(H,\tau )^1(\text{coefficient of }(|0^{\mathrm{gh}})^N\text{ in }v).$$
(ii) $`𝒞^{\mathrm{bos},\mathrm{hor}}(D,𝛔_{\stackrel{}{\mu }dx})`$ is one-dimensional if and only if there exists $`\mu (t)dt(\pi _{𝔛/S})_{}(𝔥^{}\mathrm{\Omega }_𝔛^1(D))`$ such that each $`\mu _i(x_i)dx_i`$ is a Laurent expansion of $`\mu (t)dt`$ at $`t=z_i`$ with respect to $`x_i=tz_i`$. In this case, $`𝒞^{\mathrm{bos},\mathrm{hor}}(D,𝛔_{\stackrel{}{\mu }dx})=\mathrm{\Phi }^{\mathrm{bos}}(v)`$, where
$$\mathrm{\Phi }^{\mathrm{bos}}(v)=(\text{coefficient of }|\stackrel{}{\mu }dx\text{ in }v).$$
Otherwise, $`𝒞^{\mathrm{bos}}(D,𝛔_{\stackrel{}{\mu }dx})=0`$.
Assume $`\lambda _i`$ ($`i=1,\mathrm{},N`$) satisfies the condition (1.8), and put $`\mu _i=\lambda _i`$ ($`i=1,\mathrm{},N`$), $`\mu _{N+j}=\alpha _{i(j)}`$ ($`j=1,\mathrm{},M`$). Lemma 4.2 guarantees that $`𝒞^{\mathrm{free},\mathrm{hor}}(D+D^{},𝝈_{\stackrel{}{\mu }dx})`$ is one-dimensional if we define $`D=P_1+\mathrm{}+P_N`$ ($`P_i`$ has the coordinate $`z_i`$), $`D^{}=Q_1+\mathrm{}+Q_M`$ ($`Q_j`$ has the coordinate $`t_j`$) and $`\mu _i(x)dx`$ ($`i=1,\mathrm{},N+M`$) as the Laurent expansion of $`\mu (t)dt`$ defined by
(4.11)
$$\mu (t)dt=\underset{i=1}{\overset{N}{}}\lambda _i\zeta _{11}(tz_i)\underset{j=1}{\overset{M}{}}\alpha _{i(j)}\zeta _{11}(tt_j).$$
The basis of this one-dimensional space is $`\mathrm{\Phi }^{\mathrm{free}}(H;v)=\mathrm{\Phi }^{\mathrm{gh}}(H;v)\mathrm{\Phi }^{\mathrm{bos}}(v)`$.
We assume that the parameters $`z_i`$ ($`i=1,\mathrm{},N`$) and $`t_j`$ ($`j=1,\mathrm{},M`$) satisy the Bethe Ansatz equations (1.9). Then by Lemma 2 of \[FFR\] (or by Corollary 5.2 of \[K\]), the screening vector $`scr_j`$ (cf. (5.21) of \[K\]) in $`\mathrm{W}ak_{\mu _{N+j}(x)dx}`$ is a singular vector of imaginary weight. Thanks to this property and (3.10), the linear functional defined by
(4.12)
$$\mathrm{\Psi }(H;v):=\mathrm{\Phi }^{\mathrm{free}}(H;vscr_1\mathrm{}scr_M)$$
for $`v_{i=1}^N\mathrm{W}ak_{\mu _i(x)dx}^0_{i=1}^NM_{\lambda _i}^{}`$ has the same property as (2.18):
(4.13)
$$\stackrel{~}{\mathrm{\Psi }}(H;vS[2]|0)=\stackrel{~}{\tau }(u)\mathrm{\Psi }(H;v).$$
On the other hand, we can compute the left hand side of (4.13), using the expression (3.9) and the Ward identity (2.9) for the free boson. The result is
(4.14)
$$\stackrel{~}{\mathrm{\Psi }}(H;vS[2]|0)=\tau _\mathrm{\Psi }(u)\mathrm{\Psi }(H;v),$$
namely $`\mathrm{\Psi }(H;v)`$ is the eigenvector of $`\widehat{\tau }(u)`$ with the eigenvalue $`\tau _\mathrm{\Psi }(u)`$ defined by (1.12).
The explicit form of $`\mathrm{\Psi }(H;v)`$, (1.10), is derived by the same argument as that of the rational case. See Lemma 3 of \[FFR\].
## Appendix A Elliptic functions
We follow the notations of \[M\] for theta functions and denote the odd theta function by
(A.1)
$$\theta _{11}(z,\tau )=\underset{n}{}\mathrm{exp}\left(\pi i\tau \left(n+\frac{1}{2}\right)^2+2\pi i\left(z+\frac{1}{2}\right)\left(n+\frac{1}{2}\right)\right).$$
We use a multiplicatively and additively quasi-periodic function,
(A.2)
$$w_c(z):=\frac{\theta _{11}^{}(0)\theta _{11}(zc)}{\theta _{11}(z)\theta _{11}(c)},\zeta _{11}(z):=\frac{d}{dz}\mathrm{log}\theta _{11}(z).$$
These functions are characterised by the properties
(A.3)
$$w_c(z+1)=w_c(z),w_c(z+\tau )=e^{2\pi ic}w_c(z),w_c(z)z^1\text{ around }z=0.$$
(A.4)
$$\zeta _{11}(z+1)=\zeta _{11}(z),\zeta _{11}(z+\tau )=\zeta _{11}(z)2\pi i,\zeta _{11}(z)z^1\text{ around }z=0.$$
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# 1 Introduction
## 1 Introduction
String field theory has been one of the most fundamental and the most mysterious subjects in string theory. In the course of the development, it has been clarifying the gauge interactions among higher excited states , the moduli problem at least for the open string . Originally it was regarded as the only candidate to describe the non-perturbative aspects of string theory.
The revolutionary developments in these years, however, the new ideas such as D-brane or M-theory turned out to play more fundamental rôle. One of the shortcoming of string field theory may be that it does not has direct means to describe D-branes dynamics although there were some attempts . In this respect, people pay more interests in the alternative approaches such as the matrix models where D-brane itself becomes the dynamical variable.
Some years ago, a novel approach to string field theory was evolved from the matrix model view point. In the infrared, the theory is described as a conformal field theory on the symmetric product $`S^nM=M^n/S_n`$. The orbifold singularities of the target are described as the twisted sectors. Excitations belonging to each twisted sectors can be physically interpreted as the collections of the “long strings” which are composite of the fundamental string variable (“short string” or “string bit”). The theory therefore contains a mechanism of the splitting/joining interactions of the closed string naturally in its definition. The partition function is expressed as the exponential of the one body partition function of the long string . This is typical structure of the partition function of the quantum field theory in the discrete lightcone quantization (DLCQ) which is not restricted to the string theory. This is the analog of the fact that the vacuum amplitude can be expressed as the exponential of the contributions from the connected Feynmann diagrams in the conventional quantum field theories. These facts support the idea that the matrix string theory describes the string field theory in DLCQ.
In this direction, a steady progress was made. For example, the four point amplitudes of the string theory was directly calculated by this method to reproduce Virasoro amplitude. It is applied to describe the little string theory to reproduce the black-hole entropy formula. It was generalized to heterotic matrix strings to describe the second quantized lightcone heterotic string. Some aspects of the $`S_N`$ orbifold CFT such as the modular properties and the fusion rule coefficients are studied in . A lot of the developments are made in the context of the moduli space. In particular, the instanton sectors of two dimensional Yang-Mills theory is related to the nontrivial topology of the matrix string world sheet .
From the mathematical viewpoint, it is originated from the calculation of the elliptic genera and has a direct relation with Götsche’s formula for Hilbert scheme of points and generalized Kac-Moody algebras.
In this paper, we study the open string version of the matrix string theory. The motivation of this subject should be obvious since we can not escape from dealing with D-branes in the matrix strings. We use the explicit calculation based on the boundary conformal field theory on the orbifold and give the some of the explicit analysis which should be made in BCFT. The new material is the appearance of the long open (closed) strings in the twisted sector of the open string<sup>1</sup><sup>1</sup>1 While we are finishing this manuscript, we noticed the work by Johnson where the notion of the long open strings as the twisted sectors was already mentioned. His strategy is to split the closed string amplitude into a product of the open string amplitude .. We give the classification theorem of all the possible form of such twisted sectors. We calculate the partition function for each twisted boundary conditions and show that it can be reducible to the amplitude of the one long string. An interesting feature is the world sheet topology of the short string is in general different from that of the long string. We develop also the boundary state formalism and reproduced the amplitude. If we sum up all possible boundary conditions, the partition function can be written as the exponential of the sum of the long string partition functions. This is the typical form of the partition function in the discrete lightcone gauge. Finally we confirmed that the dilaton tadpole cancellation occurs when the gauge group is famous $`SO(2^{13})`$ for the bosonic string.
Let us explain the organization of this paper. We put the main claims at the beginning of each section. One may first read these parts and skip the detailed explanation or the proof until it becomes necessary.
In section 2, we give a review of the basic structure of the orbifold theory on the symmetric product. We describe it in detail since some of the explicit calculations become essential later. We emphasize the aspect that it can be formulated as the conventional non-abelian orbifold theory. Namely for the torus amplitude, the consistency conditions for the twisted boundary condition contain all the information necessary to reproduce its characteristic feature of the string field theory. We also give a review of the discrete lightcone gauge and derived the typical form of its partition function.
In section 3, we investigate the constraint on the twisted boundary conditions for annulus/Möbius strip/Klein bottle amplitudes and relate it to the various boundary/cross-cap states . Open string twisted sectors were discussed in literature mainly for the abelian case. For non-abelian case, we need some extra care because of the non-commutativity. Because the open string twisted sector is the main object in this paper, we describe it in detail. The content of this section is generic and can be applied to arbitrary non-abelian orbifold models.
In section 4, we exactly solve the constraints for the symmetric product orbifold. The solution for the Klein bottle amplitude is similar to the torus case. In the Annulus and Möbius strip cases, there are some extra series of the solutions which will be interpreted as the contributions of the closed string sector. The content of this section is mathematical but is one of our main claim in this article.
In section 5, we calculate the one loop amplitude in two ways. First we do it by using the explicit operator formalism for the flat background. The calculation itself is technically similar to that of section 2. Second we calculate it by using multiple cover of the world sheet. The abstract combinatorial solutions in section 4 are translated into the form of the physically clear interpretation as the long string amplitudes. One interesting feature is the appearance of four sectors for the long string in each of the annulus and Möbius string amplitude for the short string. Namely, the topology of the world sheet seen from the long string is in general different from that for the short string as we already mentioned.
In section 6, we give the explicit form of the boundary/cross-cap states for the arbitrary twisted sectors. We argue that the boundary state for the short string can be classified into three types. The first one is the conventional boundary state for the long string. The second one turns out to be the cross-cap state for the long string. This is one of the origin of the topology change. The third one describes the connection of the two short strings. It encodes the nature of the string field theory of the orbifold theory. The cross-cap states for the short string have the similar classification. We calculate the inner product between them in order to examine the “modular invariance” or the tadpole condition in the next section.
In section 7, we first prove that each of three open string sectors can be expressed as the exponential of the one-body amplitudes of the long strings of the various scale. This is quite natural as the DLCQ partition function. The annulus amplitude for the short string contains all types of the amplitudes for the long string. One interesting aspect is that the torus amplitude which is contained there has the complex (but discrete) moduli parameter while the original annulus has only imaginary part.
In section 8, we examine the tadpole cancellation of the bosonic string in our context. We use only the annulus amplitude (for short string) to derive the tadpole condition. In this case there are cancellations among the massless parts of the boundary states (of the short string) alone. Since one-string partition function is exponentiated, one may reduce the tadpole condition to each one body problem for the long string. In this form, one can immediately reproduce the famous relation such as $`SO(2^{13})`$.
## 2 Review of closed string on symmetric product
### 2.1 Orbifold CFT on symmetric product
Let $`𝐗^I(\tau ,\sigma )`$ ($`𝐗^IM`$, $`I=1,\mathrm{},N`$) be the bosonic coordinates which define the string embedding on symmetric product $`S^NM`$. The twisted sectors $`_h`$ ($`hS_N`$) are defined by the boundary condition of $`𝐗`$,
$$𝐗(\sigma _0,\sigma _1+2\pi )=h𝐗(\sigma _0,\sigma _1),$$
(2.1)
where $`(h𝐗)^I_Jh^{IJ}𝐗^J`$.
The modular invariant partition function on torus is,
$`Z_N(\tau ,\overline{\tau })`$ $`=`$ $`{\displaystyle \frac{1}{|S_N|}}{\displaystyle \underset{\begin{array}{c}g,hS_N\\ gh=hg\end{array}}{}}\chi _{h,g}(\tau ,\overline{\tau })`$ (2.4)
$`\chi _{h,g}(\tau ,\overline{\tau })`$ $`=`$ $`\text{Tr}__h(g𝐞\left[\tau L_0\overline{\tau }\overline{L_0}\right]),`$ (2.5)
$`𝐞\left[x\right]`$ $``$ $`e^{2\pi ix}.`$ (2.6)
The summation in $`gS_N`$ is needed to define a projection onto the $`S_N`$ invariant subspace. The constraint
$$gh=hg$$
(2.7)
in (2.4) is the consistency condition of the path integral to assure that the twists in time and space directions commute.
In the non-abelian orbifold, only the conjugacy class of $`h`$ has the invariant meaning since $`g𝐗_{ghg^1}`$ if $`𝐗_h`$. The summation in (2.4) over $`h`$ is then replaced by the summation over the conjugacy class $`𝒞_i`$ of $`S_N`$. For a particular element $`h𝒞_i`$, the solutions of (2.7) are the elements of the centralizer group $`𝒩_i`$. With the relation $`|S_N|=|𝒞_i||𝒩_i|`$, (2.7) can be rewritten as,
$$Z_N(\tau ,\overline{\tau })=\underset{i}{}\frac{1}{|𝒩_i|}\underset{g𝒩_i}{}\chi _{h,g}(\tau ,\overline{\tau }),h𝒞_i$$
(2.8)
This formula is the generic expression for the arbitrary non-abelian orbifold.
For the permutation group $`S_N`$, the conjugacy group is labeled by the partition of $`N`$, since any group element can be written as a product of elementary cycles $`(n)`$ of length $`n`$,
$$[h]=(1)^{N_1}(2)^{N_2}\mathrm{}(k)^{N_k},\underset{n>0}{}nN_n=N.$$
(2.9)
The centralizer of such an element is a semi-direct product of factors $`S_{N_n}`$ and $`𝐙_n`$,
$$g𝒩_h=S_{N_1}\times (S_{N_2}𝐙_2^{N_2})\times \mathrm{}(S_{N_k}𝐙_k^{N_k}).$$
(2.10)
The factors $`S_{N_n}`$ permute the $`N_n`$ cycles $`(n)`$, while the factors $`𝐙_n`$ rotate each cycle $`(n)`$. The order of the centralizer group is,
$$|𝒩_h|=\underset{n=1}{\overset{k}{}}n^{N_n}N_n!$$
(2.11)
The physical interpretation of these factors are well-known. Let us first consider the case where $`h`$ is the element of the cyclic permutation of $`N`$ elements,
$$h=𝒯_N:𝐗^I𝐗^{I+1},I\{0,1,\mathrm{},N1\}.$$
(2.12)
Here superscript $`I`$ is defined by mod $`N`$. The twist by $`h`$ then gives the boundary condition,
$$𝐗^I(\sigma _0,\sigma _1+2\pi )=𝐗^{I+1}(\sigma _0,\sigma _1).$$
(2.13)
It means that $`N`$ short closed strings $`𝐗^I`$ are connected with each other to form one long string of length $`n`$. For the general situation (2.9), we will have $`N_n`$ long strings of length $`n`$ for $`n=1,\mathrm{},k`$. The short strings that form a long string are sometimes called “string bits”.
In this language, the element of the centralizer group (2.10) has a clear interpretation. $`𝐙_n`$ factors are the rotations of the string bits that constitute a long string of length $`n`$. $`S_{N_n}`$ then reshuffle the long strings of the same length $`n`$ as a whole (figure 1).
### 2.2 Partition function
The long string interpretation can be established further by the calculation of the partition function. In this subsection, we would like to give somewhat more explicit computation compared with the literature for the preparation for later sections. To make our argument clear, we first restrict the situation where target space is flat $`M=𝐑^1`$. We will then give a generic argument for the arbitrary target space.
#### 2.2.1 Irreducible Diagram
In the calculation of the partition function (2.8), we need to divide $`N`$ free fields into small subgroups. Each subgroup consists of the collection of free fields which are mixed up by the action of two twists $`h`$ and $`g`$. It is obvious that there are minimal sets of free field which can not divided into subgroups. We call such a set of free fields as “irreducible set”. The partition function $`\chi _{h,g}`$ is identified as the product of the contributions of each irreducible sets.
In the torus case, it is known that such irreducible sets can be always reduced to $`nm`$ free fields where $`h`$ and $`g`$ act as,
$`h`$ $`=`$ $`diag(𝒯_n,\mathrm{},𝒯_n)`$
$`g`$ $`=`$ $`diag(𝒯_n^{p_1},\mathrm{},𝒯_n^{p_m})\overline{𝒯}_m,`$ (2.14)
where $`\overline{𝒯}_m`$ acts on $`𝐗^{I,J}`$ as $`(\overline{𝒯}_m𝐗)^{I,J}=𝐗^{I,J+1}`$ and $`𝒯_n`$ acts as $`(𝒯_n𝐗)^{I,J}=𝐗^{I+1,J}`$. We label $`nm`$ free fields as $`𝐗^{I,J}`$ ($`I(0,1,\mathrm{},n1)`$ and $`J(0,1,\mathrm{},m1)`$). $`I`$ (resp. $`J`$) is defined modulo $`n`$ (resp. $`m`$). $`p_J`$’s represent the rotations each long strings and take their values in $`0,1,\mathrm{},n1`$.
The action of $`h`$ and $`g`$ on $`𝐗`$ in the component is given as,
$`(h𝐗)^{I,J}`$ $`=`$ $`𝐗^{I+1,J}`$
$`(g𝐗)^{I,J}`$ $`=`$ $`𝐗^{I+p_J,J+1}`$ (2.15)
To make $`h`$ action diagonal, we introduce the discrete Fourier transformation with respect to $`I`$,
$$\stackrel{~}{𝐗}^{a,J}\frac{1}{\sqrt{n}}\underset{I=0}{\overset{n1}{}}𝐞\left[\frac{aI}{n}\right]𝐗^{I,J},a(0,1,\mathrm{},n1)$$
(2.16)
In this basis, the actions of $`h,g`$ are modified to,
$`(h\stackrel{~}{𝐗})^{a,J}`$ $`=`$ $`𝐞\left[{\displaystyle \frac{a}{n}}\right]\stackrel{~}{𝐗}^{a,J}`$
$`(g\stackrel{~}{𝐗})^{a,J}`$ $`=`$ $`𝐞\left[{\displaystyle \frac{ap_J}{n}}\right]\stackrel{~}{𝐗}^{a,J+1}`$ (2.17)
Since $`h`$ action is diagonalized, the periodicity of $`\stackrel{~}{𝐗}`$ becomes well-defined,
$$\stackrel{~}{𝐗}^{a,J}(\sigma _0,\sigma _1+2\pi )=𝐞\left[\frac{a}{n}\right]\stackrel{~}{𝐗}^{a,J}(\sigma _0,\sigma _1).$$
(2.18)
We have the mode expansion of $`\stackrel{~}{𝐗}^{a,J}`$ as ($`z=e^{i(\sigma _1+i\sigma _0)}`$),
$$\stackrel{~}{𝐗}^{a,J}=\alpha _0^J\delta _{a,0}\sigma _0+i\underset{r𝐙}{}\left(\frac{1}{ra/n}\alpha _{ra/n}^{a,J}z^{r+a/n}+\frac{1}{r+a/n}\stackrel{~}{\alpha }_{r+a/n}^{a,J}\overline{z}^{ra/n}\right)$$
(2.19)
with the commutation relation,
$$[\alpha _{ra/n}^{a,I},\alpha _{sb/n}^{b,J}]=(ra/n)\delta _{I,J}\delta _{r+s}\delta _{a+b}.$$
(2.20)
In order to diagonalize $`g`$ action, we combine the oscillators further as,
$$𝐘^a\underset{J=0}{\overset{m1}{}}C_J\stackrel{~}{𝐗}^{a,J}.$$
(2.21)
Eigenstate equation $`g𝐘^a=\mu 𝐘^a`$ gives a relation between the neighboring coefficients,
$$\mu C_{J+1}=𝐞\left[\frac{ap_J}{n}\right]C_J.$$
(2.22)
The periodicity condition $`C_m=C_0`$ implies,
$$\mu ^m=𝐞\left[\frac{ap}{n}\right],p\underset{J}{}p_J$$
(2.23)
The eigenvalues for $`g`$ action are evaluated as,
$$\mu ^{a,b}=𝐞\left[\frac{b}{m}+\frac{ap}{nm}\right],a=0,1,\mathrm{},n1,b=0,1,\mathrm{},m1.$$
(2.24)
In other word, the action of $`g`$ is diagonalized as
$$(g\alpha _{ra/n})^{a,b}=\mu ^{a,b}\alpha _{ra/n}^{a,b},$$
(2.25)
where $`\alpha ^{a,b}`$ are the linear combinations of $`\alpha ^{a,J}`$.
The (chiral) oscillator part of the partition function can be now evaluated by using mode expansion and eigenvalues of $`g`$,
$`{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{a=0}{\overset{n1}{}}}{\displaystyle \underset{b=0}{\overset{m1}{}}}\left(1𝐞\left[{\displaystyle \frac{b}{m}}+{\displaystyle \frac{ap}{nm}}\right]𝐞\left[(ra/n)\tau \right]\right)^1`$
$`={\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{a=0}{\overset{n1}{}}}\left(1𝐞\left[{\displaystyle \frac{ap}{n}}\right]𝐞\left[m\tau (ra/n)\right]\right)^1`$
$`={\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{a=0}{\overset{n1}{}}}\left(1𝐞\left[p(r{\displaystyle \frac{a}{n}})\right]𝐞\left[m\tau (ra/n)\right]\right)^1`$
$`={\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}\left(1𝐞\left[\left({\displaystyle \frac{m\tau p}{n}}\right)r\right]\right)^1`$ (2.26)
Thus at least the chiral part of the partition function from $`nm`$ string bits are knitted together to give the partition function of one long string with modified moduli parameter
$$\frac{m\tau p}{n}\tau _{n,m,p}.$$
(2.27)
As for the momentum integration, only the $`a=0`$ component have zero mode. In the remaining $`m`$ momentum, only $`b=0`$ part gives non-trivial inner product when $`g`$ is inserted. We conclude that zero-mode contribution comes from integration of single momentum. One subtlety is how to fix the normalization constant. While changing the variable from the short string to the long one, Virasoro generators should be modified to $`L_0L_0/n`$ . Next since we have $`m`$ of such long strings moving coherently, we have to multiply it $`m`$. The kinetic term is thus modified to $`\frac{mp^2}{2n}`$. By integrating out $`p`$, we get
$$1/\sqrt{\text{Im}(\tau _{n,m,p})}.$$
(2.28)
By combining the contributions from the anti-chiral part and momentum integration, the partition function of $`nm`$ free fields becomes
$$Z(n,m,\left\{p\right\}|\tau ,\overline{\tau })=Z_1(\tau _{n,m,p},\overline{\tau }_{n,m,p}).$$
(2.29)
$`Z_1(\tau ,\overline{\tau })`$ is the standard partition function for one free boson,
$$Z_1(\tau ,\overline{\tau })=\frac{1}{\sqrt{\text{Im}\tau }}(\eta (\tau )\overline{\eta }(\overline{\tau }))^1.$$
(2.30)
As we see in the next paragraph, (2.29) is the generic feature of the symmetric space orbifold. It may be physically interpreted as a kind of the renormalization. Namely path integral over the short string variables are replaced by that of the long string. The two formula coincides if we replace the moduli parameter.
#### 2.2.2 Generalization of the target space
Our consideration in the previous subsection essentially depends on the flatness of the target space. We first review the general arguments on the arbitrary target space by using the path integral method.
Let $`M`$ be a general manifold where a string can live and consider the path integral on $`𝐗^{I,J}M`$ which satisfies
$$𝐗^{I,J}(w+1)=𝐗^{I+1,J}(w),𝐗^{I,J}(w+\tau )=𝐗^{I+p_J,J+1}(w),$$
(2.31)
for each irreducible set $`h,g`$ as in the previous subsection. Here $`𝐗`$ is defined on a small parallelogram with period $`(1,\tau )`$.
These fields $`𝐗(w)`$ can be rewritten from one field $`𝒳(w)M`$ by the identification,
$$𝐗^{I,J}(w)=𝒳\left(w+J\tau +IP_J\right),P_J=\underset{\mathrm{}=0}{\overset{J1}{}}p_{\mathrm{}}.$$
(2.32)
$`𝒳`$ has the following periodicity derived from $`𝐗^{I+n,J}=𝐗^{I,J+m}=𝐗^{I,J}`$,
$$𝒳(w+n)=𝒳(w),𝒳(w+m\tau p)=𝒳(w).$$
(2.33)
Namely, $`𝒳(w)`$ is defined on $`nm`$ combinations of the parallelogram with period $`(n,m\tau p)`$. The moduli of the bigger torus is given by $`\tau _{n,m,p}`$.
Path integral over variable $`𝐗`$ should be the same as that of $`𝒳`$. It implies the generic rule for the arbitrary target space,
$$Z_{h,g}^𝐗(\tau ,\overline{\tau })=Z^𝒳(\tau _{n,m,p},\overline{\tau }_{n,m,p}).$$
(2.34)
This type of the proof may look too abstract. One may give more concrete reasoning when the target space is the orbifold $`M=T/\mathrm{\Gamma }`$ where $`T`$ is a flat space and $`\mathrm{\Gamma }`$ is a discrete group which may be non-abelian. The total target space becomes $`T^N/S_N\mathrm{\Gamma }^N`$. The partition function is written as
$$Z_N=\frac{1}{|S_N||\mathrm{\Gamma }|^N}\underset{\begin{array}{c}h,gS_N\\ hg=gh\end{array}}{}\underset{\begin{array}{c}\left\{\alpha \right\},\left\{\beta \right\}\mathrm{\Gamma }\\ \left\{\alpha \right\}\left\{\beta \right\}^h=\left\{\beta \right\}\left\{\alpha \right\}^g\end{array}}{}\text{Tr}_{_{h,\left\{\alpha \right\}}}\left(g\beta q^{L_0}\overline{q}^{\overline{L}_0}\right).$$
(2.35)
Here $`\left\{\alpha \right\},\left\{\beta \right\}`$ define the boundary condition for $`𝐗`$ for fixed $`h,gS_N`$ as,
$$𝐗^i(w+1)=\alpha ^i𝐗^{h(i)}(w),𝐗^i(w+\tau )=\beta ^i𝐗^{g(i)}(w).$$
(2.36)
The condition $`\left\{\alpha \right\}\left\{\beta \right\}^h=\left\{\beta \right\}\left\{\alpha \right\}^g`$ in the summation signifies the “integrability condition” of these boundary conditions,
$$\alpha ^i\beta ^{h(i)}=\beta ^i\alpha ^{g(i)},i=1,2,\mathrm{},N.$$
(2.37)
For irreducible variables with $`h,g`$ given as (2.2.1), we replace the index $`i`$ to a pair $`I,J`$ and rewrite the boundary conditions and consistency conditions
$`𝐗^{I,J}(w+1)=\alpha ^{I,J}𝐗^{I+1,J}(w),𝐗^{I,J}(w+\tau )=\beta ^{I,J}𝐗^{I+p_J,J+1}(w),`$
$`\beta ^{I,J}\alpha ^{I+p_J,J+1}=\alpha ^{I,J}\beta ^{I+1,J}.`$ (2.38)
With this type of the constraint, one can always find a unique set $`\gamma ^{I,J}\mathrm{\Gamma }`$ which satisfies,
$$\gamma ^{0,0}=1,\alpha ^{I,J}=\gamma ^{I,J}(\gamma ^{I+1,J})^1,\beta ^{I,J}=\gamma ^{I,J}(\gamma ^{I+p_J,J+1})^1.$$
(2.39)
With this twist factor, one may introduce $`𝒳(w)T`$ as before which is defined on the bigger torus generated by $`(n,m\tau p)`$ and relate it $`𝐗^{IJ}`$ as
$$𝐗^{I,J}(w)=\gamma ^{IJ}𝒳(w+J\tau +IP_J).$$
(2.40)
It is easy to derive that $`𝐗`$ thus defined satisfies (2.38). In terms of $`𝒳`$ one may develop the operator formalism as before. Since $`\gamma `$ factors are just redefinition of the identification between the variables, The partition function with different $`\gamma `$ will produce identical partition function. It will give a factor $`|\mathrm{\Gamma }|^{nm1}`$. The only nontrivial element is the global monodromy factor $`A,B\mathrm{\Gamma }`$ defined by,
$$A=\gamma ^{n,0},B=\gamma ^{0,m},$$
(2.41)
which defines the boundary condition for $`\chi (w)`$,
$$𝒳(w+n)=A^1𝒳(w),𝒳(w+m\tau p)=B^1𝒳(w).$$
(2.42)
We thus arrive at the desired relation,
$$\frac{1}{|\mathrm{\Gamma }|^{nm}}\underset{\begin{array}{c}\left\{\alpha \right\},\left\{\beta \right\}\mathrm{\Gamma }\\ \left\{\alpha \right\}\left\{\beta \right\}^h=\left\{\beta \right\}\left\{\alpha \right\}^g\end{array}}{}Z_{h,\left\{\alpha \right\},g,\left\{\beta \right\}}^𝐗(\tau ,\overline{\tau })=\frac{1}{|\mathrm{\Gamma }|}\underset{\begin{array}{c}A,B\mathrm{\Gamma }\\ AB=BA\end{array}}{}Z_{A,B}^𝒳(\tau _{n,m,p},\overline{\tau }_{n,m,p}).$$
(2.43)
The partition function for $`𝒳`$ can be obtained by the explicit operator formalism and obviously identical to the single string partition function. One may conclude that the formula (2.29) should hold for any target space by replacing $`Z_1`$ by the partition function of the single string in that target space.
#### 2.2.3 Generating function of the partition function
In order to evaluate the whole partition function (2.4), we need to follow some steps.
1. As we mentioned in the previous subsection, $`\chi _{h,g}`$ in (2.4) is expressed as the product of the contributions from the irreducible sets (2.29). We need to specify how many of such irreducible sets are contained in given pair $`h,g`$.
2. As for $`h`$, we use the representation of the conjugacy class (2.9) and use the definition (2.8).
3. For each $`h`$, $`g`$ is an element of the centralizer group (2.10). If we restrict the sector $`h=𝒯_n^{N_n}`$, $`g`$ can be written as $`S_{N_n}𝐙_n^{N_n}`$. Elements in $`S_{N_n}`$ should be again decomposed into conjugacy class as $`(1)^{M_{n,1}}(2)^{M_{n,2}}(3)^{M_{n,3}}\mathrm{}`$ with the constraint,
$$\underset{s_n=1}{\overset{\mathrm{}}{}}s_nM_{n,s_n}=N_n.$$
(2.44)
4. In this decomposition, we have $`M_{n,m}`$ irreducible subsets for each $`n,m`$. By summation over elements in $`𝐙_n`$ (namely $`p_{\mathrm{}}`$ in (2.2.1)), we get the expression of $`\chi _{h,g}`$,
$`\chi _{h,g}`$ $`=`$ $`{\displaystyle \underset{n,m=1}{\overset{\mathrm{}}{}}}(n^mZ(n,m))^{M_{n,m}}`$ (2.45)
$`Z(n,m)`$ $`=`$ $`{\displaystyle \frac{1}{n^m}}{\displaystyle \underset{\left\{p\right\}}{}}Z(n,m,\left\{p\right\}|\tau ,\overline{\tau })={\displaystyle \frac{1}{n}}{\displaystyle \underset{p}{}}Z(n,m,p|\tau ,\overline{\tau }).`$ (2.46)
We used the fact that the partition function (2.29) depends only on the sum of $`\left\{p\right\}`$.
5. For each, $`(\left\{N\right\},\left\{M\right\})`$, the weight factor is
$$\frac{1}{_{n=1}^{\mathrm{}}n^{N_n}N_n!}\underset{n=1}{\overset{\mathrm{}}{}}\left(\frac{N_n!}{_{s_n=1}^{\mathrm{}}\left(s_n^{M_{n,s_n}}M_{n,s_n}!\right)}\right)$$
(2.47)
The first term is the order of the centralizer group in (2.4). The second term is the product of the number of the conjugacy class for each $`\left\{M\right\}`$. The factor $`N_n!`$ is canceled between these two terms. The factor $`n^{N_n}`$ can be absorbed in $`n^m`$ in (2.46) since $`N_n=_mmM_{n,m}`$.
6. Assembling every term, we have the following combination,
$$Z_N(\tau ,\overline{\tau })=\underset{\begin{array}{c}N_n,M_{n,m}1\\ {\scriptscriptstyle nN_n}=N\\ {\scriptscriptstyle mM_{n,m}}=N_n\end{array}}{}\underset{n,m1}{}\left(\frac{1}{m}Z(n,m|\tau ,\overline{\tau })\right)^{M_{n,m}}.$$
(2.48)
7. In the generating function of the partition function,
$$Z(\zeta |\tau ,\overline{\tau })\underset{N=0}{\overset{\mathrm{}}{}}\zeta ^NZ_N(\tau ,\overline{\tau }),$$
the constraint in (2.48) can be removed to give a free summation over $`\left\{M\right\}`$,
$`Z(\zeta |\tau ,\overline{\tau })`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{M_{n,m}=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{M_{n,m}!}}\left({\displaystyle \frac{1}{m}}Z(n,m|\tau ,\overline{\tau })\zeta ^{nm}\right)^{M_{n,m}}`$ (2.49)
$`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}\mathrm{exp}\left({\displaystyle \frac{1}{m}}Z(n,m|\tau ,\overline{\tau })\zeta ^{nm}\right)`$
$`=`$ $`\mathrm{exp}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{nm}}{\displaystyle \underset{p=0}{\overset{n1}{}}}Z(n,m,p|\tau ,\overline{\tau })\zeta ^{nm}\right)`$
$`=`$ $`\mathrm{exp}\left({\displaystyle \underset{N=1}{\overset{\mathrm{}}{}}}\zeta ^N𝒱_NZ_1(\tau ,\overline{\tau })\right),`$
$`𝒱_Nf(\tau ,\overline{\tau })`$ $``$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{\begin{array}{c}a,d=1,\mathrm{},N\\ b=0,\mathrm{},d1\\ ad=N\end{array}}{}}f({\displaystyle \frac{a\tau +b}{d}},{\displaystyle \frac{a\overline{\tau }+b}{d}})`$ (2.53)
The operator $`𝒱_N`$ which appeared in the final expression is called Hecke operator. It maps a modular form to another one with the same weight and it typically appears in this type of calculation.
We may summarize the computation in this subsection into a theorem,
Theorem 1: The generating function of the partition functions (2.4) are given in the form,
$$\underset{N=0}{\overset{\mathrm{}}{}}\zeta ^NZ_N(\tau ,\overline{\tau })=e^{(\zeta |\tau ,\overline{\tau })},(\zeta |\tau ,\overline{\tau })=\underset{N=1}{\overset{\mathrm{}}{}}\zeta ^N𝒱_NZ_1(\tau ,\overline{\tau }).$$
(2.54)
$``$ may be regarded as the free energy which consists of the contributions from the irreducible diagrams. The partition function $`𝒱_NZ_1`$ for each irreducible diagram (which is the contribution from $`N`$ string bits) is given by applying Hecke operator (2.53) to the modular invariant partition function for the single string. It is physically interpreted as the contribution of a single “long” string.
### 2.3 Discrete lightcone quantization
The physical interpretation of theorem 1 as a string field theory can be established by introducing the notion of the discrete lightcone quantization (DLCQ) (see also for a review).
Consider first the usual quantum field theory (free boson theory) in $`d`$-dimension. We take $`x^+`$ as the time variable and write it as $`\tau `$. Lagrangian of the system is given by,
$$S=d^dx(_\tau \mathrm{\Phi }_{}\mathrm{\Phi }+\frac{1}{2}\underset{i=2}{\overset{d1}{}}_i\mathrm{\Phi }_i\mathrm{\Phi })$$
(2.55)
In DLCQ, the light-like direction is compactified by $`S^1`$ of radius $`R`$,
$$x^{}x^{}+2\pi R.$$
(2.56)
Consequently the momentum along that direction is quantized
$$p^+=N/R.$$
(2.57)
We expand the (correctly normalized) wave functions along the transverse coordinates by the eigenfunction of Hamiltonian for the transverse degree of freedom,
$$H^t\psi _a(𝐱)=\lambda _a\psi _a(𝐱),a\mathrm{\Lambda }$$
(2.58)
where $`𝐱`$ is the transverse coordinates and $`\mathrm{\Lambda }`$ is a set which labels the eigenfunctions. We introduce the partition function for the transverse degree of freedom,
$$Z_1(\tau )=\underset{a\mathrm{\Lambda }}{}𝐞\left[\tau \lambda _a\right].$$
(2.59)
Fourier transformation along $`x^{}`$ gives the mode expansion of $`\mathrm{\Phi }`$,
$$\mathrm{\Phi }(x)=\underset{n=0}{\overset{\mathrm{}}{}}\left(\mathrm{\Phi }_n^ae^{2\pi inx^{}/R}\psi _a^{}(𝐱)+\mathrm{\Phi }_n^a(𝐱)e^{2\pi inx^{}/R}\psi _a(𝐱)\right),$$
(2.60)
together with the commutation relation,
$$[\mathrm{\Phi }_n^a,\mathrm{\Phi }_m^b]=\frac{1}{n}\delta _{n,m}\delta _{a,b}.$$
(2.61)
Every oscillator has the index $`n`$ which describes the sectors classified by the lightcone momentum.
The Hilbert space $`_N`$ with definite total lightcone momentum $`p^+=N/R`$ is constructed out of these oscillators as,
$$_N\{\stackrel{N_1}{\stackrel{}{\mathrm{\Phi }_1^{a_1^1}\mathrm{}\mathrm{\Phi }_1^{a_{N_1}^1}}}\stackrel{N_2}{\stackrel{}{\mathrm{\Phi }_2^{a_1^2}\mathrm{}\mathrm{\Phi }_2^{a_{N_2}^2}}}\mathrm{}|0|a_s^{\mathrm{}}\mathrm{\Lambda },N_n0,\underset{n}{}nN_n=N.\}.$$
(2.62)
DLCQ partition function is defined as the trace over such Hilbert space with weight $`\zeta ^N`$,
$$Z(\zeta ,\tau )=\underset{N=0}{\overset{\mathrm{}}{}}\text{Tr}__N(\zeta ^{p^+}𝐞\left[\tau p^{}\right]),$$
(2.63)
where $`p^+=N/R`$. $`p^{}`$ can be expressed as $`p^{}=\frac{1}{p^+}H^t`$ from the on-shell condition. More explicitly we assign the weight factor $`\zeta ^{_nnN_n/R}𝐞\left[_{n,s}R\tau \lambda _{a_s^n}/n\right]`$ to the state in the brace in (2.62). In the following, we absorb $`R`$ by redefinition of $`\zeta ,\tau `$ and put $`R=1`$.
Since we have summation over $`N`$, it is clear that the partition function can be written as a product, $`Z(\zeta ,\tau )=_{n=1}^{\mathrm{}}Z_n(\zeta ,\tau ),`$ where $`Z_n(\zeta ,\tau )`$ is the contribution from the states generated by $`\mathrm{\Phi }_n^a`$ for fixed $`n`$. These sub-factors can be computed as follows,
$`Z_n(\zeta ,\tau )`$ $`=`$ $`{\displaystyle \underset{\left\{N_a\right\},a\mathrm{\Lambda }}{}}{\displaystyle \underset{a\mathrm{\Lambda }}{}}\zeta ^{nN_a}𝐞\left[\tau N_a\lambda _a/n\right]`$ (2.64)
$`=`$ $`{\displaystyle \underset{a\mathrm{\Lambda }}{}}{\displaystyle \frac{1}{1\zeta ^n𝐞\left[\tau \lambda _a/n\right]}}`$
$`=`$ $`\mathrm{exp}\left({\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\zeta ^{nm}}{m}}{\displaystyle \underset{a\mathrm{\Lambda }}{}}𝐞\left[\lambda _am/n\right]\right)`$
$`=`$ $`\mathrm{exp}\left({\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\zeta ^{nm}}{m}}Z_1(m\tau /n)\right).`$
In passing from the second to the third lines, we used the formula $`1/(1x)=e^{_{n=1}^{\mathrm{}}x^n/n}`$. Assembling these terms, we arrive at the formula which is already very similar to (2.54),
$$Z(\zeta ,\tau )=\mathrm{exp}(\underset{n,m=1}{\overset{\mathrm{}}{}}\frac{\zeta ^{nm}}{m}Z_1(m\tau /n)).$$
(2.65)
In case of the open string field theory in DLCQ, we expect exactly the same formula if we reinterpret $`\mathrm{\Lambda }`$ as the label for the Fock space of one-body open string generated by the transverse oscillators . We will prove it in section 7 by using only the combinatorics of the orbifold theory on the symmetric product.
For the closed string field theory, we have to take care of the winding mode for $`x^{}`$,
$$𝑑\sigma _\sigma X^{}(\sigma )=2\pi R\nu ,$$
(2.66)
where $`\nu `$ is the winding number. By using the relation $`_\sigma X^{}_\sigma X^i_\sigma X^i`$, Susskind argued that the transverse Virasoro generators must satisfy,
$$L_0\overline{L}_0=\nu n,$$
(2.67)
where $`n`$ is the eigenvalue for $`p^+=n/R`$. In this respect, we need to insert the projection operator which restrict the partition sum into the states satisfying $`L_0\overline{L}_00`$ mod $`n`$. This is achieved by the replacement,
$$Z(\frac{m\tau }{n},\frac{m\overline{\tau }}{n})\frac{1}{n}\underset{p=0}{\overset{n1}{}}Z(\frac{m\tau p}{n},\frac{m\overline{\tau }p}{n}).$$
(2.68)
It gives exactly (2.54).
At this point, we find that the partition function for the string theory in DLCQ coincides exactly to the partition function of orbifold theory if we identify $`M=𝐑^{24}`$ ($`M=𝐑^8`$ for the superstring). The dictionary between the two picture is first the length of the long string $`n`$ should be identified with the lightcone momentum $`n/R`$.
Comparing the formula in ordinary LCQ of string theory with the discretized formula, (2.54),
$$\frac{d^2\tau }{\text{Im}\tau }Z_1(\tau ,\overline{\tau })\underset{m=1}{\overset{\mathrm{}}{}}\frac{1}{mn}\underset{p=0}{\overset{n1}{}}Z_1(\tau _{n,m,p},\overline{\tau }_{n,m,p}),$$
(2.69)
it is natural to regard the summation over $`m,p`$ as the discrete version of the integral over $`\tau `$. Namely we have the identification, $`\text{Re}\tau p/n`$, $`\text{Im}\tau m/n`$ (see for example, ).
We note that the moduli parameter $`\tau `$ for the short string becomes rather “redundant” variable. We might say that we do not need integration over moduli for the constituent string bits and may simply set $`\tau =i`$.
## 3 Open string CFT on non-abelian orbifold
In this paper, we aim to extend the analysis in the previous section to the open string theory. For that purpose, we need to find the analog of the consistency condition (2.7) for various open string one-loop diagrams (Klein bottle, Annulus, Möbius strip). Although this is an elementary issue, we could not find a literature where such conditions for non-abelian orbifold were examined<sup>2</sup><sup>2</sup>2 Conditions for the abelian orbifold were studied in many papers. For the incomplete list, see . in detail. We note that an important difference between abelian and non-abelian situations is that we will have non-trivial open string twisted sectors in non-abelian case. In the abelian orbifold the twisted sector reduces to the standard Dirichlet-Neumann sector. On the other hand, in the permutation orbifold case, there are infinite varieties of open string twisted sectors where they will be interpreted to give the long open (or closed) strings.
In this section, we give a summary of the consistency conditions for the generic non-abelian orbifold $`M/\mathrm{\Gamma }`$ where $`\mathrm{\Gamma }`$ is a non-abelian discrete group.
### 3.1 Open string Hilbert space
The usual boundary conditions for the open string are the following two,
$$(𝐗_R(\sigma _0,\sigma _1)\pm 𝐗_L(\sigma _0,\sigma _1))|_{\sigma _1=0}=0,$$
(3.1)
where $`+`$ (resp. $``$) sign is assigned for Neumann (resp. Dirichlet) boundary condition if $`𝐗`$ is decomposed as $`𝐗=𝐗_L(z)+𝐗_R(\overline{z})`$. In orbifold case, they are generalized to the twisted reflection relation between the left and right movers
$$\left(𝐗_R(\sigma _0,\sigma _1)\pm f𝐗_L(\sigma _0,\sigma _1)\right)|_{\sigma _1=0}=0,f\mathrm{\Gamma }.$$
(3.2)
As long as the consistency of the boundary state which we will examine later, this condition is consistent for the arbitrary element $`f\mathrm{\Gamma }`$. However, such a general twist leads us to the unequal footing for the left and the right movers. This is the situation which appears in the asymmetric orbifold which we would not like discuss in this article. In this sense, we will impose a constraint on $`f`$,
$$f^2=1.$$
(3.3)
When $`f`$ is not identity in $`\mathrm{\Gamma }`$, taking Dirichlet type boundary condition means that the open string endpoint is fixed at the fixed point of the orbifold action. On the other hand, the interpretation of Neumann boundary condition is not very clear.
The open string Hilbert space is specified by the boundary twists at two boundaries $`\sigma _1=0,\pi `$. Let us write it as $`_{f_1,f_2}`$ where $`f_1\mathrm{\Gamma }`$ (resp.$`f_2`$) is the twist at $`\sigma _1=0`$ (resp. $`\sigma _1=\pi `$). We note that by changing the basis for $`𝐗`$, the boundary condition reduces to $`f_{1,2}=\pm 1`$, namely Dirichlet or Neumann boundary conditions for the abelian orbifold.
For the permutation group, the boundary condition (3.2) has a clear physical interpretation. Because of the condition (3.3), $`f`$ should belong to the conjugacy class of the form, $`(1)^{n_1}(2)^{n_2}`$. If $`f=𝒯_1`$, it simply means that the open string bit has Neumann or Dirichlet boundary conditions at that boundary. On the other hand, if $`f=𝒯_2`$ or the permutation of (12), the twisted boundary condition reads
$$𝐗_L^1\pm 𝐗_R^2|_{\sigma _1=0}=0,𝐗_L^2\pm 𝐗_R^1|_{\sigma _1=0}=0.$$
For Dirichlet type condition, it means that two open string bits are connected smoothly at that boundary.
In the following sections, we will mainly use Neumann boundary condition when $`f`$ defines the open boundary and Dirichlet condition when $`f`$ describes the connection of two edges. This restriction is not essential in our discussion and can be easily generalized to include the mixed boundary conditions.
By combining $`f_1`$ and $`f_2`$, one can realize many varieties of the configurations for the open string bits. For example, we illustrate the situation for a long open string (figure 2, left), $`f_1=(1)(23)(4)`$, $`f_2=(12)(34)`$, and for a long closed string (figure 2, right), $`f_1=(14)(23)`$, $`f_2=(12)(34)`$.
It is clear that by choosing a suitable pair, $`(f_1,f_2)`$, one may realize both the long open and closed strings of arbitrary length.<sup>3</sup><sup>3</sup>3For the long closed string, it should be constructed out of even number of open string bits.
### 3.2 Annulus diagram and boundary state
The annulus diagram is represented as the trace of open string Hilbert space $`_{f_1,f_2}`$ with the projection operator onto the $`\mathrm{\Gamma }`$ invariant state,
$$\text{Tr}_{_{f_1,f_2}}\left(g𝐞\left[\tau L_0\right]\right)$$
(3.4)
Here and in the following sections, the moduli parameter $`\tau `$ is pure imaginary.
It is straightforward to check that the periodicity along time direction (twist by $`g`$) and the twist at the boundary (3.2) is consistent only if,
$$gf_1=f_1g,gf_2=f_2g$$
(3.5)
For the oscillator representation of the open string Hilbert space, the standard method is to introduce a chiral field on the double cover (see for example ).
For the annulus case depicted in figure 3, we impose the path integral variable $`𝒳`$ to have the twisted boundary condition,
$$𝒳(\sigma _0,\sigma _1+2\pi )=h𝒳(\sigma _0,\sigma _1)𝒳(\sigma _0+2\pi ,\sigma _1)=g𝒳(\sigma _0,\sigma _1)$$
(3.6)
$`g,h\mathrm{\Gamma }`$ should satisfy $`gh=hg.`$ We will identify it with $`𝐗`$ by (for $`0\sigma _1\pi `$)
$$𝐗_L(\sigma _0,\sigma _1)=𝒳(\sigma _0,\sigma _1),𝐗_R(\sigma _0,\sigma _1)=f_2𝒳(\sigma _0,2\pi \sigma _1).$$
(3.7)
With this assignment, $`𝐗`$ satisfies the boundary condition at $`\sigma _1=\pi `$ automatically. On the other hand, the boundary condition at $`\sigma _1=0`$ requires,
$$𝐗_R(\sigma _0,0)=f_2𝒳(\sigma _0,2\pi )=f_2h𝒳(\sigma _0,0).$$
(3.8)
Since the last formula should be same as $`f_1𝒳(\sigma _0,0)`$,
$$h=f_2^1f_1.$$
(3.9)
$`h`$ gives the actual twist in the open string sector and satisfies $`h^gghg^1=h`$.
The variant of the modular invariance in the open string case is that the annulus partition function can be alternatively expressed as the inner product between the boundary states. In the orbifold case, the boundary state depends on two elements $`(g,f)`$ in the group $`\mathrm{\Gamma }`$. We will denote it as $`|B:g,f`$.
The first element $`g`$ specifies which twisted sector the closed string variable belongs to.
$$(\overline{𝐗}(\overline{\sigma }_0,\overline{\sigma }_1+2\pi )g\overline{𝐗}(\overline{\sigma }_0,\overline{\sigma }_1))|_{\overline{\sigma }_0=0}|B:g,f=0.$$
(3.10)
(We used different notations for embedding functions $`\overline{𝐗}`$ and world sheet coordinates $`\sigma `$ in order to explicitly shows that we are considering the closed string sector.) The second element $`f`$ specifies the twist at the boundary,
$$(\overline{𝐗}_R(\overline{\sigma }_0,\overline{\sigma }_1)f\overline{𝐗}_L(\overline{\sigma }_0,\overline{\sigma }_1))|_{\overline{\sigma }_0=0}|B:g,f=0.$$
(3.11)
From the boundary conditions of the open string, we need to impose the constraints,
$$[g,f]=0,f^2=1.$$
(3.12)
The modular transformation implies the relation,
$$\text{Tr}_{_{f_1,f_2}}\left(g𝐞\left[\tau L_0^{open}\right]\right)B:g,f_2|𝐞\left[(1/2\tau )(L_0^{closed}+\overline{L}_0^{closed})\right]|B:g,f_1.$$
(3.13)
### 3.3 Möbius strip and cross-cap state
In Möbius diagram, we need to calculate the trace such as,
$$\text{Tr}_{_{f_1,f_2}}\left(\mathrm{\Omega }g𝐞\left[\tau L_0^{open}\right]\right),$$
(3.14)
where $`\mathrm{\Omega }`$ is the open string flip operator. As in the annulus case, we need to impose some constraints on $`f_1,f_2,g\mathrm{\Gamma }`$ to have a non-vanishing result.
First let us investigate the action of $`\mathrm{\Omega }g`$ on $`_{f_1,f_2}`$. It acts on $`𝐗_L,𝐗_R`$ as
$`𝐗_L`$ $``$ $`𝐘_L(\sigma _0,\sigma _1)g𝐗_R(\sigma _0,\pi \sigma _1)`$
$`𝐗_R`$ $``$ $`𝐘_R(\sigma _0,\sigma _1)g𝐗_L(\sigma _0,\pi \sigma _1).`$ (3.15)
The boundary condition for the field $`𝐘`$ becomes,
$`𝐘_R(\sigma _0,0)`$ $`=`$ $`g𝐗_L(\sigma _0,\pi )`$ (3.16)
$`=`$ $`gf_2^1𝐗_R(\sigma _0,\pi )`$
$`=`$ $`gf_2^1g^1𝐘_L(\sigma _0,0)`$
In other word, $`𝐗_{f_1,f_2}`$ implies $`\mathrm{\Omega }g𝐗_{gf_2^1g^1,gf_1^1g^1}`$. In order to have non-vanishing trace (3.14), we need to impose these two Hilbert spaces are the same,
$$f_2=gf_1^1g^1,f_1=gf_2^1g^1.$$
(3.17)
If one examines these two conditions carefully, one may notice that $`g,f_1,f_2`$ should satisfy further conditions,
$$[g^2,f_1]=[g^2,f_2]=0$$
(3.18)
To summarize, the non-vanishing Möbius strip amplitude is characterized by $`g,f_1\mathrm{\Gamma }`$ with constraints,
$$f_1^2=1,[g^2,f_1]=0.$$
(3.19)
$`f_2`$ is then determined from (3.17). The actual twist in the open string is given by,
$$h=f_2^1f_1=gf_1g^1f_1,$$
(3.20)
which satisfies the condition typical in unorientable cases, $`h^g=h^1`$.
In the standard double covering of Möbius strip (figure 4), these conditions are explained as follows. We introduce the path integral variable $`𝒳`$ satisfies the twisted boundary condition,
$`𝒳(\sigma _0+\pi ,\sigma +\pi )`$ $`=`$ $`\alpha 𝒳(\sigma _0,\sigma _1)`$
$`𝒳(\sigma _0+2\pi ,\sigma )`$ $`=`$ $`\beta 𝒳(\sigma _0,\sigma _1)\alpha ,\beta \mathrm{\Gamma }.`$ (3.21)
From this field, we would like to construct the left and right movers of the open string which satisfy the boundary condition,
$`𝐗_R(\sigma _0,0)`$ $`=`$ $`f_1𝐗_L(\sigma _0,0)`$
$`𝐗_R(\sigma _0+\pi ,\pi /2)`$ $`=`$ $`g𝐗_L(\sigma _0,\pi /2),`$
$`𝐗_L(\sigma _0+\pi ,\pi /2)`$ $`=`$ $`g𝐗_R(\sigma _0,\pi /2).`$ (3.22)
The boundary conditions in the second and the third lines are (twisted) cross-cap type conditions. It is consistent with (3.3) only when $`\beta =g^2`$.
We use the following identification between $`𝐗`$ and $`𝒳`$,
$`𝐗_L(\sigma _0,\sigma _1)`$ $`=`$ $`𝒳(\sigma _0,\sigma _1)`$
$`𝐗_R(\sigma _0,\sigma _1)`$ $`=`$ $`f_3𝒳(\sigma _0+\pi ,\pi \sigma _1).`$ (3.23)
Boundary condition at $`\sigma _1=\pi /2`$ implies,
$`𝐗_R(\sigma _0+\pi ,\pi /2)`$ $`=`$ $`f_3𝒳(\sigma _0+2\pi ,\pi /2)`$
$`=`$ $`f_3g^2𝒳(\sigma _0,\pi /2)`$
$`g𝐗_L(\sigma _0,\pi /2)`$ $`=`$ $`g𝒳(\sigma _0,\pi /2).`$ (3.24)
Namely,
$$f_3=g^1.$$
(3.25)
Similarly boundary condition at $`\sigma _1=0`$ is consistent with (3.3) only when,
$$\alpha =gf_1$$
(3.26)
Periodicity and the twisted boundary condition along $`\sigma _1=0`$ implies,
$$[f_1,\beta ]=[f_1,g^2]=0.$$
(3.27)
This calculation can be summarized by the introduction of the cross-cap state $`|C:g,f`$ which satisfies,
$`(\overline{𝐗}_L(0,\overline{\sigma }_1+\pi )f\stackrel{~}{𝐗}_R(0,\overline{\sigma }_1))|C:g,f`$ $`=`$ $`0`$
$`(\overline{𝐗}_R(0,\overline{\sigma }_1+\pi )f\stackrel{~}{𝐗}_L(0,\overline{\sigma }_1))|C:g,f`$ $`=`$ $`0`$ (3.28)
The twisted boundary condition of the closed string is specified by $`g`$. However, it is clear that these conditions automatically implies that $`\overline{𝐗}`$ belongs to the twisted sector defined by $`f^2`$,
$$\overline{𝐗}(\overline{\sigma }_0,\overline{\sigma }_1+2\pi )=f^2\overline{𝐗}(\overline{\sigma }_0,\overline{\sigma }_1).$$
(3.29)
Namely we always have $`g=f^2`$. The modular transformation of Möbius strip can be now written as
$$\text{Tr}_{_{f_1,f_1^g}}\left(g𝐞\left[\tau L_0^{open}\right]\right)=B:g^2,f_1|𝐞\left[(1/8\tau )(L_0^{closed}+\overline{L}_0^{closed})\right]|C:g^2,g.$$
(3.30)
### 3.4 Klein bottle
In the Klein bottle amplitude, we need to evaluate
$$\text{Tr}__h\left(\mathrm{\Omega }^{closed}g𝐞\left[\tau (L_0^{closed}+\overline{L}_0^{closed})\right]\right)$$
(3.31)
As in the Möbius strip case,
$$𝐗_h\mathrm{\Omega }^{closed}g𝐗_{gh^1g^1}.$$
(3.32)
Therefore, to have non-vanishing trace, we need impose constraint on $`g,h`$,
$$h=gh^1g^1\text{or}hg=gh^1.$$
(3.33)
In the double covering (figure 5), we define a chiral field $`𝒳`$ which has twisted boundary condition,
$$𝒳(\sigma _0,\sigma _1+2\pi )=h𝒳(\sigma _0,\sigma _1),𝒳(\sigma _0+2\pi ,\sigma _1)=\beta 𝒳(\sigma _0,\sigma _1).$$
(3.34)
$`𝐗`$ satisfies cross-cap type boundary condition at $`\sigma _1=0,\pi `$,
$$𝐗_{R,L}(\sigma _0+\pi ,0)=f_1𝐗_{L,R}(\sigma _0,0)𝐗_{R,L}(\sigma _0+\pi ,\pi )=f_2𝐗_{L,R}(\sigma _0,\pi )$$
(3.35)
We identify
$$𝐗_L(\sigma _0,\sigma _1)=𝒳(\sigma _0,\sigma _1)𝐗_R(\sigma _0,\sigma _1)=f_2𝒳(\sigma _0\pi ,2\pi \sigma _1)$$
(3.36)
It automatically satisfies boundary condition at $`\sigma _1=\pi `$. Boundary condition at $`\sigma _1=0`$ is satisfied if
$$h=f_2^1f_1.$$
(3.37)
Twist in time direction requires,
$$\beta =f_1^2=f_2^2.$$
(3.38)
If we identify $`g=f_1`$, two conditions (3.33) and (3.38) are equivalent. The modular invariance in this case can be written as,
$$\text{Tr}__h\left(g\mathrm{\Omega }^{closed}𝐞\left[\tau (L_0+\overline{L}_0)\right]\right)=C:f_2^2,f_2|𝐞\left[(1/4\tau )(L_0+\overline{L}_0)\right]|C:f_1^2,f_1$$
(3.39)
with $`g=f_1`$ and $`h=f_2^1f_1`$.
## 4 Classification of irreducible boundary conditions
In this section, we explicitly solve the constraints of the twists in the previous section for three diagrams. We classify all the possible irreducible solutions together with the free parameters. Our result in this section is summarized in the following theorem.
Theorem 2
(i) Klein bottle: irreducible solutions for (3.33) are given by,
$$h=diag(\stackrel{m}{\stackrel{}{𝒯_n,\mathrm{},𝒯_n}}),g=diag(𝒯_n^{p_0}𝒮_n,\mathrm{},𝒯_n^{p_{m1}}𝒮_n)\overline{𝒯}_m,$$
(4.1)
where $`𝒮_n`$ is the inversion permutation of $`n`$ elements $`(n1,n2,n3,\mathrm{},0)`$. $`p_i`$ ($`i0,\mathrm{},m1`$ mod $`m`$) takes their values in $`0,1,\mathrm{},n1`$ (mod $`n`$). We will refer this solution as (K). (ii) Annulus: there are three types of the irreducible solutions for (3.5,3.9) together with $`f_i^2=1`$.
* ($`I_A`$):
$`h`$ $`=`$ $`diag(\stackrel{m}{\stackrel{}{𝒯_n,\mathrm{},𝒯_n}})`$
$`g`$ $`=`$ $`diag(𝒯_n^{p_0},\mathrm{},𝒯_n^{p_{m1}})\overline{𝒯}_m,`$
$`f_1`$ $`=`$ $`diag(𝒯_n^{q_0}𝒮_n,\mathrm{},𝒯_n^{q_{m1}}𝒮_n),`$
$`f_2`$ $`=`$ $`diag(𝒯_n^{q_0+1}𝒮_n,\mathrm{},𝒯_n^{q_{m1}+1}𝒮_n).`$ (4.2)
$`p_i,q_i`$ ($`i0,\mathrm{},m1`$ mod $`m`$) takes their values in $`0,1,\mathrm{},n1`$ (mod $`n`$). They should satisfy the constraint,
$$q_{\mathrm{}}q_{\mathrm{}+1}2p_{\mathrm{}},(\mathrm{}=0,\mathrm{},n1\text{ mod}n)$$
(4.3)
* ($`II_A`$): for even $`m`$, with the same $`h,g`$ as in (4), together with
$`f_1`$ $`=`$ $`diag(𝒯_n^{q_0}𝒮_n,\mathrm{},𝒯_n^{q_{m1}}𝒮_n)(\overline{𝒯}_{(m)})^{m/2},`$
$`f_2`$ $`=`$ $`diag(𝒯_n^{q_0+1}𝒮_n,\mathrm{},𝒯_n^{q_{m1}+1}𝒮_n)(\overline{𝒯}_{(m)})^{m/2},`$ (4.4)
with $`q_k=q_{k+m/2}`$. The constraint for $`p,q`$ are given by
$$q_{\mathrm{}}q_{\mathrm{}+1}p_{\mathrm{}}+p_{\mathrm{}+m/2}(\text{mod}n)$$
(4.5)
$`\mathrm{}`$ in this equation is defined in modulo $`m`$.
* ($`\stackrel{~}{II}_A`$): for even $`m`$, with the same $`h,f_1,f_2`$ as in ($`II_A`$), together with
$$g=diag(𝒯_n^{p_0},\mathrm{},𝒯_n^{p_{m1}})diag(\overline{𝒯}_{m/2},\overline{𝒯}_{m/2}).$$
(4.6)
The constraint for $`p,q`$ are given by
$$q_{\mathrm{}}q_{\mathrm{}+1}p_{\mathrm{}}+p_{\mathrm{}+m/2}(\text{mod}n)$$
(4.7)
The index $`\mathrm{}`$ in this equation is defined in modulo $`m/2`$.
(iii) Möbius strip: there are three types of the irreducible solutions for (3.19,3.20).
* ($`I_M`$):
$`h`$ $`=`$ $`diag(\stackrel{m}{\stackrel{}{𝒯_n,\mathrm{},𝒯_n}})`$
$`g`$ $`=`$ $`diag(𝒯_n^{p_0}𝒮_n,\mathrm{},𝒯_n^{p_{m1}}𝒮_n)\overline{𝒯}_m,`$
$`f_1`$ $`=`$ $`diag(𝒯_n^{q_0}𝒮_n,\mathrm{},𝒯_n^{q_{m1}}𝒮_n).`$ (4.8)
The constraint for $`p,q`$ is
$$q_{\mathrm{}}+q_{\mathrm{}+1}2p_{\mathrm{}}1,(\mathrm{}=0,\mathrm{},n1\text{ mod }n)$$
(4.9)
* ($`II_M`$): for even $`m`$, with the same $`h,g`$ as in (4), together with
$$f_1=diag(𝒯_n^{q_0}𝒮_n,\mathrm{},𝒯_n^{q_{m1}}𝒮_n)(\overline{𝒯}_{(m)})^{m/2}.$$
(4.10)
The constraint for $`p,q`$ are given by
$$q_{\mathrm{}}+q_{\mathrm{}+1}p_{\mathrm{}}+p_{\mathrm{}+m/2}1(\text{mod}n)$$
(4.11)
$`\mathrm{}`$ in this equation is defined in modulo $`m`$.
* ($`\stackrel{~}{II}_M`$): for even $`m`$, with the same $`h,f_1`$ as in ($`II_M`$), together with
$$g=diag(𝒯_n^{p_0}𝒮_n,\mathrm{},𝒯_n^{p_{m1}}𝒮_n)diag(\overline{𝒯}_{m/2},\overline{𝒯}_{m/2}).$$
(4.12)
The constraint for $`p,q`$ are given by
$$q_{\mathrm{}}+q_{\mathrm{}+1}p_{\mathrm{}}+p_{\mathrm{}+m/2}1(\text{mod}n)$$
(4.13)
The index $`\mathrm{}`$ in this equation is defined in modulo $`m/2`$.
While $`𝒯_n^p`$ represents the cyclic rotation of string bits in length $`n`$ long string, $`𝒮_n`$ flips the orientation of the long string. Four sectors, $`(II_A)`$, $`(\stackrel{~}{II}_A)`$, $`(II_M)`$, $`(\stackrel{~}{II}_M)`$, can be interpreted as giving the closed string sectors with/without orientation flip which appears in open string sector.
The rest of this section is devoted to straightforward but rather lengthy proof of this theorem. Before we embark on the calculation of indivisual cases, we first mention a simple lemma
Lemma 1: In the irreducible sets, $`h`$ should always be the direct product of the cyclic permutations of the same length $`n`$. Proof: We already mentioned that for each diagram, $`h`$ and $`g`$ satisfy,
$`h^g=`$ $`h,`$ Annulus
$`h^g=`$ $`h^1,`$ $`\text{Klein bottle, Möbius strip},`$ (4.14)
where $`h^gghg^1`$. Let us assume that the conjugacy class of $`h`$ is given by the partition $`(n)(m)`$ with $`nm`$. LHS of (4) belongs to the same conjugacy class with each element permuted by $`g`$. Since $`nm`$, $`g`$ can not mix elements in $`(n)`$ and $`(m)`$.
For open string sectors (Annulus and Möbius), we also determine $`f_i`$ ($`i=1,2`$) with $`f_i^2=1`$ and $`h=f_2^1f_1`$. Because of $`f_i^2=1`$, one may easily derive that $`h^{f_1}=h^1`$. By repeating our argument on $`g`$, $`f_1`$ (and also $`f_2`$) can not mix the elements in $`(n)`$ and $`(m)`$.
It proved that when $`h=(n)(m)`$ with $`nm`$, there are no the irreducible sets. QED
### 4.1 Klein bottle
From lemma 1, we may restrict $`h,g`$ to the following form,
$$h=diag(\stackrel{m}{\stackrel{}{𝒯_n,\mathrm{},𝒯_n}}),g=diag(\alpha _0,\mathrm{},\alpha _{m1})\overline{𝒯}_m,$$
(4.15)
where $`\alpha _{\mathrm{}}`$ is any permutations of $`n`$ elements. $`h^g=h^1`$ implies, $`\alpha _{\mathrm{}}𝒯_n\alpha _{\mathrm{}}^1=𝒯_n^1`$. To derive conditions in (K), it is enough to prove the following lemma, Lemma 2: The general solution to
$$g𝒯_ng^1=𝒯_n^1,$$
(4.16)
is given by $`g=𝒯_n^p𝒮_n`$ for $`p=0,1,\mathrm{},n`$. Proof: Assume that $`g`$ maps $`(0,1,2,\mathrm{},n1)`$ to $`(g_0,g_1,\mathrm{},g_{n1})`$. $`g𝒯_n=𝒯_n^1g`$ implies
$$g(i+1)g(i)+10\text{mod}n.$$
(4.17)
General solution to this difference equation is clearly (4.16). QED.
### 4.2 Annulus
Because $`[h,g]=0`$, we may write general irreducible solutions in the following form,
$`h`$ $`=`$ $`diag(\stackrel{m}{\stackrel{}{𝒯_n,\mathrm{},𝒯_n}})`$
$`g`$ $`=`$ $`diag(𝒯_n^{p_0},\mathrm{},𝒯_n^{p_{m1}})G,`$
$`f_i`$ $`=`$ $`diag(\alpha _1^{(i)},\mathrm{},\alpha _{m1}^{(i)})F_i.`$ (4.18)
where $`G`$ and $`F_i`$ belong to $`S_m`$. The constraints $`[g,f_i]=0`$, $`h=f_2^1f_1`$ and $`f_i^2=1`$ implies,
$`[G,F_i]`$ $`=`$ $`0`$
$`F_1^2=F_2^2=F_1F_2`$ $`=`$ $`1.`$ (4.19)
The equation in the second line implies $`F_1=F_2`$. We can use Lemma 1 to show that $`G`$ is written as a direct product of the cyclic permutations of the same length $`(\overline{𝒯}_{m/s})^s`$ if we use the irreducibility. From the first equation, $`s`$ must be either 1 or 2. If $`s=1`$, $`F_i`$ can be either 1 ($`I_A`$),
$$G=\overline{𝒯}_mF=1_m,$$
(4.20)
or $`\overline{𝒯}_{m/2}`$ if $`m`$ is even ($`II_A`$)
$$G=\overline{𝒯}_mF=G^{m/2}=\left(\begin{array}{cc}0& 1_{m/2}\\ 1_{m/2}& 0\end{array}\right),$$
(4.21)
If $`s=2`$, $`G`$ and $`F_i`$ must be the following form ($`\stackrel{~}{II}_A`$),
$$G=\left(\begin{array}{cc}\overline{𝒯}_{m/2}& 0\\ 0& \overline{𝒯}_{m/2}\end{array}\right),F=\left(\begin{array}{cc}0& 1_{m/2}\\ 1_{m/2}& 0\end{array}\right),$$
(4.22)
with even $`m`$.
#### 4.2.1 $`I_A`$
$`f_i^2=1`$ and $`f_2^1f_1=h`$ implies
$$(\alpha _{\mathrm{}}^{(2)})^1\alpha _{\mathrm{}}^{(1)}=𝒯_n,(\alpha _{\mathrm{}}^{(i)})^2=1.$$
(4.23)
(4) is easily derived by using the following lemma. (4.3) comes from the constraint $`[f_i,g]=0`$.
Lemma 3: For $`f_1,f_2S_n`$, the general solution to
$$f_1^2=f_2^2=1,f_2^1f_1=𝒯_n$$
(4.24)
is given by
$$f_1=𝒯_n^p𝒮_n,f_2=𝒯_n^{p+1}𝒮_n,p=0,\mathrm{},n1\text{mod}n.$$
(4.25)
Proof: If we write $`f_2^1f_1=h`$ and $`f_1=g`$,
$$h^g=f_1f_2=h^1=𝒯_n^1.$$
(4.26)
This is exactly the same condition as Lemma 2. It permits us to write $`f_i`$ in the form (4.25). By using the relation,
$$𝒯_n^p𝒮_n=𝒮_n𝒯_n^p,$$
(4.27)
it is straightforward to prove that $`f_i`$s thus defined satisfy $`f_i^2=1`$. QED.
#### 4.2.2 $`II_A`$, $`\stackrel{~}{II}_A`$
We write $`m/2k`$. The constraints $`f_i^2=1`$ are written in components, $`\alpha _{\mathrm{}}^{(i)}\alpha _{\mathrm{}+k}^{(i)}=1`$, namely $`\alpha _{\mathrm{}+k}^{(i)}=(\alpha ^{(i)})_{\mathrm{}}^1`$. By plugging it into $`f_2^1f_1=h`$, we get
$$\alpha _{\mathrm{}}^{(2)}(\alpha _{\mathrm{}}^{(1)})^1=(\alpha _{\mathrm{}}^{(2)})^1\alpha _{\mathrm{}}^{(1)}=𝒯_n.$$
From this equation, it is not difficult to prove $`\alpha _{\mathrm{}}^{(i)}𝒯_n\alpha _{\mathrm{}}^{(i)1}=𝒯_n^1`$. By using lemma 2, we conclude that $`f_i`$ must be of the form, (4). Finally constraints (4.5,4.7) are obtained by imposing $`[g,f_i]=0`$.
### 4.3 Möbius
With the help of Klein bottle calculation, one may seek the general irreducible solution in the following form,
$`h`$ $`=`$ $`diag(\stackrel{m}{\stackrel{}{𝒯_n,\mathrm{},𝒯_n}})`$
$`g`$ $`=`$ $`diag(𝒯_n^{p_0}𝒮_n,\mathrm{},𝒯_n^{p_{m1}}𝒮_n)G,`$
$`f_1`$ $`=`$ $`diag(\alpha _0,\mathrm{},\alpha _{m1})F_1,`$ (4.28)
with $`G,F_1S_m`$. In this form $`h^g=h^1`$ is automatically satisfied. (3.19,3.20) lead to
$$[F_1,G]=0,F_1^2=1.$$
(4.29)
This is exactly the same constraint in the annulus (4.2). We may use the same solutions, (4.20,4.21,4.22) and call the solutions associated with each of them as $`I_M`$, $`II_M`$, $`\stackrel{~}{II}_M`$.
#### 4.3.1 $`I_M`$
$`f_1^2=1`$ implies $`\alpha _{\mathrm{}}^2=1`$. Constraint (3.20) gives
$$𝒯_n^p_{\mathrm{}}𝒮_n\alpha _{\mathrm{}}^1𝒯_n^p_{\mathrm{}}𝒮_n\alpha _{\mathrm{}+1}=𝒯_n.$$
(4.30)
Since $`\alpha _{\mathrm{}}^2=1`$ and $`(𝒯_n^p_{\mathrm{}}𝒮_n\alpha _{\mathrm{}}^1𝒯_n^p_{\mathrm{}}𝒮_n)^2=1`$, lemma 3 permits us to write
$$\alpha _{\mathrm{}}=𝒯_n^q_{\mathrm{}}𝒮_n,$$
(4.31)
which proved (4). The constraint (4.3) comes from putting this value into (4.30) again.
#### 4.3.2 $`II_M`$ and $`\stackrel{~}{II}_M`$
In these cases, $`f_1`$ has the following form ($`m=2k`$),
$$f_1=diag(\alpha _0,\mathrm{},\alpha _{m1})(\overline{𝒯}_{2k})^k.$$
(4.32)
$`f_1^2=1`$ gives $`\alpha _{\mathrm{}+k}=\alpha _{\mathrm{}}^1`$. We note that $`f_2=(f_1)^g`$ have the same form. As in the $`II_A`$ case, one can show that $`\alpha _{\mathrm{}}`$ have the following form,
$$\alpha _{\mathrm{}}=𝒯_n^q_{\mathrm{}}𝒮_n.$$
(4.33)
(3.20) gives the same type of constraint in both $`II_M`$ and $`\stackrel{~}{II}_M`$<sup>4</sup><sup>4</sup>4 The difference between those two cases is whether the variable $`\mathrm{}`$ is counted as mod $`2k`$ ($`II_M`$) or as mod $`k`$ ($`\stackrel{~}{II}_M`$). ,
$$(𝒯_n^p_{\mathrm{}}𝒮_n)\alpha _{\mathrm{}}(𝒯_n^{p_{k+\mathrm{}}}𝒮_n)\alpha _{\mathrm{}+1}^1=𝒯_n,$$
(4.34)
which gives (4.11,4.13).
The proof of theorem 2 is completed.
## 5 Partition functions
In this section, we explicitly calculate the partition functions for the irreducible boundary conditions discussed in previous section. One interesting feature is that the sectors for the long strings are in general different from those of the short strings. Actually this is clear since we already mentioned there are the long closed string sectors in the annulus or Möbius strip amplitude in the short string.
Since the correspondence looks rather complicated, we summarize our result in the following table.
| Short string sector | $`n`$ | $`m`$ | Long string sector | Partition function |
| --- | --- | --- | --- | --- |
| KB | $``$ | odd | Klein Bottle | $`𝒵^{KB}(\tau _{n,m,0})`$ |
| | $``$ | even | Torus | $`𝒵^T(\tau _{n,m,p},\overline{\tau }_{n,m,p})`$ |
| Annulus: $`I_A`$ | odd | $``$ | Annulus | $`𝒵^A(\tau _{n,m,0})`$ |
| | even | $``$ | Annulus+Möbius | $`𝒵^A(\tau _{n,m,0})+𝒵^M(\tau _{n,m,0})`$ |
| Annulus: $`II_A`$ | $``$ | $`2\times `$ | Klein Bottle | $`𝒵^{KB}(\tau _{n,m/2,0})`$ |
| Annulus: $`\stackrel{~}{II}_A`$ | $``$ | $`2\times `$ | Torus | $`𝒵^T(\tau _{n,m/2,p},\overline{\tau }_{n,m/2,p})`$ |
| Möbius: $`I_M`$ | odd | odd | Möbius | $`𝒵^M(\tau _{n,m,0})`$ |
| | odd | even | Annulus | $`𝒵^A(\tau _{n,m,0})`$ |
| | even | even | Annulus+Möbius | $`𝒵^A(\tau _{n,m,0})+𝒵^M(\tau _{n,m,0})`$ |
| | even | odd | | 0 |
| Möbius: $`II_M`$ | $``$ | $`2\times `$odd | Torus | $`𝒵^T(\tau _{n,m/2,p^{}},\overline{\tau }_{n,m/2,p^{}})`$ |
| | $``$ | $`2\times `$even | Klein Bottle | $`𝒵^{KB}(\tau _{n,m/2,0})`$ |
| Möbius: $`\stackrel{~}{II}_M`$ | $``$ | $`2\times `$even | Torus | $`𝒵^T(\tau _{n,m/2,p},\overline{\tau }_{n,m/2,p})`$ |
| | $``$ | $`2\times `$odd | Klein Bottle | $`𝒵^{KB}(\tau _{n,m/2,0})`$ |
For Klein Bottle/Annuls/Möbius strip cases, the moduli parameter $`\tau `$ is pure imaginary. $`\tau _{n,m,p}`$ is the long string moduli (2.27). $`p`$ is an integer from $`0`$ to $`n1`$ and $`p^{}`$ is the half-odd integer from $`0`$ to $`n`$.
In this table we used the partition functions for a single string in each sector. We will first prove this table by employing the explicit operator formalism for the simplest target space $`𝐑^1`$. In this case, they have the following standard form,
$`𝒵^T(\tau ,\overline{\tau })`$ $`=`$ $`{\displaystyle \frac{𝐞[(\tau \overline{\tau }))/24]}{\sqrt{\text{Im}\tau }}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left((1𝐞\left[n\tau \right])(1𝐞\left[n\overline{\tau }\right])\right)^1,`$ (5.1)
$`𝒵^{KB}(\tau )`$ $`=`$ $`{\displaystyle \frac{𝐞\left[\tau /12\right]}{\sqrt{\text{Im}\tau }}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(1𝐞\left[2n\tau \right]\right)^1,`$ (5.2)
$`𝒵^A(\tau )`$ $`=`$ $`{\displaystyle \frac{𝐞\left[\tau /24\right]}{\sqrt{2\text{Im}\tau }}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(1𝐞\left[n\tau \right]\right)^1,`$ (5.3)
$`𝒵^M(\tau )`$ $`=`$ $`{\displaystyle \frac{𝐞\left[\tau /24\right]}{\sqrt{2\text{Im}\tau }}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(1(1)^n𝐞\left[n\tau \right]\right)^1.`$ (5.4)
In the first subsection, we give the examination of this table by using the explicit operator formalism. In the most cases, we omit the analysis of the momentum integration since they are the same as our discussion in section 2.
In the second subsection, we present a different proof based on multiple cover of the world sheets. This is similar to our discussion in the section 2.2.2 and can be applied to the arbitrary target space. Although it may be less rigorous compared to the analysis by the operator formalism, it is much better to explain the topological nature of this table.
### 5.1 Analysis by Operator formalism
#### 5.1.1 Klein Bottle
As in our calculation in the torus amplitude, we make the discrete Fourier transformation (2.16) for the component fields. The action of $`h`$ was determined in (2.2.1) and it is possible to use the same mode expansion (2.19). One novelty is to determine the action of $`𝒮_n`$,
$`(𝒮_n\stackrel{~}{𝐗})^{a,J}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{n}}}{\displaystyle \underset{I=0}{\overset{n1}{}}}𝐞\left[{\displaystyle \frac{aI}{n}}\right]𝐗^{nI1,J}`$ (5.5)
$`=`$ $`𝐞\left[{\displaystyle \frac{a}{n}}\right]\stackrel{~}{𝐗}^{a,J}.`$
The action of $`g`$ in (4.1) is evaluated as,
$$(g\stackrel{~}{𝐗})^{a,J}=𝐞\left[\frac{ap_J}{n}\right]\stackrel{~}{𝐗}^{a,J+1}.$$
(5.6)
Since the orientation flip interchange the left and the right movers, $`(\mathrm{\Omega }g)`$ acts on the oscillators as
$`(\mathrm{\Omega }g)\alpha _{ra/n}^{a,J}`$ $`=`$ $`𝐞\left[{\displaystyle \frac{ap_J}{n}}\right]\stackrel{~}{\alpha }_{ra/n}^{a,J+1}`$
$`(\mathrm{\Omega }g)\stackrel{~}{\alpha }_{ra/n}^{a,J}`$ $`=`$ $`𝐞\left[{\displaystyle \frac{ap_J}{n}}\right]\alpha _{ra/n}^{a,J+1}.`$ (5.7)
We need to find a combination of the left and right movers which is invariant up to scalar multiplication under $`\mathrm{\Omega }g`$. To this end, we define a linear combination
$$\beta _{ra/n}^a=\underset{J=0}{\overset{m1}{}}C_J\alpha _{ra/n}^{a,J},$$
(5.8)
and define $`\stackrel{~}{\beta }_{ra/n}^a(\mathrm{\Omega }g)\beta _{ra/n}^a`$. If one can find appropriate coefficients $`C_J`$ such that $`\beta `$ satisfies
$$\mathrm{\Omega }g\stackrel{~}{\beta }_{ra/n}^a=\lambda \beta _{ra/n}^a,$$
(5.9)
$`\beta _{ra/n}^a\stackrel{~}{\beta }_{ra/n}^a`$ becomes diagonal under the action of $`\mathrm{\Omega }g`$,
$$\mathrm{\Omega }g\left(\beta _{ra/n}^a\stackrel{~}{\beta }_{ra/n}^a\right)=\lambda \left(\beta _{ra/n}^a\stackrel{~}{\beta }_{ra/n}^a\right).$$
(5.10)
In terms of $`C_J`$, (5.9) becomes
$$\lambda C_{J+2}=C_J𝐞\left[\frac{a}{n}(p_Jp_{J+1})\right].$$
(5.11)
Since $`J`$ takes its value in $`0,\mathrm{},m1`$ (modulo $`m`$), the solution becomes essentially different depending on whether $`m`$ is even or odd. odd $`m`$:
In this case, the recursion relation gives
$`C_{2m}`$ $`=`$ $`{\displaystyle \frac{1}{\lambda }}𝐞\left[{\displaystyle \frac{a}{n}}(p_{2m2}p_{2m1})\right]C_{2(m1)}`$ (5.12)
$`=`$ $`\mathrm{}`$
$`=`$ $`{\displaystyle \frac{1}{\lambda ^m}}𝐞\left[{\displaystyle \frac{a}{n}}\left({\displaystyle \underset{\mathrm{}=0}{\overset{m1}{}}}p_{2m2\mathrm{}}{\displaystyle \underset{\mathrm{}=0}{\overset{m1}{}}}p_{2m12\mathrm{}}\right)\right]C_0`$
$`=`$ $`{\displaystyle \frac{1}{\lambda ^m}}C_0.`$
$`C_{2m}=C_0`$ implies $`\lambda =𝐞\left[\frac{b}{m}\right]`$ with $`b=0,1,\mathrm{}m1`$.
Along the same line of calculation (2.2.1), we obtain the oscillator contribution to the partition function of Klein bottle,
$`\text{Tr}__h(\mathrm{\Omega }g𝐞\left[\tau (L_0+\overline{L}_0)\right])`$ $`=`$ $`{\displaystyle \underset{s=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{a=0}{\overset{n1}{}}}{\displaystyle \underset{b=0}{\overset{m1}{}}}{\displaystyle \frac{1}{1𝐞\left[\frac{b}{m}+2\tau \left(s\frac{a}{n}\right)\right]}}`$ (5.13)
$`=`$ $`{\displaystyle \underset{s=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{a=0}{\overset{n1}{}}}{\displaystyle \frac{1}{1𝐞\left[2m\tau \left(s\frac{a}{n}\right)\right]}}`$
$`=`$ $`{\displaystyle \underset{s=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{1𝐞\left[\frac{2m\tau }{n}\right]}}.`$
The calculation of the zero mode contribution is the same as the torus. It proved the table.
even $`m`$:
In this case, the recursion relation is split into two sequences $`\left\{C_{\mathrm{}}\right\}_{\mathrm{}_{even}}`$ and $`\left\{C_{\mathrm{}}\right\}_{\mathrm{}_{odd}}`$. The first one can be solved as,
$$C_m=\frac{1}{\lambda ^{m/2}}𝐞\left[\frac{ap}{n}\right]C_0.$$
(5.14)
with
$$p=\underset{\mathrm{}=0}{\overset{m/21}{}}p_{m2\mathrm{}}\underset{\mathrm{}=0}{\overset{m/21}{}}p_{m12\mathrm{}}.$$
(5.15)
It gives quantized eigenvalues for $`\lambda `$,
$$\lambda _{even}^{a,b}=𝐞\left[\frac{2b}{m}+\frac{2ap}{nm}\right],b=0,\mathrm{},\frac{m}{2}1.$$
(5.16)
Similarly, the recursion relation for $`\left\{C_{\mathrm{}}\right\}_{\mathrm{}_{odd}}`$, gives
$$\lambda _{odd}^{a,b}=𝐞\left[\frac{2b}{m}\frac{2ap}{nm}\right],b=0,\mathrm{},\frac{m}{2}1.$$
(5.17)
The calculation of the partition function is now completely parallel to (2.2.1),
$`\text{Tr}__h(\mathrm{\Omega }g𝐞\left[\tau (L_0+\overline{L}_0)\right])`$ $`=`$ $`{\displaystyle \underset{s=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{a=0}{\overset{n1}{}}}{\displaystyle \underset{b=0}{\overset{m/21}{}}}\left(1𝐞\left[{\displaystyle \frac{2b}{m}}{\displaystyle \frac{2ap}{nm}}+2\tau \left(s{\displaystyle \frac{a}{n}}\right)\right]\right)^1`$ (5.18)
$`\left(1𝐞\left[{\displaystyle \frac{2b}{m}}+{\displaystyle \frac{2ap}{nm}}+2\tau \left(s{\displaystyle \frac{a}{n}}\right)\right]\right)^1`$
$`=`$ $`{\displaystyle \underset{s=1}{\overset{\mathrm{}}{}}}\left(1𝐞\left[{\displaystyle \frac{m\tau p}{n}}s\right]\right)^1\left(1𝐞\left[{\displaystyle \frac{m\tau +p}{n}}s\right]\right)^1`$
Together with the zero mode contribution, this is exactly the same as (2.29) except for that the short string moduli $`\tau `$ is pure imaginary. Interestingly the moduli for the long string $`\tau _{n,m,p}`$ has the real part as $`p/n`$.
#### 5.1.2 Annulus
##### $`I_A`$
From the data (4), the mode expansion that satisfies the boundary conditions (3.2) at the two boundaries is given by,
$`\stackrel{~}{𝐗}_L^{a,J}`$ $`=`$ $`i{\displaystyle \underset{s𝐙}{}}{\displaystyle \frac{1}{sa/n}}\alpha _{sa/n}^{a,J}𝐞\left[\sigma _1(s+a/n)\right]`$
$`\stackrel{~}{𝐗}_R^{a,J}`$ $`=`$ $`i𝐞\left[{\displaystyle \frac{aq_J}{n}}\right]{\displaystyle \underset{s𝐙}{}}{\displaystyle \frac{1}{s+a/n}}\alpha _{s+a/n}^{a,J}𝐞\left[\sigma _1(s+a/n)\right].`$ (5.19)
The action of $`g`$ on the other hand is given by,
$$\left(g\stackrel{~}{𝐗}_L\right)^{a,J}=𝐞\left[\frac{ap_J}{n}\right]\stackrel{~}{𝐗}_L^{a,J+1}=i𝐞\left[\frac{ap_J}{n}\right]\underset{s𝐙}{}\frac{1}{sa/n}\alpha _{sa/n}^{a,J+1}𝐞\left[\sigma _1(s+a/n)\right].$$
(5.20)
Comparing these two equations, one gets the action of $`g`$ on the oscillators,
$$(g\alpha _{sa/n})^{a,J}=𝐞\left[\frac{ap_J}{n}\right]\alpha _{sa/n}^{a,J+1}.$$
(5.21)
Since this is the same expression that we met in torus amplitude, one obtain the partition function immediately.
Unlike the torus case, $`p=_Jp_J`$ cannot take arbitrary integer. By summing over $`J`$ in (4.3), one obtains $`2p0`$ modulo $`n`$. When $`n`$ is odd, $`p=0`$ is the only solution. Putting $`p=0`$, one obtains the standard annulus contribution of one long open string,
$$\text{Tr}_{_{f_1,f_2}}\left(g𝐞\left[\tau L_0\right]\right)=\underset{r=0}{\overset{\mathrm{}}{}}\frac{1}{1𝐞\left[\frac{m\tau }{n}r\right]}.$$
(5.22)
On the other hand, when $`n`$ is even, we have two solutions, $`p=0,\frac{n}{2}`$. The first one is the annulus. The second solution gives,
$$\underset{r=1}{\overset{\mathrm{}}{}}\frac{1}{1+𝐞\left[\frac{m\tau }{n}r+\frac{r}{2}\right]}=\underset{r=1}{\overset{\mathrm{}}{}}\frac{1}{1+(1)^r𝐞\left[\frac{m\tau }{n}r\right]}.$$
(5.23)
This is the standard Möbius strip amplitude for a long string.
##### $`II_A`$ and $`\stackrel{~}{II}_A`$
In these cases we write $`m=2k`$ since $`m`$ must be even. The mode expansion and $`g`$ action is almost the same as $`I_A`$ case,
$`\stackrel{~}{𝐗}_L^{a,J}`$ $`=`$ $`i{\displaystyle \underset{s𝐙}{}}{\displaystyle \frac{1}{sa/n}}\alpha _{sa/n}^{a,J}𝐞\left[\sigma _1(s+a/n)\right]`$
$`\stackrel{~}{𝐗}_R^{a,J}`$ $`=`$ $`i𝐞\left[{\displaystyle \frac{aq_J}{n}}\right]{\displaystyle \underset{s𝐙}{}}{\displaystyle \frac{1}{s+a/n}}\alpha _{s+a/n}^{a,J+k}𝐞\left[i\sigma _1(s+a/n)\right].`$ (5.24)
The action of $`g`$ on the other hand is the same as (5.21). The difference between $`II_A`$ and $`\stackrel{~}{II}_A`$ is the definition of $`J`$. In $`II_A`$ (resp. $`\stackrel{~}{II}_A`$) case it is defined modulo $`2k=m`$ (resp. $`k`$).
In $`II_A`$ case, the constraint (4.5) gives $`_{J=0}^{2k1}p_J=0`$. The partition function becomes,
$$\underset{r=1}{\overset{\mathrm{}}{}}\frac{1}{1𝐞\left[\frac{2k\tau }{n}r\right]},$$
(5.25)
which can be identified with Klein bottle amplitude.
In $`\stackrel{~}{II}_A`$ case, $`p=_{J=0}^{k1}p_J`$ can take any value in $`0,1,\mathrm{},n1`$. The partition function therefore becomes,
$$\underset{r=1}{\overset{\mathrm{}}{}}\frac{1}{(1𝐞\left[\frac{k\tau +p}{n}r\right])(1𝐞\left[\frac{k\tau p}{n}r\right])},$$
(5.26)
which gives the torus partition function.
#### 5.1.3 Möbius strip
##### $`I_M`$
The mode expansion is given by (5.1.2). The action of $`g`$ together with the orientation flip is given by,
$`(\mathrm{\Omega }g\stackrel{~}{𝐗})_L^{a,J}(\sigma _1)`$ $`=`$ $`(g\stackrel{~}{𝐗}_R)^{a,J}(\pi \sigma _1)`$ (5.27)
$`=`$ $`𝐞\left[{\displaystyle \frac{ap_J}{n}}\right]\stackrel{~}{𝐗}_R^{a,J+1}(\pi \sigma _1)`$
$`=`$ $`𝐞\left[{\displaystyle \frac{ap_J}{n}}\right]𝐞\left[{\displaystyle \frac{aq_{J+1}}{n}}\right]\stackrel{~}{𝐗}_L^{a,J+1}(\sigma _1\pi ).`$
In terms of the oscillators, it can be evaluated as,
$$(g\mathrm{\Omega }\alpha _{sa/n})^{a,J}=𝐞\left[\frac{a}{2n}(2p_J2q_{J+1}1)\right](1)^s\alpha _{sa/n}^{a,J+1}.$$
(5.28)
If we denote the eigenvalues of $`g\mathrm{\Omega }`$ action as $`\mu `$, the same line of argument as in section 2 gives, $`\mu ^m=(1)^{sm}𝐞\left[\frac{a}{2n}R\right]`$ with
$$P=\underset{J=0}{\overset{m1}{}}(2p_J2q_{J+1}1).$$
(5.29)
The constraint (4.9) implies that $`P0`$ mod $`n`$. However, we need to evaluate it in mod $`2n`$ to get the accurate phase factor. The answer is,
$$P\{\begin{array}{cc}n& n\text{,}m\text{ odd}\hfill \\ 0& n\text{ odd, }m\text{ even}\hfill \\ 0,n& n\text{,}m\text{ even}\hfill \end{array}$$
(5.30)
When $`n`$ is even and $`m`$ odd, we do not have any solution to mod n relation
$`2_Jp_J2_Jq_Jm0`$. Although we need some care for $`(1)^s`$ factor, the rest of the calculation is almost the same. When $`P=0`$ it gives the annulus partition function (5.22) and when $`P=n`$, it gives Möbius amplitude, (5.23).
##### $`II_M`$ and $`\stackrel{~}{II}_M`$
Let us denote $`m=2k`$. The mode expansion is same as $`II_A`$ cases (5.1.2). In the similar line of calculation as in $`I_M`$, one obtains the action of $`\mathrm{\Omega }g`$ on the oscillators,
$$\left(\mathrm{\Omega }g\alpha _{sa/n}\right)^{a,J}=𝐞\left[\frac{a}{2n}(2p_J2q_{J+1}1)\right](1)^s\alpha _{sa/n}^{a,J+k+1}.$$
(5.31)
The difference between case $`II_M`$ and $`\stackrel{~}{II}_M`$ is that index $`J`$ in this equation is defined mod $`2k`$ ($`II_M`$) or $`k`$ ($`\stackrel{~}{II}_M`$). If we write $`\alpha ^{a,k+J}=\overline{\alpha }^{a,J}`$, the above relation resembles the Klein bottle case (5.1.1). This is exactly the case for $`\stackrel{~}{II}_M`$. If we use the result in that section, one gets
1. Klein bottle amplitude when $`k`$ is odd with partition function:
$$\underset{r=1}{\overset{\mathrm{}}{}}\frac{1}{1𝐞\left[\frac{2k\tau }{n}\right]},$$
(5.32)
2. Torus amplitude when $`k`$ is even,
$$\underset{r=1}{\overset{\mathrm{}}{}}\frac{1}{(1𝐞\left[\frac{k\tau +p}{n}r\right])(1𝐞\left[\frac{k\tau p}{n}r\right])},$$
(5.33)
with $`p=_{\mathrm{}=0}^{k/21}(p_2\mathrm{}+p_{k+2\mathrm{}+1})_{J=0}^{k1}q_J\frac{k}{2}`$.
In $`II_M`$ case, we need to use “twisted identification” $`\alpha ^{a,k}\stackrel{~}{\alpha }^{a,0}`$. Because of this twist, the correspondence with torus/Klein is toggled between $`k`$ being odd/even.
1. When $`k`$ is even, we get the Klein bottle amplitude (5.32).
2. When $`k`$ is odd, we get the torus amplitude,
$$\underset{r=1}{\overset{\mathrm{}}{}}\frac{1}{(1𝐞\left[\frac{2k\tau +p}{2n}r\right])(1𝐞\left[\frac{2k\tau p}{2n}r\right])},$$
(5.34)
with $`p2_{\mathrm{}=0}^{k1}p_2\mathrm{}2_{J=0}^{k1}q_Jk`$ mod $`2n`$. We have to note that this quantity is always odd as long as we evaluate it in mod $`2n`$.
### 5.2 Derivation by multiple cover
In this subsection, we present the alternative derivation of the table at the beginning of this section based on the multiple cover of the short string world sheet. We first recall that the four sectors are graphically represented in Figure 6.
Here the white (shaded) box represents the world sheet for the left (resp. right) mover. In the torus world sheets, the left and the right movers are independent and are detached from each other. In the Klein bottle, two boxes are piled vertically because of the orientation projection in the trace. In the annulus, they are piled horizontally because of the reflections at the boundaries. In the Möbius strip, they are piled horizontally and vertical at the same time. The fundamental parallelogram of Figure 4 is drawn by dashed line.
In the torus case discussed in section 2.2.2, we first introduce infinite plane with lattice points. We then assign the name of each short string to each rectangle by the rule determined from $`h`$ and $`g`$. In the open case, we need to put the left and right movers in the same plane. In the Klein bottle amplitude, there are toggles between left/right movers in the vertical direction (horizontal stripes). In the annulus, we have the toggling in the horizontal direction (vertical stripes). In the Möbius case, we have the toggling in both direction (checker board). We then put the names of the short string on each rectangle by using rule from $`g`$, $`h`$, and $`f`$ and determine the “fundamental region”. In the following, we will not try to exhaust all combinations appearing in the table. Rather we will illustrate the simple situation which will illuminate the topology change between the world sheets of the short and the long strings.
##### Klein bottle
In this case, there is toggling between torus/Klein bottle amplitude when viewed as the long string diagram. We explain it by taking the simplest situations $`n=1`$ and $`m=2,3`$.
In these cases, the left (resp. right) mover of $`I+1`$’th short string world sheet should be piled over the right (resp. left) mover of $`I`$’th. When $`m=2`$, we have independent two vertical loops, $`L1R2L1`$, $`R1L2R1`$ where $`LI`$ ($`RI`$) means the world sheet of the left (resp. right) mover of $`I`$’th short string. As we illustrate it in Figure 7 left, we have two independent rectangles of size 2. From the viewpoint of the long string, they should be identified as the left and right moving sectors of the long string. It is easy to observe that the similar phenomena occurs whenever $`m`$ is even.
If $`m=3`$, the six boxes should be attached with each other to form one big rectangle. It should be identified as the Klein bottle world sheet of the long string. It is easily generalized that the fat Klein bottle type world sheet appear whenever $`m`$ is odd.
Toggling between Annulus/Möbius strip long string amplitudes in the $`I_M`$ sector can be similarly understood.
##### Annulus
We first explain the appearance of the long open and closed strings by taking the simplest example which consists of two short strings in annulus diagram. We have two choices for $`f_1`$ and $`f_2`$, (A) $`f_1=(1)(2)`$, $`f_2=(12)`$ (B) $`f_1=f_2=(12)`$.
In the first case, the horizontal attachments are defined as in Figure 8 (A) and we have only one chiral world sheet. This is the situation which describes annulus diagram for the long string. In the second case, we have two independent groups which are not attached each other Figure 8 (B). This is again the torus world sheet for the long string. We therefore meet the world sheet of the long closed string.
This is not the end of the story. Even in such a simple situation, we have a degree of freedom of the twist in the vertical direction. Since we are considering the annulus diagram, the world sheet of the left (right) mover should be piled vertically over that of the left (right) mover. Since we have two boxes for each, we have two choices for $`g`$, $`g=1`$ and $`g=(12)`$. In the first case, the diagrams is Figure 8(A,B) themselves. On the other hand, if we take $`g=(12)`$, we get two new world sheets (C,D). In diagram (C), as we draw a dashed line, the obtained diagram should be interpreted as the Möbius strip for the long string. For the diagram (D), two independent rectangles in (B) are piled vertically and gives the Klein bottle world sheet for the long string. In the table, (A) and (C) are classified as $`I_A`$ with $`n=2`$, $`m=1`$. (B) and (D) are classified as $`\stackrel{~}{II}_A`$ and $`II_A`$ respectively with $`n=1`$ and $`m=2`$.
In this way, we get all the four diagrams of the long string amplitude from the annulus of the short string.
##### Möbius strip
In this case, we have the toggling of left and right movers in both horizontal and vertical directions. As in the annulus case, we explain the essence of the correspondence by using two short strip configurations. In figure 9, we illustrate three possible configurations.
The first one (A) corresponds to $`I_M`$ case ($`n=1,m=2`$) with $`f_1=(1)(2)`$, $`g=(12)`$. Because $`m`$ is even, we get the world sheet of annulus diagram. The second one corresponds to $`II_M`$ ($`n=1,m=2`$) with $`f_1=g=(12)`$. As illustrated in the figure (B), we have two independent parallelogram region with a twist by one block. It corresponds to the torus amplitude with $`\tau _{1,1,1/2}`$. The appearance of the half an odd integer is the characteristic feature in $`II_M`$ case. The third one corresponds to $`\stackrel{~}{II}_M`$ case with $`f_1=(12)`$, $`g=(1)(2)`$. As depicted in the figure (C), it describes the world sheet of the Klein bottle
## 6 Boundary states
In this section, we show that there are basically two types of the boundary/cross-cap states in the symmetric product orbifold.
The first one is the conventional boundary/cross-cap states of the long string. One non-trivial point is that the boundary state for the short string $`|B:g,f`$ sometimes describes the cross-cap state of the long string. This is one of the origin of the change of the world sheet topology in the long string.
The second one describes the joint of two short strings. It helps to organize arbitrary long string world sheet from that of the short strings. This is clearly needed if our orbifold CFT has the character of the second quantized string theory. We will call such state as the joint state.
In usual description of the string theory, the boundary state describes the dynamics of the D-brane. In our approach, it is naturally unified into the interaction of the string fields.
### 6.1 Boundary states of short strings
Let us first derive the boundary states for the irreducible sets. Here we use the terminology “irreducible” to mean that it can not written as the direct product of the boundary states for the subset fields.
Let us start from a generic irreducible combination which satisfies $`[g,f]=0`$,
$$g=diag(\stackrel{\overline{m}}{\stackrel{}{\overline{𝒯}_{\overline{n}},\mathrm{},\overline{𝒯}_{\overline{n}}}}),f=diag(\overline{𝒯}_{\overline{n}}^{\overline{p}_0},\mathrm{},\overline{𝒯}_{\overline{n}}^{\overline{p}_{\overline{m}1}})𝒯_{\overline{m}}.$$
(6.1)
The condition $`f^2=1`$ firstly imposes the condition $`\overline{m}=1,2`$. When $`\overline{m}=1`$ (i.e. $`f=\overline{𝒯}_{\overline{n}}^{\overline{p}_0}`$), $`f^2=1`$ implies $`2\overline{p}_00`$ mod $`\overline{n}`$. If $`\overline{n}`$ is odd, the only solution is $`\overline{m}0`$ mod $`\overline{n}`$. If $`\overline{n}`$ is even, we have two solutions, $`\overline{m}0,\overline{n}/2`$. When $`\overline{m}=2`$, $`f^2=1`$ is satisfied if $`\overline{p}_0+\overline{p}_10`$. We are left with only three types of the boundary states,
1. Boundary state of long string: $`(g,f)=(\overline{𝒯}_{\overline{n}},1)`$. If we use the mode expansion of the closed string (2.19) (while exchanging $`n`$ by $`\overline{n}`$), we get the explicit expression for the boundary state,
$$|B:\overline{𝒯}_n,1=\mathrm{exp}(\underset{r=1}{\overset{\mathrm{}}{}}\underset{a=0}{\overset{\overline{n}1}{}}\frac{1}{ra/\overline{n}}\alpha _{r+a/\overline{n}}^{(a)}\stackrel{~}{\alpha }_{r+a/\overline{n}}^{(\overline{n}a)})|0_{\overline{n}}.$$
(6.2)
We introduce the long string oscillator of length $`\overline{n}`$ as
$$𝒜_{\overline{n}pa}\sqrt{\overline{n}}\alpha _{pa/\overline{n}}^{(a)},\stackrel{~}{𝒜}_{\overline{n}pa}\sqrt{\overline{n}}\stackrel{~}{\alpha }_{pa/\overline{n}}^{(\overline{n}a)}.$$
(6.3)
They satisfy a standard commutation relation $`[𝒜_n,𝒜_m]=n\delta _{n+m,0}`$. The commutation relation with Hamiltonian is modified to $`[L_0,𝒜_r]=\frac{r}{\overline{n}}𝒜_r`$. In terms of this variable, the boundary state is rewritten as
$$\mathrm{exp}\left(\underset{r=1}{\overset{\mathrm{}}{}}\frac{1}{r}𝒜_r\stackrel{~}{𝒜}_r\right)|0_{\overline{n}}|B_{\overline{n}}.$$
(6.4)
As for the zero mode, since we are considering the Neumann type boundary condition, we need to impose (writing $`P`$ for the momentum for $`𝒜`$)
$$P|0_{\overline{n}}=0.$$
(6.5)
This is nothing but the standard expression of the boundary state. We will denote the vacuum state for the $`J`$’s long oscillator satisfying (6.5) as $`|0_{(J)}_{\overline{n}}`$ in the following.
2. Cross-cap state for long string: $`(g,f)=(\overline{𝒯}_{\overline{n}},\overline{𝒯}_{\overline{n}}^{\overline{n}/2})`$
This state exist only when $`\overline{n}`$ is even. (3.10) is satisfied by
$`|B:\overline{𝒯}_{\overline{n}},\overline{𝒯}_{\overline{n}}^{\overline{n}/2}`$ $``$ $`|C_{\overline{n}}`$ (6.6)
$`=`$ $`\mathrm{exp}\left({\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{a=0}{\overset{\overline{n}1}{}}}{\displaystyle \frac{(1)^a}{ra/\overline{n}}}\alpha _{r+a/\overline{n}}^{(a)}\stackrel{~}{\alpha }_{r+a/\overline{n}}^{(\overline{n}a)}\right)|0_{\overline{n}}`$
$`=`$ $`\mathrm{exp}\left({\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}(1)^r{\displaystyle \frac{1}{r}}𝒜_r\stackrel{~}{𝒜}_r\right)|0_{\overline{n}}.`$
This is again the standard expression for the cross-cap state of the long string. This is the origin of the mixture of annulus/Möbius amplitude which we have observed in the previous section.
3. Joint state:
$$g=\left(\begin{array}{cc}\overline{𝒯}_{\overline{n}}& 0\\ 0& \overline{𝒯}_{\overline{n}}\end{array}\right),f=\left(\begin{array}{cc}0& \overline{𝒯}_{\overline{n}}^{\overline{p}}\\ \overline{𝒯}_{\overline{n}}^{\overline{p}}& 0\end{array}\right).$$
(6.7)
This boundary state actually interconnects two long strings at the boundary. The boundary condition (3.10) can be easily solved to give,
$`|B:g,f`$ $``$ $`|J(12),\overline{p}_{\overline{n}}`$
$`=`$ $`\mathrm{exp}({\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{a=0}{\overset{\overline{n}1}{}}}{\displaystyle \frac{1}{ra/\overline{n}}}(𝐞\left[{\displaystyle \frac{a\overline{p}}{\overline{n}}}\right]\alpha _{r+a/\overline{n}}^{(a,1)}\stackrel{~}{\alpha }_{r+a/\overline{n}}^{(\overline{n}a,2)}`$
$`+𝐞\left[{\displaystyle \frac{a\overline{p}}{\overline{n}}}\right]\alpha _{r+a/\overline{n}}^{(a,2)}\stackrel{~}{\alpha }_{r+a/\overline{n}}^{(\overline{n}a,1)}))|0_{\overline{n}}`$
$`=`$ $`\mathrm{exp}\left({\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{r}}\left(𝐞\left[{\displaystyle \frac{r\overline{p}}{\overline{n}}}\right]𝒜_r^{(1)}\stackrel{~}{𝒜}_r^{(2)}+𝐞\left[{\displaystyle \frac{r\overline{p}}{\overline{n}}}\right]𝒜_r^{(1)}\stackrel{~}{𝒜}_r^{(2)}\right)\right)|0_{\overline{n}},`$
where $`𝒜^{(I)}`$ ($`I=1,2`$) represents two long string variables.
As for the zero mode, the constraint from Dirichlet type boundary condition is written in terms of the zero mode as,
$$(x_0^1x_0^2)|0_{\overline{n}}=(P^1+P^2)|0_{\overline{n}}=0.$$
(6.9)
We will denote the vacuum state which satisfy above condition for $`I`$th and $`J`$th long string as $`|0_{(IJ)}`$.
### 6.2 Cross-cap states of short strings
In the cross-cap states, $`g`$ and $`f`$ satisfies $`f^2=g`$. Since $`[g,f]=[f^2,f]=0`$, $`g`$ and $`f`$ should be written in the form (6.1). As in the previous subsection, the constraint $`f^2=g`$ implies $`\overline{𝒯}_{\overline{m}}^2=1`$. So the only possibility is $`\overline{m}=1,2`$.
When $`\overline{m}=1`$, $`f=\overline{𝒯}^{\overline{p}}`$ and $`f^2=g`$ gives $`2\overline{p}1`$ mod $`\overline{n}`$. If $`\overline{n}`$ is odd, we have one solution $`\overline{p}=(\overline{n}+1)/2`$. On the other hand if $`\overline{n}`$ is even we have no solution. When $`\overline{m}=2`$, the general solution to $`f^2=g`$ is (6.1) with $`p_0+p_11`$ mod $`n`$. We end up with two class of solutions,
1. Cross-cap states of the long string: $`g=\overline{𝒯}_{\overline{n}}`$ and $`f=\overline{𝒯}_{\overline{n}}^{(\overline{n}+1)/2}`$ for odd $`\overline{n}`$.
In this case one may reshuffle the basis to write $`f=\overline{𝒯}_{\overline{n}}`$ and $`g=\overline{𝒯}_{\overline{n}}^2`$. The mode expansion (2.19) is slightly modified to
$$\stackrel{~}{𝐗}^a=\alpha _0\delta _{a,0}\sigma _0+i\underset{r𝐙}{}\left(\frac{1}{r2a/\overline{n}}\alpha _{r2a/\overline{n}}^az^{r+2a/\overline{n}}+\frac{1}{r+2a/\overline{n}}\stackrel{~}{\alpha }_{r+2a/\overline{n}}^a\overline{z}^{r2a/\overline{n}}\right)$$
(6.10)
The boundary condition (3.28) can be written in terms of oscillators as
$$(\alpha _{r2a/\overline{n}}^a+(1)^r\stackrel{~}{\alpha }_{r+2a/\overline{n}}^{\overline{n}a})|C:f^2,f=0,$$
(6.11)
where the twist factor $`f`$ (3.28) is exactly canceled by the translation of $`\sigma _1`$. Now it is straightforward to write down the cross-cap state,
$$|C:\overline{𝒯}_{\overline{n}}^2,\overline{𝒯}_{\overline{n}}=\mathrm{exp}(\underset{r=1}{\overset{\mathrm{}}{}}\underset{a=0}{\overset{\overline{n}1}{}}\frac{(1)^r}{r2a/\overline{n}}\alpha _{r+2a/\overline{n}}^{(a)}\stackrel{~}{\alpha }_{r+2a/\overline{n}}^{(\overline{n}a)})|0_{\overline{n}}$$
It is not difficult from this point to show that after redefinition of the long string variable, it reduces to (6.6).
2. Joint state:
$$g=\left(\begin{array}{cc}\overline{𝒯}_{\overline{n}}& 0\\ 0& \overline{𝒯}_{\overline{n}}\end{array}\right),f=\left(\begin{array}{cc}0& \overline{𝒯}_{\overline{n}}^{\overline{p}+1}\\ \overline{𝒯}_{\overline{n}}^{\overline{p}}& 0\end{array}\right).$$
(6.12)
After repeating the similar calculation, we arrive at (3) with $`\overline{p}`$ replaced by $`\overline{p}+1/2`$. Clearly this is the origin of half-twist which appeared in the fat torus amplitude (5.34).
In either cases, the condition for the zero-mode are the same as in the boundary states (6.5,6.9).
### 6.3 Inner product between boundary states
In this subsection, we calculate the inner product between boundary states of the short strings to see the modular property of the long open strings. In the following, we will restrict our attraction to the boundary states since the calculation of the cross-cap states are completely analogous.
In order to reproduce the open string amplitudes from the boundary state, we need to loosen the the irreducibility of the boundary state. We therefore start from the general form (6.1) and impose the condition $`f^2=1`$.
The general solution is given by
$$f=diag(\overline{𝒯}_{\overline{n}}^{\overline{p}_0},\mathrm{},\overline{𝒯}_{\overline{n}}^{\overline{p}_{\overline{m}1}})𝒯_{\overline{m}}^q𝒮_{\overline{m}}.$$
(6.13)
The constraint $`f^2=1`$ further imposes
$$\overline{p}_{\mathrm{}}+\overline{p}_{\overline{m}1\mathrm{}+q}0\text{mod}\overline{n}.$$
(6.14)
The corresponding boundary state $`|B:g,f`$ can be decomposed into the product of the (long string) boundary, cross-cap, and joint states. When
$$\mathrm{}\overline{m}1\mathrm{}+q\text{ mod}\overline{m},$$
(6.15)
it is described by the boundary and cross-cap states and otherwise it is given by the joint state. Let us first count the number of the boundary and cross-cap states.
We remark first that $`\overline{m}`$ is the number of long open strings which have their open end at the boundary. When $`\overline{m}`$ is even and $`q`$ is odd, we have two solutions to (6.15), $`\mathrm{}=(q1)/2,(q1+\overline{m})/2`$. In this case, the long string has two loose end at that boundary and the others are connected each other (Figure 10 left). When $`\overline{m}`$ is even and $`q`$ is even, we have no solution and every long strings are jointed each other. On the other hand, when $`\overline{m}`$ is odd, we always have only one loose end (Figure 10 right).
#### 6.3.1 Annulus/Möbius strip/Klein bottle amplitudes
Let us assume that $`f_1`$ and $`f_2`$ (reflection factors at each boundary) to have the general form (6.13)
$$f_1=diag(\overline{𝒯}_{\overline{n}}^{\overline{p}_0^{(1)}},\mathrm{},\overline{𝒯}_{\overline{n}}^{\overline{p}_{\overline{m}1}^{(1)}})𝒯_{\overline{m}}^{q_1}𝒮_{\overline{m}},f_2=diag(\overline{𝒯}_{\overline{n}}^{\overline{p}_0^{(2)}},\mathrm{},\overline{𝒯}_{\overline{n}}^{\overline{p}_{\overline{m}1}^{(2)}})𝒯_{\overline{m}}^{q_2}𝒮_{\overline{m}}.$$
(6.16)
Annulus, Möbius strip and Klein bottle amplitudes can be obtained if there are two loose ends at the boundaries. Such situations can be obtained if (i) $`\overline{m}`$ is odd, or (ii) $`\overline{m}`$ is even and $`q_1q_2`$ is odd. To get the irreducible diagram, one may always obtain a basis where $`q_1=1`$ and $`q_2=0`$ in either cases.
In the explicit evaluation, we use the following formula,
$$0|e^{f_2^{IJ}\alpha _I\stackrel{~}{\alpha }_Jq}e^{f_1^{KL}\alpha _K^{}\stackrel{~}{\alpha }_L^{}}|0=\underset{a}{}\frac{1}{1\mu _aq},$$
(6.17)
where oscillators satisfy the commutation relations $`[\alpha _I,\alpha _J^{}]=[\stackrel{~}{\alpha }_I,\stackrel{~}{\alpha }_J^{}]=\delta _{IJ}`$ and $`\mu _a`$ are the eigenvalues of $`f_2f_1^t`$.
In our case, $`f_2f_1^t=f_2f_1`$ is given by
$$diag(\overline{𝒯}_{\overline{n}}^{\overline{p}_0^{(2)}+\overline{p}_{\overline{m}1}^{(1)}},\overline{𝒯}_{\overline{n}}^{\overline{p}_1^{(2)}+\overline{p}_{\overline{m}2}^{(1)}},\mathrm{})𝒯_{\overline{m}}^1.$$
(6.18)
As we have been doing in section 2, we first diagonalize the $`\overline{𝒯}_{\overline{n}}`$ action by discrete Fourier transformation. In the subspace where $`\overline{𝒯}_{\overline{n}}𝐗^{a,J}=𝐞\left[\frac{a}{\overline{n}}\right]𝐗^{a,J}`$, the eigenvalues of this matrix should satisfy
$$\mu ^{\overline{m}}=𝐞\left[\frac{a}{\overline{n}}\left(\underset{\mathrm{}=0}{\overset{\overline{m}1}{}}(\overline{p}_{\mathrm{}}^{(2)}+\overline{p}_{\overline{m}1\mathrm{}}^{(1)})\right)\right]$$
(6.19)
Because of (6.14), almost all the terms except for $`\overline{p}`$ at the loose ends cancel each other in the right hand side. For the annulus/Klein bottle amplitude, the phase at the loose end cancel each other and $`\mu ^{\overline{m}}=1`$. For the Möbius strip amplitude, we have non-trivial phase $`\mu ^{\overline{m}}=(1)^a`$ since $`\overline{p}=\overline{n}/2`$.
For the oscillator contribution of the annulus and Klein bottle amplitude, we obtain the inner product as,
$`B:g,f_2|q^{L_0+\overline{L}_0}|B:g,f_1^{oscillator}`$ $`=`$ $`{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{a=0}{\overset{\overline{n}1}{}}}{\displaystyle \underset{b=0}{\overset{\overline{m}1}{}}}{\displaystyle \frac{1}{1𝐞\left[b/\overline{m}\right]q^{2(ra/\overline{n})}}}`$ (6.20)
$`=`$ $`{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{1q^{2\overline{m}r/\overline{n}}}}`$
Similarly for the Möbius strip amplitude ($`\overline{n}`$ should be even in this case),
$`B:g,f_2|q^{L_0+\overline{L}_0}|B:g,f_1^{oscillator}`$ $`=`$ $`{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{a=0}{\overset{\overline{n}1}{}}}{\displaystyle \underset{b=0}{\overset{\overline{m}1}{}}}{\displaystyle \frac{1}{1𝐞\left[\frac{b}{\overline{m}}+\frac{a}{2\overline{m}}\right]q^{2(ra/\overline{n})}}}`$ (6.21)
$`=`$ $`{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{1(1)^rq^{2\overline{m}r/\overline{n}}}}`$
For the zero mode contribution to the inner product, we use the momentum representation of the vacuum states,
$$P^I|0_{(I)}_{\overline{n}}=\delta (P^I),P^I,P^J|0_{(IJ)}_{\overline{n}}=\delta (P^I+P^J).$$
(6.22)
If $`\overline{m}`$ is even, the inner product can be written in the following form,
$`{}_{\overline{n}}{}^{}0_{(1)(23)\mathrm{}(\overline{m}2,\overline{m}1)(\overline{m})}|q^{\frac{1}{2\overline{n}}(L_0+\overline{L}_0)}|0_{(12)\mathrm{}(\overline{m}1,\overline{m})}_{\overline{n}}^{}`$
$`={\displaystyle }d^{\overline{m}}P{}_{\overline{n}}{}^{}0_{(1)(23)\mathrm{}(\overline{m}2,\overline{m}1)(\overline{m})}|\stackrel{}{P}q^{\frac{1}{2\overline{n}}_IP_I^2}\stackrel{}{P}|0_{(12)\mathrm{}(\overline{m}1,\overline{m})}_{\overline{n}}`$
$`={\displaystyle d^{\overline{m}}P\delta (P_1)\delta (P_1+P_2)\delta (P_2+P_3)\mathrm{}\delta (P_{\overline{m}1}+P_{\overline{m}})\delta (P_{\overline{m}})q^{\frac{1}{2\overline{n}}_IP_I^2}}`$
$`=\delta (0)=V.`$ (6.23)
Here $`V`$ is the volume and $`|0_{(12)\mathrm{}(\overline{m}1,\overline{m})}_{\overline{n}}`$ is the short hand notation of $`|0_{(12)}_{\overline{n}}\mathrm{}|0_{\overline{m}1,\overline{m}}_{\overline{n}}`$. Calculation for odd $`\overline{m}`$ is mostly the same and gives the same answer.
#### 6.3.2 Torus
In order to get the torus amplitude, $`\overline{m}`$ should be even since the both boundary should have no loose ends. Similarly both $`q_1`$ and $`q_2`$ are even to fulfill this relation. In order to get the irreducible amplitude, one may put $`q_1=2`$ and $`q_2=0`$ without losing generality. In this case $`f_2f_1`$ is given by
$$diag(\overline{𝒯}_{\overline{n}}^{\overline{p}_0^{(2)}+\overline{p}_{\overline{m}1}^{(1)}},\overline{𝒯}_{\overline{n}}^{\overline{p}_1^{(2)}+\overline{p}_{\overline{m}2}^{(1)}},\mathrm{})𝒯_{\overline{m}}^2.$$
(6.24)
Since $`\overline{m}`$ is even, we have two sets of eigenvalue equations,
$$\mu ^{\overline{m}/2}=𝐞\left[\pm \frac{a\overline{p}}{\overline{n}}\right],\overline{p}\underset{\mathrm{}=0}{\overset{\overline{m}/21}{}}(\overline{p}_2\mathrm{}^{(2)}+\overline{p}_{\overline{m}12\mathrm{}}^{(1)}).$$
(6.25)
Unlike the situation in the previous subsection, $`\overline{p}`$ can take any integer.
Writing $`k=\overline{m}/2`$, we obtain the oscillator part of the torus amplitude,
$`B:g,f_2|q^{L_0+\overline{L}_0}|B:g,f_1^{oscillator}`$
$`={\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{a=0}{\overset{\overline{n}1}{}}}{\displaystyle \underset{b=0}{\overset{k1}{}}}{\displaystyle \frac{1}{(1𝐞\left[\frac{b}{\overline{k}}+\frac{a\overline{p}}{\overline{k}}\right]q^{2(ra/\overline{n})})(1𝐞\left[\frac{b}{\overline{k}}\frac{a\overline{p}}{\overline{k}}\right]q^{2(ra/\overline{n})})}}`$
$`={\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(1𝐞\left[\frac{2\overline{k}\tau +\overline{p}}{\overline{n}}\right])(1𝐞\left[\frac{2\overline{k}\tau \overline{p}}{\overline{n}}\right])}}`$ (6.26)
For the zero mode calculation, the calculation is parallel to (6.23). The only difference is the product of the delta functions of $`P`$. The integral produces,
$`{\displaystyle d^{\overline{m}}P\delta (P_1+P_2)\delta (P_2+P_3)\mathrm{}\delta (P_{\overline{m}1}+P_{\overline{m}})\delta (P_{\overline{m}}+P_1)q^{\frac{1}{2\overline{n}}_IP_I^2}}`$
$`=\delta (0){\displaystyle 𝑑Pq^{\frac{\overline{m}}{2\overline{n}}P^2}}=V(\text{Im}\tau _{\overline{n},\overline{m},0})^{1/2}.`$ (6.27)
The moduli dependence appearing here is necessary to have the correct modular property.
## 7 DLCQ partition function
In this section, we prove that the various amplitudes of the open string are organized to give the simple formula similar to (2.54). Such a structure is needed to identify our model as the DLCQ of the open string field theory as we reviewed in section 2.3.
Unlike the closed string partition function, the requirement of the modular invariance does not necessarily determine the combinations of the twisted open string sectors (3.4, 3.14). It is quite encouraging that the summation over the all possible twists indeed gives the DLCQ type partition function.
Usually in BCFT, the combination of various sectors is determined by the tadpole cancellation condition. In terms of the boundary state and the cross cap state, it is written as the cancellation of the massless part of the boundary states,
$$|B_0+|C_0=0.$$
(7.1)
In the bosonic string in the flat target space, this condition determine the gauge group should be $`SO(2^{13})`$. In other situations, this gives a crucial constraint on the model building. While we try to apply the tadpole condition (7.1) naively, we meet one difficulty. Namely the generic one loop amplitude is reducible and its irreducible components have a tachyonic part. Multiplying them usually produces the higher negative modes which complicate the constraint (7.1). It usually becomes a nonlinear relation and the analysis would be very difficult. Such a situation will be remedied if we have the DLCQ type partition function. In this case, products of the various string amplitudes can be organized as the exponential of the sum of the irreducible ones. We can examine the tadpole condition for each of the single long string amplitudes.
We will summarize our results in the following theorem,
Theorem 3 : (i) Klein bottle<sup>5</sup><sup>5</sup>5 While we are typing this manuscript, we noticed the work where authors independently derived the partition function of Klein bottle amplitude. :
The generating function of the partition function,
$$Z_N^{KB}(\tau )=\frac{1}{N!}\underset{g,hS_N}{}\text{Tr}__h(g\mathrm{\Omega }^{closed}𝐞\left[\tau (L_0+\overline{L}_0)\right]),$$
(7.2)
with a constraint $`hg=gh^1`$, is written as
$`\mathrm{ln}({\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}\zeta ^NZ_N^{KB}(\tau ))`$ (7.3)
$`={\displaystyle \underset{n,m=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\zeta ^{n(2m1)}}{2m1}}𝒵^{KB}(\tau _{n,2m1})+{\displaystyle \frac{\zeta ^{2nm}}{2nm}}{\displaystyle \underset{p=0}{\overset{n1}{}}}𝒵^{Torus}(\tau _{n,2m,p},\overline{\tau }_{n,2m,p})\right).`$
(ii) Annulus :
The generating function of the partition function,
$$Z_N^{Annulus}=\frac{1}{N!}\underset{h,g,f_1,f_2S_N}{}\text{Tr}_{_{f_1,f_2}}\left(g𝐞\left[\tau L_0\right]\right),$$
(7.4)
with the constraint, $`[g,f_1]=[g,f_2]=0`$, $`f_1^2=f_2^2=1`$, $`h=f_2^1f_1`$ is written in the following form,
$`\mathrm{ln}({\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}\zeta ^NZ_N^{Annulus}(\tau ))`$
$`={\displaystyle \underset{n,m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\zeta ^{nm}}{m}}𝒵^{Annulus}(\tau _{n,m,0})+{\displaystyle \underset{n,m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\zeta ^{2nm}}{m}}𝒵^{Mobius}(\tau _{2n,m,0})`$
$`+{\displaystyle \underset{n,m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\zeta ^{2nm}}{2nm}}{\displaystyle \underset{p=0}{\overset{n1}{}}}𝒵^{Torus}(\tau _{n,m,p},\overline{\tau }_{n,m,p})+{\displaystyle \underset{n,m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\zeta ^{2nm}}{2m}}𝒵^{KB}(\tau _{n,m,0}).`$
(iii) Möbius strip :
The generating function of the partition function,
$$Z_N^{Mobius}=\frac{1}{N!}\underset{h,g,fS_N}{}\text{Tr}_{_{f,gfg^1}}\left(g\mathrm{\Omega }^{open}𝐞\left[\tau L_0\right]\right),$$
(7.6)
with the constraint, $`h=gfg^1f`$, $`f^2=1`$, $`[g^2,f]=0`$, is given as follows,
$`\mathrm{ln}({\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}\zeta ^NZ_N^{Mobius}(\tau ))`$
$`={\displaystyle \underset{n,m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\zeta ^{2nm}}{2m}}𝒵^{Annulus}(\tau _{n,2m,0})+{\displaystyle \underset{n,m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\zeta ^{(2n1)m}}{m}}𝒵^{Mobius}(\tau _{2n1,m,0})`$ (7.7)
$`+{\displaystyle \underset{n,m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\zeta ^{2nm}}{2nm}}{\displaystyle \underset{p=0}{\overset{n1}{}}}𝒵^{Torus}(\tau _{n,m,p+n/2},\overline{\tau }_{n,2m,p+n/2})+{\displaystyle \underset{n,m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\zeta ^{2nm}}{2m}}𝒵^{KB}(\tau _{n,m,0}).`$
Proof:
Klein bottle: The strategy is completely parallel to our discussion in section 2.2.3. We first choose an element $`h`$ which represents each conjugacy class. For each $`h`$, we need to count the number of the elements $`g`$ which satisfies this constraint. By comparing (2.2.1) and (4.1), we have the same degree of freedom for $`g`$ including the introduction of parameter $`p`$. The only deference is that the Klein bottle amplitude at the end does not depend on $`p`$. Therefore the summation $`\frac{1}{n}_p`$ does not exist. We can conclude that it has the same type of the generating functional (7.3). Annulus and Möbius strip:
We have three class of solutions (4,4,4.6). For given $`N`$, the annulus partition function is the sum of the products of all possible combinations of three irreducible solutions. If we combine them into the generating function, the contributions from each diagram are factorized and we can count them independently. Although we have extra summation over $`f_i`$, the number of the solutions are almost the same if we look at (4,4,4.6) carefully.
For type $`I_A`$ solutions, there are one constraint on $`p`$, $`2_{\mathrm{}}p_{\mathrm{}}0`$. For each value of $`_{\mathrm{}}p_{\mathrm{}}`$, we have $`n`$ solutions for $`q`$ which satisfies (4.3). Since the final expression (5.22) does not depend on $`q`$, we have the same number of degree of freedom as in the torus case. This type of counting holds exactly the same fashion for type $`II_A`$ solutions.
For type $`\stackrel{~}{II}_A`$, we need to be more careful to count the combinations. However, the final expression turns out to be the same as in $`II_A`$ case.
By combining all types of the solutions, we arrive at (7).
The calculation of the combinatorics is exactly the same for the Möbius strip case. QED
## 8 Tadpole Condition
In the previous section, we have seen that four sectors (Torus, Klein bottle, Annulus, Möbius strip) of the long string are actually contained in the annulus diagram of the short open string. This leads us to suspect that we may accomplish the tadpole condition (7.1) by combining the boundary states of the short strings. One subtle issue is that there is no cross-cap state when $`\overline{n}`$ is odd (it is described as the cross-cap state of the short string). In order to achieve the tadpole condition in terms of boundary states alone, we need to restrict the length of the long strings ($`\overline{n}`$) to be even. This condition may be imposed by setting normalization factor for those boundary states is zero.
In the following, we will restrict our discussion to the tadpole cancellation between various amplitudes for the single long string. It is not necessarily equivalent to the tadpole condition in conventional BCFT . Usually the tadpole condition is used to derive consistent compactification of the target space. In our case, however, the symmetric product is used to describe the string field theory. We believe that the our treatment is the appropriate one in this physically different context.
Before we start the discussion, we should comment on the dimension of the target space. In the following we will discuss on the standard bosonic string. In order to recover the Lorentz covariance, we need put the transverse dimension of the target space to be 24. The change of the various amplitude is straightforward.
As was seen in section 6.3, the boundary and cross-cap states of long strings are constructed by the products of irreducible boundary, cross-cap and joint states. The normalizations of these states should be also factorized into those of each irreducible boundary states. We will denote them as $`\kappa _B,\kappa _C`$ and $`\kappa _J`$. The structure of the boundary and cross-cap states of the long strings are slightly different if $`\overline{m}`$ is even or $`\overline{m}`$ is odd.
First, we will discuss odd $`\overline{m}`$ case (Figure 10 right). There is one loose end on each boundary. Long string boundary states are written as
$$|B:g,f^B=\kappa _J^{(\overline{m}1)/2}\kappa _B\underset{\left\{\overline{p}\right\}}{}(\underset{i=1}{\overset{(\overline{m}1)/2}{}}|J(2i1,2i),\overline{p}_i_{\overline{n}})|B_{\overline{n}},$$
(8.1)
where $`g=diag(\stackrel{\overline{m}}{\stackrel{}{\overline{𝒯}_{\overline{n}},\mathrm{},\overline{𝒯}_{\overline{n}}}})`$. $`f^B=diag(\overline{𝒯}_{\overline{n}}^{\overline{p}_0^{(B)}},\mathrm{},\overline{𝒯}_{\overline{n}}^{\overline{p}_{\overline{m}1}^{(B)}})𝒯_{\overline{m}}^{q_B}𝒮_{\overline{m}}`$. It contains a loose end (boundary) on $`\overline{m}`$’th long string. Long string cross-cap states are written as
$$|B:g,f^C=\kappa _J^{(\overline{m}1)/2}\kappa _C\underset{\left\{\overline{p}\right\}}{}(\underset{i=1}{\overset{(\overline{m}1)/2}{}}|J(2i1,2i),\overline{p}_i_{\overline{n}})|C_{\overline{n}},$$
(8.2)
where $`f^C=diag(\overline{𝒯}_{\overline{n}}^{\overline{p}_0^{(C)}},\mathrm{},\overline{𝒯}_{\overline{n}}^{\overline{p}_{\overline{m}1}^{(C)}})𝒯_{\overline{m}}^{q_C}𝒮_{\overline{m}}`$. It contains a cross-cap on $`\overline{m}`$’th long string.
Tadpole cancellation condition 7.1 is factorized into the following conditions,
$`\kappa _B|B_{\overline{n}0}+\kappa _C|C_{\overline{n}0}=0.`$ (8.3)
$`{\displaystyle \underset{\overline{p}_i}{}}|J(2i1,2i),\overline{p}_i_{\overline{n}0}=0.`$ (8.4)
The condition for the joint states (8.4) is satisfied automatically since,
$`{\displaystyle \underset{\overline{p}_i}{}}|J(2i1,2i),\overline{p}_i_{\overline{n}0}`$
$`=\left({\displaystyle \underset{p}{}}𝐞\left[p/\overline{n}\right]\right)(𝒜_r^{(2i1)}\stackrel{~}{𝒜}_r^{(2i)}+𝒜_r^{(2i)}\stackrel{~}{𝒜}_r^{(2i1)})|0_{2i1,2i}=0.`$ (8.5)
There are no constraints for $`\kappa _J`$. For the boundary and the cross-cap states, tadpole condition is satisfied when
$$\kappa _B\kappa _C=0.$$
(8.6)
Let us move to even $`\overline{m}`$ case (Figure 10 left). There are two loose ends on one boundary and no loose ends on the other. The long string boundary states (written as $`|B:g,f^J`$) which does not have any loose ends, don’t make any contributions for the tadpole cancellation conditions since they are factorized into the joint states.
We introduce the long string boundary state for Annulus $`|B:g,f^A`$, Möbius strip $`|B:g,f^{MS}`$, and Klein bottle $`|B:g,f^{KB}`$ as follows,
$`|B:g,f^A`$ $`=`$ $`\kappa _J^{(\overline{m}2)/2}\kappa _B^2({\displaystyle \underset{i=1}{\overset{\overline{m}/21}{}}}{\displaystyle \underset{\overline{p}^i}{}}|J(2i,2i+1),\overline{p}_i_{\overline{n}})|B(1)_{\overline{n}}|B(\overline{m})_{\overline{n}},`$
$`|B:g,f^{MS}`$ $`=`$ $`\kappa _J^{(\overline{m}2)/2}\kappa _B\kappa _C({\displaystyle \underset{i=1}{\overset{\overline{m}/21}{}}}{\displaystyle \underset{\overline{p}_i}{}}|J(2i,2i+1),\overline{p}_i_{\overline{n}})`$ (8.7)
$`(|B(1)_{\overline{n}}|C(\overline{m})_{\overline{n}}+|C(1)_{\overline{n}}|B(\overline{m})_{\overline{n}}),`$
$`|B:g,f^{KB}`$ $`=`$ $`\kappa _J^{(\overline{m}2)/2}\kappa _C^2({\displaystyle \underset{i=1}{\overset{\overline{m}/21}{}}}{\displaystyle \underset{\overline{p}_i}{}}|J(2i,2i+1),\overline{p}_i_{\overline{n}})|C(1)_{\overline{n}}|C(\overline{m})_{\overline{n}}.`$
The name of these states comes from the inner product formulae,
$`Z^A`$ $``$ $`B:g,f^J|q^{L_0+\overline{L_0}}|B:g,f^A,`$
$`Z^M`$ $``$ $`B:g,f^J|q^{L_0+\overline{L_0}}|B:g,f^{MS},`$
$`Z^{KB}`$ $``$ $`B:g,f^J|q^{L_0+\overline{L_0}}|B:g,f^{KB}.`$ (8.8)
Tadpole cancellation condition (7.1) becomes
$`|B:g,f^A_0+|B:g,f^{MS}_0+|B:g,f^{KB}_0`$
$`=\kappa _J^{(\overline{m}2)/2}(({\displaystyle \underset{i=1}{\overset{\overline{m}/21}{}}}{\displaystyle \underset{\overline{p}_i}{}}|J(2i,2i+1)\overline{p}_i_{\overline{n}})`$ (8.9)
$`(\kappa _B|B(1)_{\overline{n}}+\kappa _C|C(1)_{\overline{n}})(\kappa _B|B(\overline{m})_{\overline{n}}+\kappa _C|C(\overline{m})_{\overline{n}}))_0=0.`$
Since it is factorized, it does not produce any new constraints on $`\kappa `$’s.
To determine Chan-Paton factor, we calculate the modular properties of various open string amplitudes. In section 6 we determine the inner products between various boundary states. Together with the normalization factors, annulus/Möbius strip/Klein bottle amplitudes with length $`\overline{m}`$ are given as (6.20, 6.21),
Annulus : $`(\overline{n}\kappa _J)^{\overline{m}1}(\kappa _B)^2𝐞\left[{\displaystyle \frac{\overline{m}}{\overline{n}\tau }}\right]{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}\left(1𝐞\left[{\displaystyle \frac{\overline{m}}{\overline{n}\tau }}r\right]\right)^{24},`$
Möbius : $`2(\overline{n}\kappa _J)^{\overline{m}1}(\kappa _B)^2𝐞\left[{\displaystyle \frac{\overline{m}}{\overline{n}\tau }}\right]{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}\left(1(1)^r𝐞\left[{\displaystyle \frac{\overline{m}}{\overline{n}\tau }}r\right]\right)^{24}`$
KB : $`(\overline{n}\kappa _J)^{\overline{m}1}(\kappa _B)^2𝐞\left[{\displaystyle \frac{\overline{m}}{\overline{n}\tau }}\right]{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}\left(1𝐞\left[{\displaystyle \frac{\overline{m}}{\overline{n}\tau }}r\right]\right)^{24}.`$ (8.10)
To achieve the tadpole condition (8.3), we need to restrict $`\overline{n}`$ to be even. We will write $`\overline{n}=2\overline{k}`$ in the following.
After the modular transformation, these amplitudes are rewritten as,
Annulus : $`(\overline{n}\kappa _J)^{\overline{m}1}(\kappa _B)^2\left({\displaystyle \frac{2\overline{k}\tau }{i\overline{m}}}\right)^{12}𝐞\left[{\displaystyle \frac{2\overline{k}\tau }{\overline{m}}}\right]{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}\left(1𝐞\left[{\displaystyle \frac{2\overline{k}\tau }{\overline{m}}}r\right]\right)^{24},`$
Möbius : $`2(\overline{n}\kappa _J)^{\overline{m}1}(\kappa _B)^2\left({\displaystyle \frac{\overline{k}\tau }{i\overline{m}}}\right)^{12}𝐞\left[{\displaystyle \frac{\overline{k}\tau }{2\overline{m}}}\right]{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}\left(1(1)^r𝐞\left[{\displaystyle \frac{\overline{k}\tau }{2\overline{m}}}r\right]\right)^{24},`$
KB : $`(\overline{n}\kappa _J)^{\overline{m}1}(\kappa _B)^2\left({\displaystyle \frac{2\overline{k}\tau }{i\overline{m}}}\right)^{12}𝐞\left[{\displaystyle \frac{2\overline{k}\tau }{\overline{m}}}\right]{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}\left(1𝐞\left[{\displaystyle \frac{2\overline{k}\tau }{\overline{m}}}r\right]\right)^{24}.`$ (8.11)
These expression should be compared to the partition functions obtained in section 5.
Annulus : $`{\displaystyle \frac{N^2}{4}}\left({\displaystyle \frac{2m\tau }{in}}\right)^{12}𝐞\left[{\displaystyle \frac{m\tau }{n}}\right]{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}\left(1𝐞\left[{\displaystyle \frac{m\tau }{n}}r\right]\right)^{24},`$
Möbius : $`N\eta \left({\displaystyle \frac{m\tau }{ik}}\right)^{12}𝐞\left[{\displaystyle \frac{m\tau }{2k}}\right]{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}\left(1(1)^r𝐞\left[{\displaystyle \frac{m\tau }{2k}}r\right]\right)^{24},`$
KB : $`\left({\displaystyle \frac{\tau k}{in}}\right)^{12}𝐞\left[{\displaystyle \frac{2k\tau }{n}}\right]{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}\left(1𝐞\left[{\displaystyle \frac{2k\tau }{n}}r\right]\right)^{24}.`$ (8.12)
Here $`N`$ is the Chan-Paton factor for the long open strings and we write $`m=2k`$ in Möbius and KB amplitudes since they appear only when $`m`$ is even. By comparing expressions, we first need impose $`\kappa _J=\overline{n}^1`$ since there are no length dependent factors in (8). By comparing the oscillators, we need to impose,
* Annulus: $`m=2\overline{k}`$, $`n=\overline{m}`$
* Möbius: $`m=\overline{k}`$, $`k=\overline{m}`$
* Klein bottle: $`k=\overline{k}`$, $`n=\overline{m}`$
It shows that we need to project to even $`m`$ sector for Annulus in (7) because the restriction that $`\overline{n}`$ is even. For other sectors, it reproduces every terms in (7). By comparing the normalization factor, we get
$`{\displaystyle \frac{1}{4}}N^22^{12}`$ $`=`$ $`(\kappa _B)^2,`$
$`N\eta `$ $`=`$ $`2(\kappa _B)^2,`$
$`1`$ $`=`$ $`2^{12}(\kappa _B)^2.`$ (8.13)
It produces, $`\kappa _B=2^6`$, $`\eta =1`$ and $`N=2^{13}`$. We thus have the standard gauge group $`SO(2^{13})`$ for the bosonic string.
## 9 Discussion
As we mentioned in the introduction, one of the main goal of the current project is to construct the second quantized open string theory which has the powerful handling of D-brane. For this purpose, we studied the detailed combinatorial aspects of the open matrix string theory. We hope that our argument is convincing enough that the theory have quite reasonable structure as the second quantized open string theory.
One distinct merit of current approach to conventional string field theory is the description of D-brane. From the boundary conformal field theoretical viewpoint, the classification of the possible boundary states should be interpretable as the possible geometric configuration of D-branes (for example see ). In our formalism, it is very straightforward to include various D-brane configurations as the description of the loose ends of the long open strings. They can be deformed by introducing the marginal transformations of the short strings.
On the other hand, in the string field theory, we need the information of the D-branes in the very definition of the string fields. Introduction of several D-branes may force us to introduce new string fields and the string vertex operators for each of D-branes. Although this approach is useful in the calculation of the tachyon condensation , the string field theory may not give an economic description of the multiple D-brane background.
This approach is also economical in the description of the string interaction vertices. As already mentioned in , there is only one open string vertex operator $`\mathrm{\Phi }_{KL}`$ which interchanges $`K`$’th and $`L`$’th open strings at the boundary. In terms of the boundary states, this operator mixes the boundary/cross-cap states and the joint states,
$`|J(KL)`$ $``$ $`|B(K)|B(L)`$
$`|J(KM)|J(LN)`$ $``$ $`|J(LM)|J(KN).`$ (9.1)
In the open/closed string field theory, we need to introduce seven types of the string interaction vertices . This is because the the global topology of interaction vertex becomes rather involved and we need the vertex operator for each of them. In the matrix string approach, all of these vertices can be described in terms of one vertex $`\mathrm{\Phi }_{KL}`$ . This is again a great benefit of the current formalism.
In case of the string field theory, the gauge invariance of the string fields requires the gauge group should be $`SO(2^{13})`$ . Although we have not attempted the consistency of the vertex operator, it should reproduce the similar condition. This is one of the most important issues which should be clarified in the future.
In our discussion in section 5, we emphasized that there are no big difference between the boundary/cross-cap states and the joint states. They are three equally possible boundary states from the viewpoint of the short strings. This aspect is clearer in the action of the interaction vertex (9.1) since it mixes three boundary states. While the boundary states are the representation of the D-brane, the joint states represents a kind of string interaction. In this way, we have seen an interesting mixture of the string dynamics and D-brane. The consistency of the interaction $`\mathrm{\Phi }_{KL}`$ will impose the possible deformation of the joint states from the knowledge of the D-branes and vice versa.
In this paper, we do not make an explicit attempt to incorporate the supersymmetry. A straightforward generalization of the closed string was already made in . These aspects are using the same type of the combinatorics as the bosonic situation does not produce extra non-triviality except for the Chan-Paton factor. One of the difficult point which we would like to indicate is that these vertex operators should intertwine Ramond-Ramond boundary states which describe the D-brane charge and the joint states.
Finally we have to mention that current development of the matrix string theory where the nontrivial world sheet topology is interpreted as the instanton sectors of 2D Yang-Mills theory. It is quite interesting to investigate if there are similar description of the open string world sheet as the topologically non-trivial sectors in the Yang-Mills theory.
Acknowledgement: The authors would like to thank N. Ishibashi for sending us his thesis which was the essential resource for us to understand the BCFT on orbifolds. One of the authors (Y.M.) is obliged to T. Kawai for the information on DLCQ type partition function and to T. Kawano for explaining the work .
Y.M is supported in part by Grant-in-Aid ($`\mathrm{}`$09640352) and in part by Grant-in-Aid for Scientific Research in a Priority Area: “Supersymmetry and Unified Theory of Elementary Particle” ($`\mathrm{}`$707) from the Ministry of Education, Science, Sports and Culture.
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# Dynamics and topology of the gauge-invariant gauge field in two-color QCD
## 1 Introduction
The investigation I am reporting concerns the implications of the constraints that apply in QCD — in particular the non-Abelian Gauss’s law — not only for the dynamics but also for the topological features of the theory. Previous work has dealt with the implementation of Gauss’s law, which we have carried out by constructing states that are annihilated by the “Gauss’s law operator” $`\widehat{𝒢}^a(𝐫)`$ for the temporal ($`A_0^a=0`$) gauge, where
$$\widehat{𝒢}^a(𝐫)=_j\mathrm{\Pi }_j^a(𝐫)+gϵ^{abc}A_j^b(𝐫)\mathrm{\Pi }_j^c(𝐫)+j_0^a(𝐫)\text{and}j_0^a(𝐫)=g\psi ^{}(𝐫)\frac{\tau ^a}{2}\psi (𝐫)$$
(1)
and where $`\mathrm{\Pi }_j^a(𝐫)`$ is the momentum conjugate to the gauge field $`A_j^a(𝐫)`$ as well as the negative chromoelectric field. We have, furthermore, constructed gauge-invariant quark and gluon operator-valued fields and have transformed the QCD Hamiltonian into a form in which it is expressed entirely in terms of these gauge-invariant fields. An operator-valued field — the so-called “resolvent field” $`\overline{𝒜_i^\gamma }(𝐫)`$, that is a non-local functional of the gauge field in the temporal gauge — has a central role in these developments. In the two-color SU(2) version of QCD, with which we will be concerned in this work, the resolvent field appears in the gauge-invariant gluon field $`A_{\mathrm{𝖦𝖨}i}^b(𝐫)`$ as shown by
$$[A_{\mathrm{𝖦𝖨}i}^b(𝐫)\frac{\tau ^b}{2}]=V_𝒞(𝐫)[A_i^b(𝐫)\frac{\tau ^b}{2}]V_𝒞^1(𝐫)+\frac{i}{g}V_𝒞(𝐫)_iV_𝒞^1(𝐫).$$
(2)
$$\text{where}V_𝒞(𝐫)=\mathrm{exp}\left(ig\overline{𝒴^\alpha }(𝐫)\frac{\tau ^\alpha }{2}\right)\mathrm{exp}\left(ig𝒳^\alpha (𝐫)\frac{\tau ^\alpha }{2}\right)$$
(3)
$$\text{and}V_𝒞^1(𝐫)=\mathrm{exp}\left(ig𝒳^\alpha (𝐫)\frac{\tau ^\alpha }{2}\right)\mathrm{exp}\left(ig\overline{𝒴^\alpha }(𝐫)\frac{\tau ^\alpha }{2}\right)$$
(4)
$$\text{with}𝒳^\alpha (𝐫)=\frac{_i}{^2}A_i^\alpha (𝐫)\text{and}\overline{𝒴^\alpha }(𝐫)=\frac{_j}{^2}\overline{𝒜_j^\alpha }(𝐫).$$
(5)
The gauge-invariant field $`A_{\mathrm{𝖦𝖨}i}^b(𝐫)`$ given in Eq. (2) can also be expressed as
$$A_{\mathrm{𝖦𝖨}i}^b(𝐫)=A_{Ti}^b(𝐫)+[\delta _{ij}\frac{_i_j}{^2}]\overline{𝒜_j^b}(𝐫),$$
(6)
where $`A_{Ti}^b(𝐫)`$ is the transverse part of the gauge field $`A_i^b(𝐫)`$ in the temporal gauge. The resolvent field is also required for defining the gauge-invariant quark field
$$\psi _{\mathrm{𝖦𝖨}}(𝐫)=V_𝒞(𝐫)\psi (𝐫)\text{and}\psi _{\mathrm{𝖦𝖨}}^{}(𝐫)=\psi ^{}(𝐫)V_𝒞^1(𝐫).$$
(7)
We have shown that gauge transformations transform $`V_𝒞(𝐫)`$ in such a way, that they exactly compensate for the transformations undergone by $`A_i^a(𝐫)`$ and $`\psi (𝐫)`$ so that $`A_{\mathrm{𝖦𝖨}i}^b(𝐫)`$ and $`\psi _{\mathrm{𝖦𝖨}}(𝐫)`$ remain untransformed by any gauge transformations.
Central to the implementation of Gauss’s law and the construction of gauge-invariant operator-valued fields, is a nonlinear integral equation that appeared in our earlier work as the defining equation for the resolvent field. In this nonlinear integral equation, functionals of the gauge field appear as inhomogeneous source terms. The significance of this integral equation is made manifest by a proof — the so-called “fundamental theorem” — that resolvent fields that solve this integral equation generate the gauge-invariant fields given in Eqs. (2)-(7) as well as the states that implement Gauss’s law.
## 2 Topological configurations of the resolvent field
At the QCD Workshop that met in Paris last year, I discussed the transformation of the QCD Hamiltonian into a form that is expressed entirely in terms of gauge-invariant fields. In this form, the QCD Hamiltonian is particularly suited for applications to low-energy phenomena, and displays quark-quark and quark-gluon interactions that are non-Abelian analogs of the Coulomb interaction in QED. More recently, we have been studying the topological implications of the equations that determine the gauge-invariant fields. In this work, we represented the resolvent field and the gauge field in the temporal gauge as functions of spatial variables that are second-rank tensors in the combined spatial and SU(2) indices. Except in so far as the forms of $`\overline{𝒜_i^\gamma }(𝐫)`$ and $`A_i^\gamma (𝐫)`$ reflect this second-rank tensor structure, they are isotropic functions of position. In this way, we can represent the longitudinal part of the gauge field in the temporal gauge as
$$A_i^\gamma {}_{}{}^{L}(𝐫)=\frac{1}{g}[\delta _{i\gamma }\frac{𝒩(r)}{r}+\frac{r_ir_\gamma }{r}\left(\frac{𝒩(r)}{r}\right)^{}]$$
(8)
and the transverse part as
$$A_i^\gamma {}_{}{}^{T}(𝐫)=\delta _{i\gamma }𝒯_A(r)+\frac{r_ir_\gamma }{r^2}𝒯_B(r)+ϵ_{i\gamma n}\frac{r_n}{r}𝒯_C(r)$$
(9)
where $`𝒩(r)`$, $`𝒯_A(r)`$, $`𝒯_B(r)`$ and $`𝒯_C(r)`$ are isotropic functions of $`r`$, the prime denotes differentiation with respect to $`r`$, and the transversality of $`A_i^\gamma {}_{}{}^{T}(𝐫)`$ requires that
$$\frac{d(r^2𝒯_B)}{dr}+r^2\frac{d𝒯_A}{dr}=0.$$
(10)
An entirely analogous representation of the resolvent field enables us to relate it to the gauge field through the nonlinear integral equation described in section 1. As a result of this analysis, we have been able to show that it is possible to represent the resolvent field as
$$\overline{𝒜_i^\gamma }(𝐫)=\left(\delta _{i\gamma }\frac{r_ir_\gamma }{r^2}\right)\left(\frac{\overline{𝒩}}{gr}+\phi _A\right)+ϵ_{i\gamma n}\frac{r_n}{r}\phi _C$$
(11)
where
$$\phi _A=\frac{1}{gr}\left[𝒩\mathrm{cos}(\overline{𝒩}+𝒩)\mathrm{sin}(\overline{𝒩}+𝒩)\right]+𝒯_A\left[\mathrm{cos}(\overline{𝒩}+𝒩)1\right]𝒯_C\mathrm{sin}(\overline{𝒩}+𝒩)$$
(12)
and
$$\phi _C=\frac{1}{gr}\left[𝒩\mathrm{sin}(\overline{𝒩}+𝒩)+\mathrm{cos}(\overline{𝒩}+𝒩)1\right]+𝒯_C\left[\mathrm{cos}(\overline{𝒩}+𝒩)1\right]+𝒯_A\mathrm{sin}(\overline{𝒩}+𝒩)$$
(13)
$$\text{with}\overline{𝒴^a}(𝐫)=\frac{_j}{^2}\overline{𝒜_j^\alpha }(𝐫)=r_a\frac{\overline{𝒩}}{gr}\text{and}𝒳^a(𝐫)=\frac{_i}{^2}A_i^\alpha (𝐫)=r_a\frac{𝒩}{gr}.$$
(14)
Similarly, the gauge-invariant gauge field can be expressed as a functional of $`\overline{𝒩}`$ and of $`𝒩`$, $`𝒯_A`$, $`𝒯_B`$ and $`𝒯_C`$ as shown by
$`A_{\mathrm{𝖦𝖨}i}^\gamma (𝐫)={\displaystyle \frac{1}{gr}}\{ϵ_{i\gamma n}{\displaystyle \frac{r_n}{r}}[\mathrm{cos}(\overline{𝒩}+𝒩)1+𝒩\mathrm{sin}(\overline{𝒩}+𝒩)]+(\delta _{i\gamma }{\displaystyle \frac{r_ir_\gamma }{r^2}})\times `$
$`\times [𝒩\mathrm{cos}(\overline{𝒩}+𝒩)\mathrm{sin}(\overline{𝒩}+𝒩)]{\displaystyle \frac{r_ir_\gamma }{r}}{\displaystyle \frac{d\overline{𝒩}}{dr}}\}`$
$`+𝒯_A\left\{\left(\delta _{i\gamma }{\displaystyle \frac{r_ir_\gamma }{r^2}}\right)\mathrm{cos}(\overline{𝒩}+𝒩)+ϵ_{i\gamma n}{\displaystyle \frac{r_n}{r}}\mathrm{sin}(\overline{𝒩}+𝒩)\right\}+{\displaystyle \frac{r_ir_\gamma }{r^2}}\left(𝒯_A+𝒯_B\right)`$
$`+𝒯_C\left\{ϵ_{i\gamma n}{\displaystyle \frac{r_n}{r}}\mathrm{cos}(\overline{𝒩}+𝒩)\left(\delta _{i\gamma }{\displaystyle \frac{r_ir_\gamma }{r^2}}\right)\mathrm{sin}(\overline{𝒩}+𝒩)\right\}.`$ (15)
With these representations, the nonlinear integral equation that relates the resolvent field to the gauge field is transformed to the nonlinear differential equation
$`{\displaystyle \frac{d^2\overline{𝒩}}{du^2}}`$ $`+`$ $`{\displaystyle \frac{d\overline{𝒩}}{du}}+2\left[𝒩\mathrm{cos}(\overline{𝒩}+𝒩)\mathrm{sin}(\overline{𝒩}+𝒩)\right]`$ (16)
$`+`$ $`2gr_0\mathrm{exp}(u)\left\{𝒯_A\left[\mathrm{cos}(\overline{𝒩}+𝒩)1\right]𝒯_C\mathrm{sin}(\overline{𝒩}+𝒩)\right\}=0`$
where $`u=\mathrm{ln}(r/r_0)`$ and $`r_0`$ is an arbitrary constant. Eq. (16) relates $`\overline{𝒩}`$ to the source terms $`𝒩`$, $`𝒯_A`$ and $`𝒯_C`$. Together with these source terms, $`\overline{𝒩}`$ completely determines the resolvent field and the gauge-invariant gauge field. In the limit in which the gauge field $`A_i^a=0`$, (i. e. $`𝒩=𝒯_A=𝒯_B=𝒯_C=0`$), Eq. (16) reduces to the autonomous “Gribov equation”
$$\frac{d^2\overline{𝒩}}{du^2}+\frac{d\overline{𝒩}}{du}2\mathrm{sin}(\overline{𝒩})=0$$
(17)
which is also the equation for a damped pendulum with $`u`$ representing the time, with the proviso that $`\overline{𝒩}`$ must remain bounded not only in the interval $`0u<\mathrm{}`$, but in the larger interval $`\mathrm{}<u<\mathrm{}`$ to include the entire space $`0r<\mathrm{}`$.
We have obtained numerical solutions of Eq. (16) with a variety of choices for the inhomogeneous terms $`𝒩`$, $`𝒯_A`$, $`𝒯_B`$ and $`𝒯_C`$. We have also defined the part of the gauge field
$$𝖠_i(𝐫)=ig\frac{\tau ^\gamma }{2}\left[A_{\mathrm{𝖦𝖨}i}^\gamma (𝐫)\right]_V=V_C(𝐫)_iV_C^1(𝐫)$$
(18)
that has the structure (but not the physical significance) of a “ pure gauge” field, and have evaluated what we have called the “winding number”
$$Q=(24\pi ^2)^1ϵ_{ijk}𝑑𝐫\text{Tr}[𝖠_i(𝐫)𝖠_j(𝐫)𝖠_k(𝐫)].$$
(19)
In evaluating $`Q`$, we have substituted Eq. (14) into Eq. (3), to obtain
$$V_C(𝐫)=\mathrm{exp}\left(i\widehat{r}_n\tau _n\frac{\left(\overline{𝒩}+𝒩\right)}{2}\right).$$
(20)
In these numerical calculations, we assume that $`𝒯_A`$, $`𝒯_B`$ and $`𝒯_C`$ — the transverse parts of the gauge field — all vanish at the origin and as $`r\mathrm{}`$; and that $`𝒩`$ vanishes at the origin and that, as $`r\mathrm{}`$, $`𝒩2\mathrm{}\pi `$ where $`\mathrm{}`$ is an integer, so that the pure-gauge field associated with the source field $`A_i^a(𝐫)`$,
$$\mathrm{exp}\left(i\widehat{r}_n\tau _n\frac{𝒩}{2}\right)_i\mathrm{exp}\left(i\widehat{r}_n\tau _n\frac{𝒩}{2}\right),$$
has an integer winding number. On the basis of these numerical calculations, we have made the following observations:
$``$ Solutions of Eq. (16) exist that are bounded in the entire interval $`0r<\mathrm{}`$ ($`\mathrm{}<u<\mathrm{}`$). When we normalize these solutions so that $`\overline{𝒩}=0`$ at $`r=0`$, we find that, for cases in which $`\mathrm{}=0`$, $`\overline{𝒩}m\pi `$, where $`m`$ is an integer. Even values of $`m`$ correspond to positions of unstable equilibrium for the equivalent pendulum, and odd value of $`m`$ correspond to stable equilibrium positions. Numerical calculations that we have carried out never terminate with unstable equilibrium positions, because of rounding errors in the numerics. Both even and odd values of $`m`$ are possible in the limit $`r\mathrm{}`$, even though numerical integrations only terminate with odd values of $`m`$ in that limit.
$``$ In the case of solutions of Eq. (16) for which $`lim_r\mathrm{}𝒩=\mathrm{}\pi `$, with $`\mathrm{}`$ an integer other than $`0`$, $`lim_r\mathrm{}\overline{𝒩}=\mathrm{tan}^1(2\pi \mathrm{})`$.
$``$ The “winding number” $`Q`$, defined in Eq. (19), has values $`Q=m/2`$ when $`\mathrm{}=0`$, (i. e. when $`lim_r\mathrm{}𝒩=0`$), and is a half-integer when $`m`$ is odd, corresponding to a position of stable equilibrium in the $`r\mathrm{}`$ limit. When the integer $`\mathrm{}0`$,
$$Q=\frac{1}{2\pi }\left\{\mathrm{arctan}(2\pi \mathrm{})+m\pi +2\pi \mathrm{}\left(1\frac{1}{\sqrt{1+4\pi ^2\mathrm{}^2}}\right)\right\}$$
(21)
where $`m`$ is an arbitrary integer and $`\mathrm{arctan}`$ is defined to represent the first sheet of the multivalued $`\mathrm{tan}^1`$ function.
$``$ The nonlinear Eq. (16) has a number of different solutions, bounded in $`0r<\mathrm{}`$, for the same set of source terms $`𝒩`$, $`𝒯_A`$, $`𝒯_B`$ and $`𝒯_C`$. Moreover, a number of numerical experiments have shown that the solutions of Eq. (16) are remarkably insensitive to these source terms. In particular, regardless of the value of $`lim_r\mathrm{}𝒩`$, in the limit $`r\mathrm{}`$, $`A_{\mathrm{𝖦𝖨}i}^\gamma (𝐫)`$ has the form of the hedgehog
$$A_{\mathrm{𝖦𝖨}i}^\gamma (𝐫)=\left(1\pm \sqrt{1+4\pi ^2\mathrm{}^2}\right)\frac{1}{gr}ϵ_{i\gamma n}\frac{r_n}{r}$$
(22)
where $`\pm `$ corresponds to negative (positive) values of $`\mathrm{sin}\{lim_u\mathrm{}\overline{𝒩}(u)\}`$ and $`\mathrm{cos}\{lim_u\mathrm{}\overline{𝒩}(u)\}`$.
## 3 Discussion
The topological features of the resolvent field — and of the gauge field — that we have discussed in this work are not due to boundary values that we have imposed arbitrarily. They are direct consequences of the integral equation that establishes the resolvent field and that accounts for the implementation of Gauss’s law and the gauge invariance of the fields described in Eqs. (2)-(7). The topological features of the resolvent field and of the gauge-invariant gauge field that we have discussed in Section 2 therefore must be ascribed to the constraint imposed by the implementation of the non-Abelian Gauss’s law and by the ansatz chosen for the representation of the SU(2) gauge field and the resolvent field.
Some of the observations made in Section 2 call for further investigations to clarify the relation between QCD dynamics and topology and between formulations of QCD in different gauges. For example, the observation that, for the same set of source terms, Eq. (16) has multiple solutions, raises questions about the relation of these multiple solutions to Gribov copies of gauge fields in the Coulomb gauge. Also, further work is required to fully understand how the “hedgehog” form of the gauge-invariant gauge field in the large-$`r`$ limit enables it to function as the “glue” in the nonlocal quark-quark and quark-gluon interactions in the representation of QCD in which the Hamiltonian is represented entirely in terms of gauge-invariant fields.
## Acknowledgments
This research was supported by the Department of Energy under Grant No. DE-FG02-92ER40716.00.
## References
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# 1 Values of the inflationary model’s parameters for different initial conditions
On the tensor field inflation in the General Relativity homogeneous cosmological model
Poghos F. Kazarian<sup>1</sup><sup>1</sup>1E-mail: pkazarian@hotmail.com
Received 19/07/2000
Department of Theoretical Physics, Yerevan State University,
1 A. Manoogian St., 375049, Yerevan, Armenia
## Abstract
The homogeneous cosmological model in GR is proposed, where the vacuum energy, which can cause the inflation, is described by tensor field rather than by commonly used in inflationary scenarios scalar field. It is shown that if the initial values of the field are sufficiently big (comparable with the Planck units), under the condition of the tensor field’s slow change in the beginning the regime of the quasiexponential inflation can exist. Numerical solutions for the inflationary stage are obtained that confirm the validity of the approximate solutions. Inflation takes place under wide range of initial conditions provided that the tensor field satisfies the condition imposed on the initial values of the tensor field $`\varphi _{}^{0}{}_{0}{}^{}(0)\varphi _{}^{i}{}_{i}{}^{}(0)`$ (i=1,2,3). That condition also arises from the requirement to satisfy existing observational data.
PACS 98.80.Bp.
Since recently mostly scalar fields were used as the inflation creating force in cosmological inflationary models. Though solving the majority of problems of standard Friedman cosmology, the first models also had their owns \[1-3\]. The chaotic inflation \[4-7\] is free of these problems, but the introduction of the anthropic principle here is necessary (For detailed discussion see, e.g. and the references there).As another way to avoid occurring problems in simple scenarios, models with two scalar fields responsible for the inflation \[9-13\] were proposed, as well as the scalar field driven inflation in non-Einsteinian theories (e.g. scalar-tensor theories \[14-17\]), where one scalar field is already present from the theory itself. Also we’d like to notice so-called bimetric scalar-tensor theory , which allows the extended inflationary regime, like scalar-tensor theories, but has less constraints on itself from observational data (for constraints from inflation on STT, for example, see ). The use of tensor field instead of scalar fields can also be an interesting possibility. On the one hand we could have more degrees of freedom to manipulate, and on the other, these are the components of the same field bound together (let’s mention that the more simple vector field obviously cannot be considered if want to satisfy the condition of the isotropy of Universe). Having more degrees of freedom brings up the requirement for the tensor field to have the initial values in the form $`\varphi _{}^{i}{}_{k}{}^{}b(t_0)diag(1;1;1;1)`$ to achieve the inflationary regime, as it will be shown below. Obviously, scalar fields have not extra degrees of freedom, and the scalar field inflation lacks such a constraint. Let’s suppose here a homogeneous inflationary model of General Relativity with the non-gravitational tensor field, introduced by the Lagrangian
$$L=L_0+L_m$$
(1)
where
$$L_0=\frac{1}{4}\varphi _{ik;l}\varphi ^{ik;l}\frac{1}{2}\varphi _{ik;l}\varphi ^{il;k}+\frac{1}{2}\varphi _{}^{i}{}_{i;k}{}^{}\varphi _{}^{k;l}{}_{l}{}^{}\frac{1}{4}\varphi _{}^{i}{}_{i;k}{}^{}\varphi _{}^{l;k}{}_{l}{}^{},$$
(2)
and
$$L_m=\frac{m^2}{4}(\varphi _{ik}\varphi ^{ik}\varphi _{}^{i}{}_{i}{}^{}\varphi _{}^{k}{}_{k}{}^{})$$
(3)
Here $`\varphi _{ik}`$ is a tensor field, which we choose in the form
$$\varphi _{}^{i}{}_{k}{}^{}=diag(b+\epsilon ;b;b;b)=b(t)\delta _{}^{i}{}_{k}{}^{}+\epsilon (t)\delta _{}^{i}{}_{0}{}^{}\delta _{}^{0}{}_{k}{}^{},$$
(4)
the semicolon denotes the covariant derivative and $`m`$ is the tensor field mass(the light velocity everywhere is $`c=1`$). As usual, let’s choose the metric as \[22-24\]
$$ds^2=dt^2a^2(t)[\frac{dr^2}{1kr^2}+r^2(d\theta ^2+sin^2\theta d\varphi ^2)],$$
(5)
where $`a(t)`$ is the scale parameter or the ”radius” of the Universe. Let’s note that the field (4) is symmetric and satisfies the following relations, where the dot denotes the time derivative:
$`\varphi _{ik}=\varphi _{ki};`$
$`\varphi _{}^{i}{}_{k;l}{}^{}=\varphi _{}^{i}{}_{k,l}{}^{}+\epsilon (\mathrm{\Gamma }_{}^{i}{}_{0l}{}^{}\delta _{}^{0}{}_{k}{}^{}\mathrm{\Gamma }_{}^{0}{}_{kl}{}^{}\delta _{}^{0}{}_{i}{}^{});`$
$`\varphi _{}^{0}{}_{0;\alpha }{}^{}=\varphi _{}^{\alpha }{}_{0;0}{}^{}=\varphi _{}^{0}{}_{\alpha ;0}{}^{}=0;`$
$`\varphi _{}^{\alpha }{}_{\beta ;l}{}^{}=\varphi _{}^{\alpha }{}_{\beta ,l}{}^{}=\dot{b}\delta _{}^{0}{}_{l}{}^{}\delta _{}^{\alpha }{}_{\beta }{}^{};`$
$`\varphi _{}^{0}{}_{0;l}{}^{}=\varphi _{}^{0}{}_{0,l}{}^{}=(\dot{b}+\dot{\epsilon })\delta _{}^{0}{}_{l}{}^{};`$
$`\varphi _{}^{0}{}_{\alpha ;l}{}^{}=\epsilon \mathrm{\Gamma }_{}^{0}{}_{\alpha l}{}^{};\varphi _{}^{\alpha }{}_{0;l}{}^{}=\epsilon \mathrm{\Gamma }_{}^{\alpha }{}_{0l}{}^{};\varphi _{}^{0}{}_{\alpha ;l}{}^{}=\epsilon \mathrm{\Gamma }_{}^{0}{}_{\alpha \beta }{}^{}\delta _{}^{\beta }{}_{l}{}^{};\varphi _{}^{\alpha }{}_{0;l}{}^{}=\epsilon \mathrm{\Gamma }_{}^{\alpha }{}_{0\beta }{}^{}\delta _{}^{\beta }{}_{l}{}^{};`$
$`\varphi _{}^{i}{}_{i}{}^{}=4b+\epsilon ;\varphi ^{ik}\varphi _{ik}=\varphi _{}^{i}{}_{k}{}^{}\varphi _{}^{k}{}_{i}{}^{}=(b+\epsilon )^2+3b^2`$ (6)
In (4) the Latin indices run from 0 to 4 and Greek indices - from 0 to 3. Considering (6), (3) can be rewritten as
$$L_m=\frac{m^2}{4}([(b+\epsilon )^2+3b^2](4b+\epsilon )^2)=\frac{3}{2}m^2b(2b+\epsilon )$$
(7)
Considering that
$`\varphi _{ik;l}\varphi ^{ik;l}=[(\dot{b}+\dot{\epsilon })^2+3\dot{b}^2]2\epsilon ^2g^{\mu \sigma }\mathrm{\Gamma }_{}^{\alpha }{}_{0\gamma }{}^{}\mathrm{\Gamma }_{}^{0}{}_{\alpha \sigma }{}^{};`$
$`\varphi _{ik;l}\varphi ^{il;k}=(\dot{b}+\dot{\epsilon })^2\epsilon ^2g^{\alpha \beta }\mathrm{\Gamma }_{}^{0}{}_{\alpha \gamma }{}^{}\mathrm{\Gamma }_{}^{\gamma }{}_{0\beta }{}^{}+2\epsilon b\dot{ln}\sqrt{g};`$
$`\varphi _{}^{i}{}_{i;k}{}^{}\varphi ^k;l_l=(4\dot{b}+\dot{\epsilon )})[(\dot{b}+\dot{\epsilon })^2\epsilon g^{\alpha \beta }\mathrm{\Gamma }_{}^{0}{}_{\alpha \beta }{}^{}];\varphi _{}^{i}{}_{i;k}{}^{}\varphi ^l;k_l=(4\dot{b}+\dot{\epsilon )})^2`$ (8)
for $`L_0`$, determined by (2), we obtain
$`L_0={\displaystyle \frac{3}{2}}\dot{b}^2+{\displaystyle \frac{1}{2}}\epsilon (2\dot{b}+\dot{\epsilon })\dot{ln}\sqrt{g};`$
$`\dot{ln}\sqrt{g}=3{\displaystyle \frac{\dot{a}}{a}}3H`$ (9)
where, as usual, $`H`$ denotes the Hubble parameter. Finally, the total action can be written as
$$S_{tot}=\frac{3}{2}\{\sqrt{g}𝑑\mathrm{\Omega }[\frac{1}{4\pi G}(\frac{\ddot{a}}{a}+\frac{\dot{a}}{a}^2+\frac{k}{a^2})+\frac{\dot{a}}{a}\epsilon (2\dot{b}+\dot{\epsilon })\dot{b}^2m^2b(2b+\epsilon )]\}$$
(10)
where $`G`$ is the Newton constant, $`k=1;0;1`$ for open, flat and closed Universe and $`\sqrt{g}=\frac{a^3(t)r^2sin\theta }{\sqrt{1kr^2}}`$. The variation of (10) will give us the field equations or the equations of motion. Varying (1) by $`\epsilon `$ , $`b`$ , and scale factor $`a`$ will provide us with the field equations:
$`2\dot{b}H=\epsilon (\dot{H}+3H^2)+m^2b;`$
$`\dot{\epsilon }H+\dot{H}\epsilon \ddot{b}+3H(\epsilon H\dot{b})+m^2(2b+\epsilon /2)=0;`$
$`{\displaystyle \frac{1}{4\pi G}}(2\dot{H}+3H^2+{\displaystyle \frac{k}{a^2}})\{\epsilon (2\dot{b}+\dot{\epsilon })\}^.\dot{b}^23m^2b(2b+\epsilon )=0`$ (11)
Using the first equation of (11), the second one can be also rewritten as
$$\ddot{b}=H(\dot{\epsilon }\dot{b})+m^2(b+\epsilon /2)$$
(12)
Assuming that in the beginning (at times comparable with the Planck time $`t_p`$ ) the tensor field was $`\varphi _{}^{i}{}_{k}{}^{}M_p`$ (where $`M_p`$ is the Planck mass) and was slowly changing, for we can ignore the second derivatives ($`\ddot{\varphi }_{}^{i}{}_{k}{}^{},(\dot{\varphi }_{}^{i}{}_{k}{}^{})^2V(\varphi _{}^{i}{}_{k}{}^{})`$) , and also considering that the scale factor $`a`$ is sufficiently big (the term $`\frac{k}{a^2}`$ can be ignored ), from (12) we obtain
$$\dot{H}=\frac{m^2}{2}\frac{2+\dot{\epsilon }/\dot{b}}{1\dot{\epsilon }/\dot{b}},$$
(13)
and, using field equations,
$$H^2=m^2(4\pi Gb(2b+\epsilon )+\frac{1}{3}\frac{2+\dot{\epsilon }/\dot{b}}{\dot{\epsilon }/\dot{b}1})$$
(14)
From here it’s obvious, that $`\dot{H}H^2`$ (for the change of $`b`$ is comparable to the one of $`\epsilon `$ at $`b,\epsilon M_p`$ ) and for slowly changing matter field ($`\dot{\varphi }_{}^{i}{}_{k}{}^{}/\varphi _{}^{i}{}_{k}{}^{}`$) in the early stages of the existence of Universe $`Hconst`$ . In this approximation the field equations (11) become
$`2\dot{\epsilon }H=3H^2\epsilon m^2(b+\epsilon );`$
$`2\dot{b}H=3H^2\epsilon +m^2b;`$
$`{\displaystyle \frac{1}{4\pi G}}H^2=m^2b(2b+\epsilon ){\displaystyle \frac{2}{3}}V(\varphi _{}^{i}{}_{k}{}^{})`$ (15)
Let’s study the behavior of $`b`$ and $`\epsilon `$ . Introducing new variables $`A=\epsilon /b,B=\dot{b}/b`$ , from (15) follows
$`\dot{A}+ABB+E(2+A)=0;`$
$`BE={\displaystyle \frac{3}{2}}HA;`$
$`E{\displaystyle \frac{m^2}{2H}}=const`$ (16)
After integration
$$t=\frac{1}{A_1A_2}ln|\frac{AA_2}{AA_1}|,$$
(17)
where
$$A_{1,2}=(\mu +\frac{1}{2})\pm \sqrt{(\mu +\frac{1}{2})^2\mu }$$
(18)
Here $`\mu \frac{m^2}{3H^2}0`$ . Finally,
$$A(t)=\frac{A_2A_1C_0exp[3H\sqrt{(\mu +\frac{1}{2})^2\mu }t]}{1C_0exp[3H\sqrt{(\mu +\frac{1}{2})^2\mu }t]},$$
(19)
and $`C_0=\frac{A(0)A_2}{A(0)A_1}`$. In the initial period $`Aconst`$ , and from (16)
$`bb_0exp({\displaystyle \frac{3}{2}}H(A+\mu )t)b_0;`$
$`\epsilon Ab;Aconst`$ (20)
Because
$`bb_0+([3A+{\displaystyle \frac{m^2}{H^2}}{\displaystyle \frac{Ht}{2}}]);`$
$`(V(\varphi _{}^{i}{}_{k}{}^{}))=V(b)m^2b_{0}^{}{}_{}{}^{2}(A2);`$ (21)
then for the scale factor $`a`$ we’ll obtain
$$aa_0exp(\frac{6\pi }{M_{p}^{}{}_{}{}^{2}}[b_{0}^{}{}_{}{}^{2}b^2(t)](A2))$$
(22)
If the term $`(A2)`$ in (20) is positive, then we have the (quasi)exponential growth of the scale factor in this model, which leads to the requirement $`A2`$. Let’s also note that the full expansion of Universe under these conditions will be
$$Pexp(\frac{3b_{0}^{}{}_{}{}^{2}(A2)}{M_{p}^{}{}_{}{}^{2}})exp(\frac{\pi \sqrt{2}M_{p}^{}{}_{}{}^{2}}{3m^2}),$$
(23)
considering that $`\epsilon _0,b_0M_p`$ , the initial values of the tensor field components are approximately of the Planck mass. We’ve shown that under certain conditions there is a possibility of quasiexponential inflation in the Universe with the vacuum energy defined by a tensor field (1)-(3). However, questions still remain about the realization of these conditions in the early Universe, like for how long the approximations of the proposed model can describe with sufficient precision the real evolution of the Universe and for how long this inflationary regime will last. The numerical solutions of the field equations (11) show that the approximation ( $`Hconst`$, $`\dot{\varphi }_{}^{i}{}_{k}{}^{}/\varphi _{}^{i}{}_{k}{}^{}const`$, and $`\ddot{\varphi }_{}^{i}{}_{k}{}^{},(\dot{\varphi }_{}^{i}{}_{k}{}^{})^2V(\varphi _{}^{i}{}_{k}{}^{})`$) is valid through the inflationary stage. The requirements necessary for the inflationary regime realization are
$`V_0(\varphi _{}^{i}{}_{k}{}^{})1;`$
$`|2b_0+\epsilon _0||b_0|,|\epsilon _0|;`$
$`0<H_0|b_0|,|\epsilon _0|;`$ (24)
(here and below $`c1`$ , $`G=M_{p}^{}{}_{}{}^{2}1`$). As far as $`\dot{b},\dot{\epsilon }`$ satisfy (24), we have considerable freedom of choice of their exact initial values, for that does not affect noticeably the numerical solutions because in the negligible time after the beginning $`|b_0|,|\epsilon _0|=const1`$ and remain so during the inflation. The choice of $`H_0`$ also does not affect qualitatively the solutions, although quantitatively $`t_{infl}`$ (the time of increase of $`a(t)`$ in $`10^{30}`$ times) is somewhat shorter for greater $`H_0`$. Otherwise, regardless to the initial value, $`H`$ rapidly becomes (and stays during the inflation) constant: $`H=const1`$. The decrease (increase) of the tensor field mass causes the increase (decrease) of $`t_{infl}`$. Also, for a greater mass the greater $`b_0,\epsilon _0`$ are required to satisfy (24) and achieve the satisfactory inflationary regime. It seems that the only strict requirement for achieving inflation is $`|2b_0+\epsilon _0||b_0|,|\epsilon _0|`$ (or, equivalently, $`\epsilon _02b_0`$ ).Any noticeable deviation from that condition makes the inflationary regime impossible. Because $`bb_0=const,\epsilon \epsilon _0=const`$ during the inflation, $`|2b_0+\epsilon _0|`$ also stays at its initial value. The behavior of scale parameter $`a(t)`$ doesn’t differ qualitatively for different sets of parameters. It increases about $`10^{30}`$ times during $`t_{infl}`$. The growth of $`a(t)`$ becomes somewhat faster (though of the same order of magnitude) with the increase of the mass and initial values of the tensor field components and the decrease of $`|2b_0+\epsilon _0|`$. Values of parameters for different initial conditions are presented in the Tables 1-3. Observations impose several constraints on initial values and field parameters (see, e.g. and the references there). For instance, the requirement of having the inflationary regime for sufficiently long time ($`t_{infl}M_p`$)imposes the upper limits on the field mass ($`m10^6`$). Also, the density fluctuations should be not greater than $`\frac{\delta \varrho }{\varrho }=10^410^5`$, or
$$\frac{\delta \varrho }{\varrho }=\frac{1}{V(\varphi _{}^{i}{}_{k}{}^{})}\frac{dV}{d\varphi _{}^{i}{}_{k}{}^{}}\delta \varphi _{}^{i}{}_{k}{}^{}\frac{\delta (2b+\epsilon )}{2b+\epsilon }\frac{H}{2\pi (2b+\epsilon )}=10^410^5;$$
(25)
or $`2b+\epsilon \frac{H}{2\pi }(10^410^5)`$ (considering $`|2b+\epsilon ||b|`$). So (25) restricts the value of $`(2b+\epsilon )`$ further for a given $`H`$, which, in its turn, restricts the values of $`b`$ and $`\epsilon `$ (therefore, the values of $`b_0`$ and $`\epsilon _0`$, for $`b,\epsilon const`$) to satisfy (24) (see the table). Numerical simulations show that for a broad range of the tensor field parameters defined by (24), (25) the proposed model allows the existence of the inflationary stage of the expansion of the Universe, and those solutions are stable for broad range of initial values (satisfying (24), (25)). Numerical solutions also confirm the predictions of the approximate solutions (19)-(22) and show that the approximation is valid during the inflation. From the comparison of the numerical and approximate solutions it also follows that the parameter $`A`$ in (19)-(22)is constant during the inflation and
$$A2$$
(26)
for achieving the inflationary regime in the model. Let’s note that the condition $`\epsilon _02b_0`$ imposed on the initial values of the field is the requirement that inflation in the model can be achieved if the field is chosen initially in the form of $`\varphi _{}^{i}{}_{k}{}^{}(0)diag(b;b;b;b)`$, where $`bb(t)b_0`$ during the inflation.
References
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# BROWN-HET-1219, DAMTP-2000-41, OUTP-00-18P Cosmological Perturbations in Brane-World Theories: Formalism
## I Introduction
In recent years, the way in which string-theory is believed to be connected to observable physics has changed dramatically. The new viewpoint is mainly due to two ideas, namely the brane-world idea \[VI\]–\[VI\] and the idea that couplings and scales of additional dimensions are much more flexible than previously assumed \[VI,VI,VI,VI\]. Not only have these ideas led to new directions in M-theory phenomenology and, more generally, string-theory inspired particle phenomenology, but also in early universe cosmology.
Much of the recent activity in brane-world cosmology is centered around five-dimensional brane-world theories, see for example \[VI\]–\[VI\]. A large class of such theories arises from heterotic M-theory \[VI,VI,VI\]. Other five-dimensional models have been introduced in ref. \[VI\]–\[VI\] which may provide an alternative solution to the hierarchy problem.
A central question is whether the possible existence of a brane-world and large additional dimensions in the early universe leads to observable consequences today. Specifically, cosmological perturbations as, for example, observed in the cosmic microwave background provide us with a window to the early universe that, perhaps, can be used to test the brane-world idea. It is with this motivation in mind, that we set out to study metric perturbations in brane-world models. It may not be immediately clear that the existence of additional dimensions and branes should have important consequences for the formation and evolution of cosmological perturbations. Let us, as a comparison, consider “traditional” Kaluza-Klein cosmologies, where the higher-dimensional universe is usually split into a product of maximally symmetric subspaces each one with an individual scale factor. Cosmological perturbations are normally treated at the linearized level and, hence, in such Kaluza-Klein cosmologies \[VI\]–\[VI\] internal and external fluctuations basically decouple. The additional dimensions effect perturbations of the three-dimensional universe only via the kinematics of the background scale factors (and/or dilaton background fields \[VI\]). The situation is quite different for brane-world theories. The branes are stretched across the three-dimensional universe and are located at specific points in the internal space. Moreover, they carry worldvolume fields that can only propagate on the brane and that are likely to be excited in the early universe, both coherently and thermally. As a consequence, the cosmological background in an early brane-world universe is highly inhomogeneous in the additional dimensions since the branes constitute localized sources of stress-energy. Even if perturbations around such a background are treated at the linearized level, the perturbations of the three-dimensional universe are effected by the non-linear distortion of the cosmological background in the internal dimensions. This constitutes a crucial difference between conventional Kaluza-Klein cosmology and brane-world cosmology which is directly related to the presence of branes. It is this difference that may lead to new predictions for cosmological perturbations in brane-world models and that motivates the present investigation.
In this paper, we will develop a formalism for metric fluctuations in brane-world theories that takes the characteristic property of brane-world cosmologies, the above mentioned inhomogeneity in the additional dimensions, into account. We understand such a formal development as a first step towards analyzing predictions for cosmological perturbations in brane-world theories. In the next section, we start out to generalize the well-known formalism of four-dimensional gauge-invariant metric perturbations \[VI\], \[VI\] to brane-world theories with an arbitrary number of additional dimensions. Subsequently, in section III, we focus on five-dimensional brane-world models on the orbifold $`S^1/Z_2`$ related to those originating from heterotic M-theory \[VI,VI,VI\]. Specifically, we consider the five-dimensional Einstein equation coupled to bulk as well as brane stress-energy. For such a theory we derive the equations of motion for scalar metric perturbations in a generalized longitudinal gauge. Consistency of these equations is used to determine the most general structure of the stress-energy on the brane. In section IV, these results are applied to find the Israel matching conditions \[VI\] for the scalar metric perturbations restricted to the branes. Finally, in section V, we show how our formalism for brane-world metric perturbations is related to the conventional one in four-dimensions. This is done by a matching procedure applied in a limit where the five-dimensional brane-world theory has an effective four-dimensional description.
## II Gauge-invariant variables
In this section, we will develop a gauge-invariant formalism for metric perturbations in brane–world models. Using such a gauge-invariant approach is particularly useful in order to identify the correct physical degrees of freedom. Once this has been done, a specific gauge can be chosen in order to simplify the subsequent equations. Specifically, we will use a generalized longitudinal gauge later on. However, as a warm-up for the higher–dimensional case, we would first like to review the well-known four-dimensional gauge-invariant formalism following ref. \[VI\].
### A The four-dimensional formalism
Starting point is a background metric with a maximally symmetric three–dimensional spatial space. This metric is of the form
$$ds_4^2=a_4^2\left\{dt^2\mathrm{\Omega }_{ij}dx^idx^j\right\},$$
(1)
where indices $`i,j,\mathrm{}=1,2,3`$ run over the three spatial indices. Indices $`\mu ,\nu ,\mathrm{}=0,1,2,3`$ are used to index four-dimensional space-time. Furthermore, $`a_4=a_4(t)`$ is the four-dimensional scale factor and $`\mathrm{\Omega }_{ij}`$ is the metric of the three–dimensional maximally symmetric space explicitly given by
$$\mathrm{\Omega }_{ij}=\delta _{ij}\left[1+\frac{1}{4}kx^lx^m\delta _{lm}\right]^2.$$
(2)
Here $`k=0,1,1`$ corresponds to a flat, closed or open universe, respectively. The idea is now to classify perturbations of the metric (1) according to their transformation properties with respect to the maximally symmetric space. This leads to the perturbed metric
$$ds_4^2=a_4^2\left\{(1+2\varphi _4)dt^2\left[(12\psi _4)\mathrm{\Omega }_{ij}+2E_{4|ij}+2F_{4(i|j)}+h_{4ij}\right]dx^idx^j+W_{4i}dtdx^i\right\}.$$
(3)
Here and in the following four-dimensional quantities are indexed by “$`4`$” to distinguish them from their higher-dimensional counterparts to be introduced later. The vertical bar refers to a covariant derivative with respect to the metric $`\mathrm{\Omega }_{ij}`$. The vector $`F_{4i}`$ has a vanishing divergence, that is $`F_{4i}^{}{}_{}{}^{|i}=0`$ and the tensor $`h_{4ij}`$ is traceless and divergence-less, that is $`h_{4i}^{}{}_{}{}^{i}=0`$ and $`h_{4i}^{j}{}_{|j}{}^{}=0`$. In addition, we can decompose the off-diagonal perturbation $`W_{4i}`$ further into the gradient of a scalar $`B_4`$ and a divergence-less vector $`S_{4i}`$. Explicitly, this reads
$$W_{4i}=B_{4|i}+S_{4i}.$$
(4)
Consequently, we have four scalar metric perturbations $`(\varphi _4,\psi _4,E_4,B_4)`$, two vector perturbations $`(F_{4i},S_{4i})`$ and a tensor perturbation $`(h_{4ij})`$. All these perturbations are functions of time as well as of the spatial coordinates $`x^i`$, of course. Next we consider an infinitesimal coordinate transformation
$$x^\mu \stackrel{~}{x}^\mu =x^\mu +\xi ^\mu ,$$
(5)
where the vector $`\xi ^\mu `$ depends on all four coordinates, in general. The corresponding infinitesimal change of the metric is given by
$$g_{4\mu \nu }\stackrel{~}{g}_{4\mu \nu }=g_{4\mu \nu }2_{(\mu }\xi _{\nu )}.$$
(6)
To understand how this coordinate transformation acts on the metric perturbations we split $`\xi ^\mu `$ as $`\xi ^\mu =(\xi ^0,\xi ^i)`$ into a time and a spatial part. The spatial component $`\xi ^i`$ can be decomposed further into a gradient and a divergence-less part as
$$\xi ^i=\xi ^{|i}+\eta ^i.$$
(7)
As a result, the transformation parameter $`\xi ^\mu `$ contains two scalar components $`(\xi ^0,\xi )`$ and one vector component $`(\xi ^i)`$. Given this setup, one can compute the transformation properties of the metric perturbations by applying eq. (6) to the perturbed metric (3) and taking into account that $`\xi _\mu =a_4^2(\xi _0,\xi _i)`$. For the scalar perturbations one finds
$`\stackrel{~}{\varphi }_4`$ $`=`$ $`\varphi H_4\xi ^0\dot{\xi }^0,`$ (8)
$`\stackrel{~}{\psi }_4`$ $`=`$ $`\psi _4+H_4\xi ^0,`$ (9)
$`\stackrel{~}{B}_4`$ $`=`$ $`B_4+\xi ^0\dot{\xi },`$ (10)
$`\stackrel{~}{E}_4`$ $`=`$ $`E_4\xi .`$ (11)
Here, $`H_4`$ is the Hubble parameter defined by $`H_4=\dot{a}_4/a_4`$. The vector perturbations transform as
$`\stackrel{~}{F}_{4i}`$ $`=`$ $`F_{4i}\eta _i,`$ (12)
$`\stackrel{~}{S}_{4i}`$ $`=`$ $`S_{4i}\dot{\eta }_i,`$ (13)
while the tensor perturbation $`h_{4ij}`$ is invariant. In these equations, spatial indices are lowered and raised with the metric $`\mathrm{\Omega }_{ij}`$, that is, for example $`\xi _{|i}=\mathrm{\Omega }_{ij}\xi ^{|j}`$ and $`\eta _i=\mathrm{\Omega }_{ij}\eta ^j`$. With these results, it is straightforward to introduce the following gauge-invariant variables.
Scalar variables
$`\mathrm{\Phi }_4`$ $`=`$ $`\varphi _4+H_4(B_4\dot{E}_4)+\dot{B}_4\ddot{E}_4`$ (14)
$`\mathrm{\Psi }_4`$ $`=`$ $`\psi _4H_4(B_4\dot{E}_4).`$ (15)
Vector variables
$`_i=S_{4i}\dot{F}_{4i}.`$ (16)
Tensor variables
$$h_{4ij}.$$
(17)
As the physical degrees of freedom, one has therefore identified two scalar perturbations, one vector perturbation and one tensor perturbation. Of particular importance are the two scalar perturbations $`\mathrm{\Phi }_4`$ and $`\mathrm{\Psi }_4`$ on which we are going to focus. The expressions (14) and (15) for these perturbations suggest the gauge choice $`B_4=E_4=0`$ in the scalar sector which is referred to as longitudinal gauge. Clearly, from the above transformation properties of the scalar perturbations such a choice can be made. Then, the gauge-invariant scalar variables coincide with the “original” variables, that is, $`\mathrm{\Phi }_4=\varphi _4`$ and $`\mathrm{\Psi }_4=\psi _4`$. This gauge choice considerably simplifies subsequent calculations and its generalization will be quite helpful to deal with the higher-dimensional case. The perturbed metric then takes the form
$$ds_4^2=a_4^2\left\{(1+2\varphi _4)dt^2(12\psi _4)\mathrm{\Omega }_{ij}dx^idx^j\right\}.$$
(18)
Finally, we need to specify the stress-energy. For the background, by the maximal symmetry of the three-dimensional spatial space, it is dictated to be of the form
$$T_{4}^{}{}_{}{}^{\mu }{}_{\nu }{}^{}=\text{diag}(\rho _4,p_4,p_4,p_4).$$
(19)
The scalar perturbations to this stress-energy can be written as
$$\delta T_{4}^{}{}_{}{}^{\mu }{}_{\nu }{}^{}=\left(\begin{array}{cc}\delta \rho _4& (\rho _4+p_4)a_4^1v_{4|j}\\ (\rho _4+p_4)a_4v_{4}^{}{}_{}{}^{|i}& \delta p_4\delta _{}^{i}{}_{j}{}^{}+\sigma _{4}^{}{}_{}{}^{|i}{}_{|j}{}^{}\end{array}\right)$$
(20)
with the potential $`v_4`$ for the velocity field $`v_{4|i}`$ and the quantity $`\sigma `$ specifying the anisotropic stress. The equations of motion for the background and the scalar perturbations subject to the above stress-energy are given in ref. \[VI\] and will not be repeated here. These equations form the basis for the study of cosmological perturbations and we now turn to develop their higher-dimensional generalization.
### B Gauge-invariant variables in brane-world theories
We would now like to proceed in close analogy with the four-dimensional case reviewed above and develop a gauge-invariant formalism of metric perturbations in brane-world theories. First, we consider the general situation of $`d`$ additional dimensions although later we will be more specific and focus on the case $`d=1`$, that is, a five-dimensional universe. The coordinates of the additional dimensions are denoted by $`(y^5,\mathrm{},y^{4+d})`$. For the purpose of this subsection, all we need to specify is that the branes are stretched across the usual four-dimensional space-time and are located at specific points (or submanifolds) in the additional dimensions. We will be more precise about the brane positions later when we consider the five-dimensional case.
How should metric perturbations be classified in such a brane-world theory? In the previous four-dimensional case we have used their tensor properties with respect to the three-dimensional spatial subspace for this classification. At first glance, one might now want to use their tensor properties with respect to the $`3+d`$-dimensional spatial space. The cosmological principle, of course, only asserts the maximal symmetry of the usual three-dimensional space but the maximal symmetry of the $`d`$-dimensional internal space may be taken as an additional, simplifying assumption. It is at this point, that the brane-world nature of the theory comes into the game. Since the branes are localized in the additional dimensions the assumption of maximal symmetry cannot, in general, be extended to those dimensions. In fact, as will become more explicit below, the branes lead to stress-energy in the Einstein equation localized in the additional dimensions and, hence, the symmetry of the background metric will typically not be enhanced with respect to the four-dimensional case. Consequently, we split the coordinates into two groups, namely the inhomogeneous coordinates $`(y^a)=(t,y^5,\mathrm{},y^{4+d})`$ on which the background metric generally depends in a non-trivial way and the usual three spatial coordinates $`(x^i)`$ corresponding to the maximally symmetric space. In the following we use indices $`a,b,\mathrm{}=0,5,\mathrm{},4+d`$ for time and the additional dimensions, indices $`i,j,\mathrm{}=1,2,3`$ for the three–dimensional space and indices $`\alpha ,\beta ,\mathrm{}=0,1,2,3,5,\mathrm{},4+d`$ for the full $`4+d`$-dimensional space-time. Then the most general higher-dimensional metric consistent with the maximally symmetric three-dimensional spatial manifold is given by
$$ds^2=a^2\left\{\gamma _{ab}dy^ady^b\mathrm{\Omega }_{ij}dx^idx^j\right\},$$
(21)
where the scale factor $`a`$ and the metric $`\gamma _{ab}`$ are functions of the coordinates $`y^a`$ only. Here $`\mathrm{\Omega }_{ij}`$ is the metric of the maximally symmetric space of constant curvature given in (2). Given this structure of the background metric, we are forced to classify metric perturbations by their three-dimensional tensor properties as in the four-dimensional case. We stress again that this is a direct consequence of the brane-world nature of the theory that we are considering. With these remarks in mind, the higher-dimensional generalization of the perturbed metric (3) can be written in the form
$$ds^2=a^2\left\{\gamma _{ac}\left(\delta _b^c+2\varphi _b^c\right)dy^ady^b\left[\left(12\psi \right)\mathrm{\Omega }_{ij}+2E_{|ij}+2F_{(i|j)}+h_{ij}\right]dx^idx^j2W_{ai}dy^adx^i\right\}.$$
(22)
As in the four-dimensional case, $`F_i`$ and $`h_{ij}`$ have a vanishing divergence and, in addition, $`h_{ij}`$ is traceless. As before, the three-vectors $`W_{ai}`$ can be split as follows,
$$W_{ai}=B_{a|i}+S_{ai},$$
(23)
where $`S_{a}^{}{}_{}{}^{i}{}_{|i}{}^{}=0`$. Observe that the perturbed metric (22), defined in this way, is completely general. In fact, this can be easily seen by counting degrees of freedom. As an example, we can consider the simplest case of only one extra dimension setting $`y=y^5`$ and $`a,b,\mathrm{}=0,5`$. Then, the most general perturbed metric contains $`15`$ degrees of freedom, which are parameterized by the seven scalar perturbations ($`\varphi _0^0`$, $`\varphi _5^0`$, $`\varphi _5^5`$, $`\psi `$, $`E`$, $`B_a`$), six degrees of freedom from the vector perturbations ($`F_i`$, $`S_{ai}`$) and two degrees of freedom from the tensor perturbation $`h_{ij}`$. Of course, in counting the degrees of freedom originating from vector and tensor perturbations we have taken the constraints on these quantities into account.
Let us now return to the general case of $`d`$ additional dimensions and consider the coordinate transformations
$`x^\alpha \stackrel{~}{x^\alpha }=x^\alpha +\xi ^\alpha `$ (24)
with infinitesimal parameters $`\xi ^\alpha `$. In accordance with the above discussion, we split these parameters as $`(\xi ^\alpha )=(\xi ^a,\xi ^i)`$. We adopt the useful convention that indices of type $`a`$ ($`i`$) are lowered, raised and contracted using the metric $`\gamma _{ab}`$ ($`\mathrm{\Omega }_{ij}`$). Furthermore, we take the vertical bar to denote the covariant derivative with respect to $`\gamma _{ab}`$ or $`\mathrm{\Omega }_{ij}`$ depending on the index type. From the transformation law
$$\stackrel{~}{g}_{\alpha \beta }=g_{\alpha \beta }2_{(\alpha }\xi _{\beta )},$$
(25)
of the metric, and taking into account that $`\xi _\alpha =a^2(\xi _a,\xi _i)`$, we find for the transformation of the scalar perturbations
$`\delta \varphi _{ab}`$ $`=`$ $`\xi _{(a|b)}H^c\xi _c\gamma _{ab},`$ (26)
$`\delta \psi `$ $`=`$ $`H^a\xi _a,`$ (27)
$`\delta E`$ $`=`$ $`\xi ,`$ (28)
$`\delta B_a`$ $`=`$ $`\xi _a\xi _{|a}.`$ (29)
where we have introduced the generalized Hubble parameters
$$H_c=\frac{a_{|c}}{a}.$$
(30)
The vector perturbations in the metric (22) change according to
$`\delta F_i`$ $`=`$ $`\eta _i,`$ (31)
$`\delta S_{ai}`$ $`=`$ $`\eta _{i|a}.`$ (32)
Finally, the tensor perturbation $`h_{ij}`$ is invariant under the first order gauge transformation (25). From these results we easily find the following gauge-invariant quantities.
Scalar variables
$`\mathrm{\Phi }_{ab}`$ $`=`$ $`\varphi _{ab}+H^c(B_cE_{|c})\gamma _{ab}+(B_{(a}E_{|(a})_{|b)}`$ (33)
$`\mathrm{\Psi }`$ $`=`$ $`\psi H^c\left(B_cE_{|c}\right).`$ (34)
Vector variables
$`_{ai}=S_{ai}F_{i|a}.`$ (35)
Tensor variables
$$h_{ij}.$$
(36)
We conclude that the physical degrees of freedom consist of the $`(d+1)(d+2)/2+1`$ gauge invariant scalar perturbations $`(\mathrm{\Phi }_{ab},\mathrm{\Psi })`$, $`d+1`$ gauge invariant vector perturbations $`_{ai}`$ and a gauge invariant tensor perturbation $`h_{ij}`$. The above gauge-invariant variables are a direct generalization of the corresponding four-dimensional ones. Specifically, restricting to no additional dimensions and setting $`\gamma _{00}=1`$, eqs. (33)–(36) exactly reproduce the four-dimensional expressions (14)–(17). However, in the case $`d>0`$ our formalism clearly has a richer structure than the conventional four-dimensional one.
### C A generalized longitudinal gauge for scalar perturbations
In the subsequent sections we will focus on the evolution of scalar perturbations. Vector and tensor perturbations will be discussed elsewhere. In order to simplify this discussion we introduce a generalized longitudinal gauge for the scalar perturbations. In analogy with the four-dimensional case, this gauge is specified by
$$B_a=0,E=0.$$
(37)
Setting these quantities to zero can indeed be achieved by an appropriate choice of the scalar transformation parameters $`\xi _a`$ and $`\xi `$ in the eqs. (26)–(29). Note that we have exactly the correct number of transformation parameters to do this and that, consequently, the gauge ambiguity in the scalar sector is complete eliminated by this choice. Then, the scalar part of the metric takes the simple form
$$ds^2=a^2\left\{\gamma _{ac}(\delta _b^c+2\varphi _b^c)dy^ady^b(12\psi )\mathrm{\Omega }_{ij}dx^idx^j\right\}.$$
(38)
Furthermore, in this gauge, the scalar perturbations $`\varphi _{ab}`$ and $`\psi `$ coincide with their gauge-invariant counterparts, that is
$$\mathrm{\Phi }_{ab}=\varphi _{ab},\mathrm{\Psi }=\psi $$
(39)
as it is the case in four dimensions.
### D The five-dimensional case
Let us restrict in this section and for the rest of the paper to the case of a single extra dimension $`y=y^5`$. Then, the indices $`a,b,\mathrm{}`$ run over the values $`0,5`$ only. Furthermore, in order to be more explicit, we choose the conformal gauge
$$(\gamma _{ab})=b^2\text{diag}(1,1)$$
(40)
for the background metric $`\gamma _{ab}`$ by performing a large gauge transformation. Here $`b=b(t,y)`$ is a new, independent scale factor. Then, the perturbed five-dimensional metric (22) reduces to
$`ds^2`$ $`=`$ $`a^2\{b^2[(1+2\varphi )dt^22Wdtdy(12\mathrm{\Gamma })dy^2]`$ (42)
$`[\mathrm{\Omega }_{ij}(12\psi )+2E_{|ij}+2F_{(i|j)}+h_{ij}]dx^idx^j2W_{0i}dtdx^i2W_{5i}dydx^i\},`$
where we have defined
$$\varphi =\varphi _0^0,\mathrm{\Gamma }=\varphi _5^5,W=2\varphi _0^5=2\varphi _5^0.$$
(43)
Recall that the scale factors $`a`$ and $`b`$ are functions of the coordinates $`t`$ and $`y`$ only while the perturbations depend on all spacetime coordinates. The scalar gauge-invariant variables defined in eq. (33) can now be written more explicitly as
$`\mathrm{\Phi }_1`$ $``$ $`\mathrm{\Phi }_0^0=\varphi +{\displaystyle \frac{1}{b^2}}\left[(H_0_0)(B_0\dot{E})+\dot{B}_0\ddot{E}(H_5+_5)(B_5E^{})\right]`$ (44)
$`\mathrm{\Phi }_2`$ $``$ $`\mathrm{\Phi }_5^5=\mathrm{\Gamma }{\displaystyle \frac{1}{b^2}}\left[(H_5_5)(B_5E^{})+B_5^{}E^{\prime \prime }(H_0+_0)(B_0\dot{E})\right]`$ (45)
$`\mathrm{\Phi }_3`$ $``$ $`\mathrm{\Phi }_0^5={\displaystyle \frac{W}{2}}{\displaystyle \frac{1}{2b^2}}\left[(B_0^{}+\dot{B}_5)2\dot{E}^{}2_5(B_0\dot{E})2_0(B_5E^{})\right]`$ (46)
$`\mathrm{\Phi }_4`$ $``$ $`\mathrm{\Psi }=\psi {\displaystyle \frac{1}{b^2}}\left[H_0(B_0\dot{E})H_5(B_5E^{})\right].`$ (47)
Here and in the following the dot (prime) denotes the derivative with respect to time (the coordinate $`y`$). Furthermore, we have introduced a second set of “Hubble”–parameters $`_a=b_{|a}/b`$. Let us specialize these results to the generalized longitudinal gauge defined by $`B_0=B_5=E=0`$. Then the above scalar gauge-invariant variables coincide with $`\varphi `$, $`\psi `$, $`\mathrm{\Gamma }`$ and $`W`$. The metric simplifies to
$$ds^2=a^2\left\{b^2\left[(1+2\varphi )dt^22Wdydt(12\mathrm{\Gamma })dy^2\right](12\psi )\mathrm{\Omega }_{ij}dx^idx^j\right\}.$$
(48)
This metric will be the starting point for our treatment of scalar perturbations in the following sections. In addition to the perturbations $`\varphi `$ and $`\psi `$ that we are familiar with from the four-dimensional case it contains two new perturbations, $`\mathrm{\Gamma }`$ and $`W`$, that are related to the presence of the fifth dimension.
## III The perturbed Einstein equation in the longitudinal gauge
As we have previously mentioned, the main application we have in mind for this paper is a compactification of a five-dimensional theory on the orbifold $`S_1/Z_2`$. We start by compactifying the fifth dimension on a circle restricting the corresponding coordinate $`y`$ to the range $`y[R,R]`$ with the endpoints identified. The action of the $`Z_2`$ orbifolding symmetry on the circle is taken to be $`yy`$. Consequently, there exist two fix points at $`y=y_1=0`$ and $`y=y_2=R`$. We assume that the three-branes, stretching across $`3+1`$–dimensional space-time, are located at these fix points in the orbifold direction. This setup is appropriate for a large class of five-dimensional heterotic M-theory models \[VI,VI,VI\] that originate from 11–dimensional Hořava-Witten theory. It also applies to the five-dimensional models introduced in ref. \[VI,VI\].
Next, we should truncate the five-dimensional metric in order to make it consistent with the orbifolding. Since the metric has to be intrinsically even under the $`Z_2`$ action its various components satisfy the constraints
$`g_{\mu \nu }(y)`$ $`=`$ $`g_{\mu \nu }(y),`$ (49)
$`g_{\mu 5}(y)`$ $`=`$ $`g_{\mu 5}(y),`$ (50)
$`g_{55}(y)`$ $`=`$ $`g_{55}(y).`$ (51)
At the same time, we have to make sure that coordinate transformations do not lead out of the class of metrics defined this way. The parameter $`\xi ^\alpha `$ for an infinitesimal coordinate transformation has, therefore, to be restricted by
$`\xi ^\mu (y)`$ $`=`$ $`\xi ^\mu (y),`$ (52)
$`\xi ^5(y)`$ $`=`$ $`\xi ^5(y),`$ (53)
which directly follows from (25). From these rules we can deduce the $`Z_2`$ properties of the various quantities in metric (48) for scalar perturbations. While the background scale factors $`a`$, $`b`$ as well as the perturbations $`\varphi `$, $`\psi `$ and $`\mathrm{\Gamma }`$ are $`Z_2`$ even, that is, for example, $`a(y)=a(y)`$, the perturbation $`W`$ is $`Z_2`$ odd, that is $`W(y)=W(y)`$. Similarly, for the scalar components in the transformation parameter $`\xi ^\alpha `$, we find that $`\xi _0`$ and $`\xi `$ are even while $`\xi _5`$ is odd. Also note that the derivative along the fifth dimension of an odd variable is even and vice versa. For instance, $`W^{}(y)=W^{}(y)`$. Based on these $`Z_2`$ truncations we should now discuss the continuity properties of all quantities. Normally, one requires the metric to be continuous across the whole of space-time in order to have a sensible notion of length and time. We will also adopt this viewpoint, however with an additional subtlety. Since the orbifolding identifies the upper and lower half of the circle in the fifth dimension already one of them, say the upper half, constitutes the whole of space-time. In fact, instead of working with the orbifold picture where one keeps the full circle as we do here, one could also use the boundary picture where only one half of the circle (a line-segment) is considered. This shows that a jump of a metric component at an orbifold fix point does not contradict the continuity requirement. Of course, such a jump is possible only for an odd component of the metric. Concretely, we therefore require that all components of the metric (48) are continuous across the full orbicircle except for the odd component $`W`$ which may jump at the fix points (but is continuous otherwise). Corresponding assumptions have to be made for the parameter $`\xi ^\alpha `$ so that coordinate transformations do not change these continuity properties of the metric. Clearly the even components $`\xi _0`$ and $`\xi `$ have to be continuous then. Is the odd component $`\xi _5`$ allowed to jump at the orbifold points? Eq. (26) shows that $`\mathrm{\Gamma }=\varphi _5^5`$, which has to be even and continuous, transforms with the derivative $`\xi _5^{}`$. Hence, if $`\xi _5`$ jumped at the fix points it would lead to a delta-function singularity in the metric which is clearly unacceptable. We, therefore, have to require that $`\xi _5`$ is continuous everywhere on the orbicircleFrom this conclusion we see, that we have glossed over a subtlety when introducing the generalized longitudinal gauge. Clearly, for continous $`\xi _5`$ the quantity $`B_5`$ can only be gauged to zero if $`E_{|5}B_5`$ is continous, as can be seen from (28) and (29). We will, therefore, in addition require the continuity of $`E_{|5}B_5`$.. In particular, this means that $`\xi _5`$ vanishes at the fix points, that is $`\xi _5(y_n)=0`$.
It is clear that the above conclusions depend somewhat on the fact that we are working with an orbifold. For example, if we had considered compactification on a circle instead, all components of the metric had to be continuous. Correspondingly, some of the conclusions below will be slightly modified for other compactifications, however, in a way that is usually rather obvious.
Given this setup the Einstein equations can be written as
$$G_{\alpha \beta }R_{\alpha \beta }\frac{1}{2}g_{\alpha \beta }R=T_{\alpha \beta }+\underset{n=1}{\overset{2}{}}T_{\alpha \beta }^{(n)}\delta (yy_n),$$
(54)
where we have set the five-dimensional Newton constant to one, for simplicity. The delta-functions in this equation are covariant with respect to the fifth dimension, that is, they include a factor of $`1/\sqrt{g_{55}}`$. Furthermore, $`T_{\alpha \beta }`$ is the bulk stress-energy tensor induced by fields that propagate in the full five-dimensional space time. The brane stress-energy tensors $`T_{\alpha \beta }^{(n)}`$, on the other hand, originate from fields that are confined to the branes at the orbifold fix points.
In order to proceed further, we need to specify these stress-energy tensors. Two requirements should be taken into account when doing this. First, one should use the fact that the background has a maximally symmetric three-dimensional space. Secondly, the brane stress-energy tensors should be restricted in a way that is consistent with their nature of representing fields on the branes. This latter requirement can be most easily implemented by using the constraints that follow from the Einstein equation (54) itself. Concretely, the delta-functions on the right-hand side of this equation have to be matched by corresponding delta-functions that appear on the left-hand side. The appearance of these latter delta-functions, however, is controlled by the structure of the equations and the continuity assumptions about the metric discussed above.
Let us see what this implies in detail. We start with the background stress-energy. For the bulk, the most general form of this tensor consistent with the three-dimensional maximal symmetry is
$$T_{}^{\alpha }{}_{\beta }{}^{}=\left(\begin{array}{ccc}\rho & 0& r\\ 0& p\delta _{}^{i}{}_{j}{}^{}& 0\\ r& 0& q\end{array}\right).$$
(55)
In particular, we note that the $`05`$ component can be non-vanishing. This possibility is, in fact, already realized for the simple case of a bulk scalar field that depends on $`t`$ and $`y`$. As far as the symmetry of the background metric is concerned, the background brane stress-energy tensors should have the same structure as (55). However, as we will see in a moment, there are two more requirements that follow from the equations of motion, namely that the $`55`$ and the $`05`$ components vanish. As a result, the background stress-energy on the branes has the form
$$T_{}^{(n)\alpha }{}_{\beta }{}^{}=\left(\begin{array}{ccc}\rho ^{(n)}& 0& 0\\ 0& p^{(n)}\delta _{}^{i}{}_{j}{}^{}& 0\\ 0& 0& 0\end{array}\right).$$
(56)
Let us now proceed to the perturbed stress-energy tensors. Since we are focusing on scalar perturbations we write the most general perturbation of the background bulk tensor (55) that can be expressed in terms of scalars on the maximally symmetric subspace. This leads to
$$\delta T_{}^{\alpha }{}_{\beta }{}^{}=\left(\begin{array}{ccc}\delta \rho & (\rho +p)b^2v_{|j}& \delta r\\ (\rho +p)v^{|i}& \delta p\delta _{}^{i}{}_{j}{}^{}+\sigma _{}^{|i}{}_{|j}{}^{}& u^{|i}\\ \delta r+2r(\varphi +\mathrm{\Gamma })(\rho +q)W& b^2u_{|j}& \delta q\end{array}\right)$$
(57)
where $`v`$ and $`u`$ are two potentials for “velocity” fields and $`\sigma `$, satisfying $`\sigma _{}^{|i}{}_{|i}{}^{}=0`$, determines the anisotropic stress. The perturbed brane stress-energy tensors should have the same structure. However, as we will see below, the equations of motion impose further constraints implying vanishing $`55`$ and $`5i`$ components as well as vanishing anisotropic stress. Therefore, the brane stress-energy perturbations are given by
$$\delta T_{}^{(n)\alpha }{}_{\beta }{}^{}=\left(\begin{array}{ccc}\delta \rho ^{(n)}& (\rho ^{(n)}+p^{(n)})b^2v_{|j}^{(n)}& \delta r^{(n)}\\ (\rho ^{(n)}+p^{(n)})v^{(n)|i}& \delta p^{(n)}\delta _{}^{i}{}_{j}{}^{}& 0\\ \delta r^{(n)}\rho ^{(n)}W& 0& 0\end{array}\right).$$
(58)
We would like to present the equations of motion based on the metric (48) and on the above stress-energy tensors that follow from the Einstein equation (54). However, in writing the metric (48) two gauge choices where involved and it is not, a priori, clear that these choices can be made while, at the same time, keeping the branes at $`y=\text{const}`$ hypersurfaces as we have conveniently assumed in writing eq. (54). As the first choice, we decided to work with the two-dimensional metric $`\gamma _{ab}`$ in conformal gauge. Fortunately, this can be achieved while keeping the branes at hypersurfaces $`y=\text{const}`$ \[VI\]. In addition, we have chosen the generalized longitudinal gauge (37) for the scalar perturbations. A brane, described by $`y=y_n`$ before the gauge transformation that leads to longitudinal gauge, is described by $`y=\stackrel{~}{y}\xi ^5(y)=y_n`$ after this gauge transformation, where $`\stackrel{~}{y}`$ is the transformed $`y`$ coordinate. However, since $`\xi ^5(y_n)=0`$, as discussed above, this equation is solved by $`\stackrel{~}{y}=y_n`$ to linear order. We conclude that, in the new coordinates that lead to the generalized longitudinal gauge, the brane location is unchanged to the relevant linear order. In summary, therefore, using the Einstein equation (54) with the branes described by $`y=y_n`$ does not restrict the generality of our results given the gauge choices that we have made.
The background equations following from (54) have been given in ref. \[VI\] for the case of stress-energy induced by scalar fields and in ref. \[VI\] for the case of ideal fluids. For completeness and in order to incorporate some of the generalizations that we have made (such as the inclusion of a $`05`$ component of the bulk stress-energy) we will nevertheless present these equations here. We find
$`a^2b^2G_{}^{0}{}_{0}{}^{}`$ $``$ $`3\left[2{\displaystyle \frac{\dot{a}^2}{a^2}}+{\displaystyle \frac{\dot{a}\dot{b}}{ab}}{\displaystyle \frac{a^{\prime \prime }}{a}}+{\displaystyle \frac{a^{}b^{}}{ab}}+kb^2\right]=a^2b^2\left[\rho +{\displaystyle \underset{n=1}{\overset{2}{}}}\rho ^{(n)}\overline{\delta }(yy_n)\right]`$ (59)
$`a^2b^2G_{}^{5}{}_{5}{}^{}`$ $``$ $`3\left[{\displaystyle \frac{\ddot{a}}{a}}{\displaystyle \frac{\dot{a}\dot{b}}{ab}}2{\displaystyle \frac{a_{}^{}{}_{}{}^{2}}{a^2}}{\displaystyle \frac{a^{}b^{}}{ab}}+kb^2\right]=a^2b^2q`$ (60)
$`a^2b^2G_{}^{0}{}_{5}{}^{}`$ $``$ $`3\left[{\displaystyle \frac{\dot{a}^{}}{a}}+2{\displaystyle \frac{\dot{a}a^{}}{a^2}}+{\displaystyle \frac{\dot{a}b^{}}{ab}}+{\displaystyle \frac{a^{}\dot{b}}{ab}}\right]=a^2b^2r`$ (61)
$`a^2b^2G_{}^{i}{}_{j}{}^{}`$ $``$ $`\left[3{\displaystyle \frac{\ddot{a}}{a}}+{\displaystyle \frac{\ddot{b}}{b}}{\displaystyle \frac{\dot{b}^2}{b^2}}3{\displaystyle \frac{a^{\prime \prime }}{a}}{\displaystyle \frac{b^{\prime \prime }}{b}}+{\displaystyle \frac{b_{}^{}{}_{}{}^{2}}{b^2}}+kb^2\right]\delta _{}^{i}{}_{j}{}^{}=a^2b^2\left[p+{\displaystyle \underset{n=1}{\overset{2}{}}}p^{(n)}\overline{\delta }(yy_n)\right]\delta _{}^{i}{}_{j}{}^{}`$ (62)
Here we have defined the delta-function $`\overline{\delta }`$ which incorporates a factor $`1/ab`$. Based on these equations we can now justify the vanishing of the $`55`$ and $`05`$ components in the brane stress-energy (56). Such components, if non-vanishing, would appear on the right-hand sides of (60) and (61) multiplied with delta-functions. We should, therefore, have corresponding delta-function terms on the left-hand sides of these equations. Since the scale factors $`a`$ and $`b`$ are assumed to be continuous, delta functions can only appear from second derivatives of these quantities with respect to $`y`$. However, there are no such terms in eq. (60) and (61). Hence, we conclude that the $`55`$ and $`05`$ components in eq. (56) must vanish.
For the perturbations, we find to linear order
$`(ab)^2\delta G_0^0`$ $``$ $`3\left[2{\displaystyle \frac{a^{}b^{}}{ab}}2{\displaystyle \frac{a^{\prime \prime }}{a}}{\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{}{y}}{\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{}{t}}\right]\mathrm{\Gamma }3\left[{\displaystyle \frac{\dot{a}^{}}{a}}+2{\displaystyle \frac{a^{}\dot{a}}{a^2}}+{\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{}{y}}\right]W6\left[2{\displaystyle \frac{\dot{a}^2}{a^2}}+{\displaystyle \frac{\dot{a}\dot{b}}{ab}}\right]\varphi `$ (64)
$`+3\left[3{\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{}{y}}{\displaystyle \frac{b^{}}{b}}{\displaystyle \frac{}{y}}3{\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{}{t}}{\displaystyle \frac{\dot{b}}{b}}{\displaystyle \frac{}{t}}+2kb^2\right]\psi +b^2\left(2\psi +\mathrm{\Gamma }\right)_{|i}^{|i}+3\psi ^{\prime \prime }`$
$`=`$ $`a^2b^2\left\{\delta \rho +{\displaystyle \underset{n=1}{\overset{2}{}}}(\delta \rho ^{(n)}+\mathrm{\Gamma }\rho ^{(n)})\overline{\delta }(yy_n)\right\}`$ (65)
$`(ab)^2\delta G_5^5`$ $``$ $`6\left[2{\displaystyle \frac{a_{}^{}{}_{}{}^{2}}{a^2}}+{\displaystyle \frac{a^{}b^{}}{ab}}\right]\mathrm{\Gamma }3\left[{\displaystyle \frac{\dot{a^{}}}{a}}+2{\displaystyle \frac{a^{}\dot{a}}{a^2}}+{\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{}{t}}\right]W+3\left[2{\displaystyle \frac{\dot{a}\dot{b}}{ab}}2{\displaystyle \frac{\ddot{a}}{a}}{\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{}{y}}{\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{}{t}}\right]\varphi `$ (67)
$`+3\left[3{\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{}{y}}+{\displaystyle \frac{b^{}}{b}}{\displaystyle \frac{}{y}}3{\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{}{t}}+{\displaystyle \frac{\dot{b}}{b}}{\displaystyle \frac{}{t}}+2kb^2\right]\psi +b^2\left(2\psi \varphi \right)_{|i}^{|i}3\ddot{\psi }`$
$`=`$ $`a^2b^2\delta q`$ (68)
$`(ab)^2\delta G_5^0`$ $``$ $`3\left[{\displaystyle \frac{a^{\prime \prime }}{a}}2{\displaystyle \frac{a^{}b^{}}{ab}}2{\displaystyle \frac{a^2}{a^2}}\right]W+3\left[2{\displaystyle \frac{\dot{a}^{}}{a}}2{\displaystyle \frac{\dot{a}b^{}}{ab}}2{\displaystyle \frac{a^{}\dot{b}}{ab}}4{\displaystyle \frac{a^{}\dot{a}}{a^2}}+{\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{}{y}}\right]\varphi `$ (70)
$`+3\left[{\displaystyle \frac{^2}{ty}}{\displaystyle \frac{b^{}}{b}}{\displaystyle \frac{}{t}}{\displaystyle \frac{\dot{b}}{b}}{\displaystyle \frac{}{y}}\right]\psi {\displaystyle \frac{b^2}{2}}W_{}^{|i}{}_{|i}{}^{}3{\displaystyle \frac{a^{}}{a}}\dot{\mathrm{\Gamma }}`$
$`=`$ $`a^2b^2\left\{\delta r{\displaystyle \underset{n=1}{\overset{2}{}}}\delta r^{(n)}\overline{\delta }(yy_n)\right\}`$ (71)
$`(ab)^2\delta G_i^0`$ $``$ $`\left\{\left[{\displaystyle \frac{3}{2}}{\displaystyle \frac{a^{}}{a}}+{\displaystyle \frac{b^{}}{b}}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{y}}\right]W+\left[3{\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{\dot{b}}{b}}\right]\varphi +\left[{\displaystyle \frac{\dot{b}}{b}}+{\displaystyle \frac{}{t}}\right]\mathrm{\Gamma }+2\dot{\psi }\right\}_{|i}`$ (72)
$`=`$ $`a^2\left\{(\rho +p)v{\displaystyle \underset{n=1}{\overset{2}{}}}(\rho ^{(n)}+p^{(n)})v^{(n)}\overline{\delta }(yy_n)\right\}_{|i}`$ (73)
$`(ab)^2\delta G_i^5`$ $``$ $`\left\{\left[3{\displaystyle \frac{a^{}}{a}}+{\displaystyle \frac{b^{}}{b}}\right]\mathrm{\Gamma }+\left[{\displaystyle \frac{b^{}}{b}}+{\displaystyle \frac{}{y}}\right]\varphi +\left[{\displaystyle \frac{3}{2}}{\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{\dot{b}}{b}}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{t}}\right]W2\psi ^{}\right\}_{|i}`$ (74)
$`=`$ $`a^2u_{|i}`$ (75)
$`(ab)^2\delta G_j^i`$ $``$ $`\{[6{\displaystyle \frac{a^{\prime \prime }}{a}}2{\displaystyle \frac{b^{\prime \prime }}{b}}+2{\displaystyle \frac{b^2}{b^2}}3{\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{}{y}}3{\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{}{t}}{\displaystyle \frac{\dot{b}}{b}}{\displaystyle \frac{}{t}}{\displaystyle \frac{b^{}}{b}}{\displaystyle \frac{}{y}}{\displaystyle \frac{^2}{t^2}}]\mathrm{\Gamma }`$ (80)
$`+\left[2{\displaystyle \frac{b^{}\dot{b}}{b^2}}2{\displaystyle \frac{\dot{b}^{}}{b}}6{\displaystyle \frac{\dot{a}^{}}{a}}3{\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{}{t}}3{\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{}{y}}{\displaystyle \frac{\dot{b}}{b}}{\displaystyle \frac{}{y}}{\displaystyle \frac{b^{}}{b}}{\displaystyle \frac{}{t}}{\displaystyle \frac{^2}{ty}}\right]W`$
$`+\left[2{\displaystyle \frac{\dot{b}^2}{b^2}}2{\displaystyle \frac{\ddot{b}}{b}}6{\displaystyle \frac{\ddot{a}}{a}}3{\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{}{y}}3{\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{}{t}}{\displaystyle \frac{\dot{b}}{b}}{\displaystyle \frac{}{t}}{\displaystyle \frac{b^{}}{b}}{\displaystyle \frac{}{y}}{\displaystyle \frac{^2}{y^2}}\right]\varphi `$
$`+[6{\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{}{y}}6{\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{}{t}}+2{\displaystyle \frac{^2}{y^2}}2{\displaystyle \frac{^2}{t^2}}]\psi +2b^2(\psi \varphi +\mathrm{\Gamma })_{|k}^{|k}\}\delta ^i_j`$
$`b^2(\psi \varphi +\mathrm{\Gamma })_{|j}^{|i}`$
$`=`$ $`a^2b^2\left\{\delta p\delta _{}^{i}{}_{j}{}^{}+\sigma _{}^{|i}{}_{|j}{}^{}{\displaystyle \underset{n=1}{\overset{2}{}}}(\delta p^{(n)}+\mathrm{\Gamma }p^{(n)})\delta _{}^{i}{}_{j}{}^{}\overline{\delta }(yy_n)\right\}`$ (81)
Given those results, we can now return to the question of why the perturbations of the brane stress-energy tensors must have the specific form (58). We recall that all quantities in our metric (48) are even except the off-diagonal perturbation $`W`$ which is odd under the $`Z_2`$ symmetry. From our continuity assumptions, delta-function terms in the perturbed Einstein tensor can, therefore, arise from first derivatives of $`W`$ with respect to $`y`$ and second derivatives with respect to $`y`$ of all other quantities. Inspection of the above equations shows that such terms are absent in the $`5i`$ and $`55`$ components of the perturbed Einstein tensor. Consequently, the corresponding components in the perturbed brane stress-energy should vanish. Furthermore, all terms in $`\delta G_{}^{i}{}_{j}{}^{}`$ that could potentially lead to delta-functions are proportional to $`\delta _{}^{i}{}_{j}{}^{}`$. This implies that the anisotropic stress on the brane, which would contribute to the traceless part of eq. (80), must vanish. As a result, the traceless part of eq. (80)
$$(\psi \varphi +\mathrm{\Gamma })_{|j}^{|i}\frac{1}{3}(\psi \varphi +\mathrm{\Gamma })_{|k}^{|k}\delta _{}^{i}{}_{j}{}^{}=a^2\sigma _{}^{|i}{}_{|j}{}^{}$$
(82)
only involves the bulk anisotropic stress as a source term. If the bulk anisotropic stress vanishes as well, as, for example, is the case for a perfect fluid in the bulk, one concludes that
$$\psi \varphi +\mathrm{\Gamma }=0.$$
(83)
The quantity $`\psi \varphi +\mathrm{\Gamma }`$ is the analog of the four-dimensional quantity $`\psi _4\varphi _4`$ that also vanishes in the absence of anisotropic stress. This correspondence will be made more explicit in section V.
## IV Density fluctuations on the brane
A systematic study of density fluctuation in five dimensions requires solving the full set of five-dimensional equations of motion given in the previous section. However, for specific questions it might be useful to have some information about the metric restricted to the brane. For example, it is this restricted metric that is felt by matter which is confined to the brane. In this section, we are going to derive such equations on the brane starting from the general equations of motion above.
For a $`Z_2`$ even field the meaning of its value on the brane is quite clear. A $`Z_2`$ odd field may jump across the brane so its value may have a sign ambiguity. To simplify the notation, we define the value of an odd field on the brane as the one that is approached from within the interval $`y[0,R]`$. This is precisely the boundary value of the field as viewed in the boundary picture and, at the same time represents one half of the jump at the fix point. We also recall that the scale factors $`a`$, $`b`$ and the perturbations $`\varphi `$, $`\psi `$ and $`\mathrm{\Gamma }`$ are even and hence continuous while the perturbation $`W`$ is odd and may jump across the fix points.
Let us start with the background equations of motion \[VI,VI\]. We have already mentioned that the delta-function sources in eq. (59) and (62) have to be matched by the terms containing second $`y`$ derivatives of the scale factors $`a`$ and $`b`$. This leads to
$$\frac{a^{}}{a}=\frac{1}{6}ab\rho ^{(n)},\frac{b^{}}{b}=\pm \frac{1}{2}ab\left(\rho ^{(n)}+p^{(n)}\right).$$
(84)
These conditions, as well as the following ones, are valid at the brane positions $`y=y_n`$ where the upper (lower) sign holds for the brane $`n=1`$ ($`n=2`$). While the two other non-vanishing equations of motion do not contain delta functions, they can still be restricted to the brane. From the $`05`$ component (61) we find
$$\dot{\rho }^{(n)}=3\frac{\dot{a}}{a}\left(\rho ^{(n)}+p^{(n)}\right)2abr.$$
(85)
which represents to an energy conservation equation on the brane. Note, however, that, in addition to intrinsic brane quantities, this equation also involves the off-diagonal bulk stress-energy component $`r`$. This reflects the simple fact that the branes are not isolated systems but can exchange energy with the bulk. Finally, we should consider the $`55`$ component (60). Restricted to the branes it results in an equation of motion for the values of the scale factors $`a`$ and $`b`$ on the brane given by
$$\frac{\ddot{a}}{a}\frac{\dot{a}\dot{b}}{ab}+kb^2=\frac{a^2b^2}{3}\left[\frac{1}{12}\rho ^{(n)}\left(\rho ^{(n)}+3p^{(n)}\right)+q\right].$$
(86)
An analogous procedure can now be applied to the perturbed equations. Observe that only the components $`\delta G_0^0`$, $`\delta G_i^i`$, $`\delta G_5^0`$ and $`\delta G_i^0`$ contain explicit delta-function terms. They should be matched by terms containing first $`y`$ derivatives of $`W`$ and second $`y`$ derivatives of all other quantities. This leads to
$`\psi ^{}`$ $`=`$ $`{\displaystyle \frac{\dot{a}}{a}}W\pm {\displaystyle \frac{1}{6}}ab\left(\delta \rho ^{(n)}\mathrm{\Gamma }\rho ^{(n)}\right)`$ (87)
$`\varphi ^{}`$ $`=`$ $`\left({\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{\dot{b}}{b}}+{\displaystyle \frac{}{t}}\right)W\pm {\displaystyle \frac{1}{3}}ab\left(\delta \rho ^{(n)}\mathrm{\Gamma }\rho ^{(n)}\right)\pm {\displaystyle \frac{1}{2}}ab\left(\delta p^{(n)}\mathrm{\Gamma }p^{(n)}\right)`$ (88)
$`W`$ $`=`$ $`{\displaystyle \frac{a}{b}}\left(\rho ^{(n)}+p^{(n)}\right)v^{(n)}`$ (89)
$`W`$ $`=`$ $`{\displaystyle \frac{\delta r^{(n)}}{\rho ^{(n)}}}`$ (90)
Interestingly, the last equation implies that the component $`\delta T_{}^{(n)5}{}_{0}{}^{}`$ of the brane stress-energy perturbation vanishes, as can be seen by comparison with eq. (58). The component $`\delta T_{}^{(n)0}{}_{5}{}^{}=\delta r^{(n)}`$, however, is non-zero and is, from eq. (89), (90) determined by
$$\delta r^{(n)}=\frac{a}{b}\rho ^{(n)}\left(\rho ^{(n)}+p^{(n)}\right)v^{(n)}.$$
(91)
We have, therefore, found an important additional constraint on the perturbed brane energy-momentum tensor (58). The quantity $`\delta r^{(n)}`$ is, in fact, uniquely fixed by the other components. For vacuum energy with $`p^{(n)}=\rho ^{(n)}`$ on the branes $`\delta r^{(n)}`$ is zero, but it is generally non-vanishing otherwise. This is, perhaps, somewhat surprising since one could have expected that a purely four-dimensional stress-energy tensor on the brane (with all $`5`$ components vanishing) should be allowed. Here we see that this is generally not the case.
Next, we deal with the odd components $`05`$ and $`5i`$ of the perturbed equations of motion given in (70) and (74). Restriction to the branes leads us, after some algebra, to
$`\dot{\delta }^{(n)}`$ $`=`$ $`\left(1+w^{(n)}\right)\left(v_{}^{(n)|i}{}_{|i}{}^{}3\dot{\psi }\right)3{\displaystyle \frac{\dot{a}}{a}}\left({\displaystyle \frac{\delta p^{(n)}}{\delta \rho ^{(n)}}}w^{(n)}\right)\delta ^{(n)}`$ (93)
$`2(1+w^{(n)})v^{(n)}a^2(\rho +q)2ab\left(\mathrm{\Gamma }+2\varphi \delta ^{(n)}+{\displaystyle \frac{\delta r}{r}}\right){\displaystyle \frac{r}{\rho ^{(n)}}}`$
and
$`{\displaystyle \frac{\dot{v}_{}^{(n)}{}_{|i}{}^{}}{b^2}}`$ $`=`$ $`\left[{\displaystyle \frac{\dot{a}}{a}}\left(13w^{(n)}\right){\displaystyle \frac{\dot{w}^{(n)}}{1+w^{(n)}}}+{\displaystyle \frac{\dot{b}}{b}}\right]{\displaystyle \frac{v_{}^{(n)}{}_{|i}{}^{}}{b^2}}{\displaystyle \frac{\delta p_{}^{(n)}{}_{|i}{}^{}}{\delta \rho ^{(n)}}}{\displaystyle \frac{\delta ^{(n)}}{1+w^{(n)}}}\varphi _{|i}`$ (95)
$`2{\displaystyle \frac{a}{b}}{\displaystyle \frac{1}{\rho ^{(n)}}}\left[{\displaystyle \frac{u}{1+w^{(n)}}}rv^{(n)}\right]_{|i},`$
where $`w^{(n)}=p^{(n)}/\rho ^{(n)}`$ and $`\delta ^{(n)}=\delta \rho ^{(n)}/\rho ^{(n)}`$ denotes the energy contrast on the branes. These equations represent the conservation of energy and momentum for scalar perturbations, including possible energy-momentum flow from the bulk onto the brane or from the brane into the bulk. They should be compared with the corresponding equations in four dimensions, eq. (30) in \[VI\], taking into account that the variable $`\theta `$ of \[VI\] is related to the peculiar velocity $`v`$ via $`\theta =v_{}^{|i}{}_{|i}{}^{}`$. Eq. (93) differs from the four-dimensional result by the last two terms. They describe the energy flux from the bulk onto the brane. In eq. (95) we have two new terms with respect to the four-dimensional equation. They describe momentum flux between bulk and brane. The coupling between bulk gravity and brane matter expressed via the above equations is one of the main results of this paper. It shows that, when considering scalar metric perturbations on the branes, the branes cannot simply be viewed as an isolated system but have to be considered together with the bulk environment. Practically, this implies that frequently one cannot simply copy four-dimensional formulae when dealing with physics on a brane that is embedded in a higher-dimensional space. Finally, we restrict the $`55`$ component of the equations of motion, eq. (67), to the brane. We find the following evolution equation for the perturbations projected onto the branes:
$`b^2(2\psi \varphi )_{|i}^{|i}`$ $``$ $`3\ddot{\psi }3{\displaystyle \frac{\dot{a}}{a}}\dot{\varphi }+3({\displaystyle \frac{\dot{b}}{b}}3{\displaystyle \frac{\dot{a}}{a}})\dot{\psi }+6kb^2(\psi +\varphi )+a^2b^2\rho _{}^{(n)}{}_{}{}^{2}[{\displaystyle \frac{1}{6}}(1+3w^{(n)})\varphi `$ (96)
$`+`$ $`{\displaystyle \frac{\delta q+2q\varphi }{\rho _{}^{(n)}{}_{}{}^{2}}}+{\displaystyle \frac{1}{6}}(1+{\displaystyle \frac{3}{2}}w^{(n)})\delta ^{(n)}+{\displaystyle \frac{\delta p^{(n)}}{4\delta \rho ^{(n)}}}\delta ^{(n)}\pm {\displaystyle \frac{a}{b}}{\displaystyle \frac{r}{\rho ^{(n)}}}(1+w^{(n)})v^{(n)}]=0.`$ (97)
## V Matching to the four-dimensional effective theory
In the previous subsection, we have derived a set of equations for the metric on the branes, essentially by restricting the five-dimensional equations of motion. These results may, for example, be useful to analyze the evolution of matter that is confined to the brane. However, the most important task is to extract predictions for structure in the late universe from our formalism of metric perturbations in brane-world theories. In this section, we will explain the basic steps in this direction.
First, we should introduce the four-dimensional effective theory, describing physics at low energy, that is associated to our five-dimensional brane-world theory (54). It is this four-dimensional theory that described the evolution of the universe “today” and that is used for the interpretation of observational results. Theoretical predictions, originating from our brane-world theory, should therefore be formulated in terms of this effective theory. The five- and the four-dimensional effective theory are related by a vacuum state that constitutes a specific solution of the five-dimensional theory and should respect the symmetries that we expect to find in the four-dimensional theory. Specifically, four-dimensional Lorentz invariance implies that the vacuum metric should have the structure
$$ds^2=A^2(y)dx^\mu dx^\nu \eta _{\mu \nu }B^2(y)dy^2.$$
(98)
The functions $`A`$ and $`B`$ have to be such that this metric solves the five-dimensional theory in the vacuum configuration. For our five-dimensional theory (54), the simplest possibility is to have no stress energy in the vacuum which results in a flat vacuum metric, $`A,B=\text{const}`$. In five-dimensional heterotic M-theory the vacuum configuration is determined by certain potentials in the bulk and on the branes that involve the dilaton \[VI,VI\]. In this case, $`A`$ and $`B`$ are non-trivial functions of $`y`$ and the deviation from the flat vacuum metric is determined by the size of the so called strong coupling expansion parameter. The vacua proposed in ref. \[VI,VI\] are based on a vacuum configuration with specific cosmological constants in the bulk and on the branes and result in an exponential function for $`A`$ in the coordinate system where $`B=\text{const}`$. Each one of these different vacuum states is associated with its specific low-energy theory. For the sake of simplicity and concreteness, we will here focus on the first possibility, namely the flat vacuum. This choice represents, at the same time, a good approximation for five-dimensional heterotic M-theory in the case of a small strong-coupling expansion parameter. The four-dimensional effective theory describes the dynamics of the collective excitations of the vacuum state. In our case, these excitation are given by a four-dimensional metric $`g_{4\mu \nu }`$ and the modulus $`\beta `$ describing the size of the fifth dimension. The vacuum metric with these collective modes put in has the structure
$$d\overline{s}^2=e^\beta g_{4\mu \nu }dx^\mu dx^\nu e^{2\beta }dy^2,$$
(99)
where $`g_{4\mu \nu }`$ and $`\beta `$ are functions of $`x^\mu `$. As usual, the effective four-dimensional description is valid as long as these functions are varying slowly enough. This is the case if all four-dimensional momenta are much smaller than the mass of the first Kaluza-Klein excitation around the vacuum state. In our case this mass is given by $`e^\beta /2R`$. Let us, therefore, consider a five-dimensional evolution that is approaching the vacuum state (99). Even though the five-dimensional metric is then close to the vacuum metric it will still have small Kaluza-Klein excitations that can be described in linear perturbation theory. A useful way to extract the zero modes from such a five-dimensional metric with small contributions from Kaluza-Klein modes is to perform an average over the fifth dimension. Doing this systematically leads to the following four-dimensional effective theory associated to the brane-world theory (54) and the vacuum state (99):
$`R_{4\mu \nu }{\displaystyle \frac{1}{2}}g_{4\mu \nu }R_4`$ $`=`$ $`{\displaystyle \frac{3}{2}}\left(_\mu \beta _\nu \beta {\displaystyle \frac{1}{2}}g_{4\mu \nu }\beta ^2\right)+T_{4\mu \nu }`$ (100)
$`_4^2\beta `$ $`=`$ $`J_4.`$ (101)
Four- and five-dimensional quantities are related by
$`e^{2\beta }`$ $`=`$ $`<g_{55}>`$ (102)
$`g_{4\mu \nu }`$ $`=`$ $`e^\beta <g_{\mu \nu }>`$ (103)
$`T_{4\mu \nu }`$ $`=`$ $`<T_{\mu \nu }>+{\displaystyle \frac{1}{2Re^\beta }}{\displaystyle \underset{n=1}{\overset{2}{}}}T_{\mu \nu }^{(n)}`$ (104)
$`J_4`$ $`=`$ $`{\displaystyle \frac{2}{3}}e^{3\beta }<T_{55}>+{\displaystyle \frac{1}{3}}g_4^{\mu \nu }T_{4\mu \nu }.`$ (105)
Here, $`<>`$ denotes the average over the fifth dimension. As stated, this four-dimensional theory is a good description as long as all momenta are small compared to $`e^\beta /2R`$, the mass of the first Kaluza-Klein mode. The Kaluza-Klein modes have decoupled from the above equations at linear order. However, due to the presence of the branes, the Kaluza-Klein modes cannot strictly be set to zero but have to be integrated out \[VI,VI\]. This leads to higher order corrections to the above four-dimensional equations that are suppressed by powers of the four-dimensional Planck scale and that we have neglected. To be consistent with this approximation, the average $`<>`$ that projects out the Kaluza-Klein excitations should be considered meaningful only at the linearized level in Kaluza-Klein excitations.
We would now like to apply the above general correspondence to our formalism for metric fluctuations. To do this, we need to assume a five-dimensional background solution that, for late time, approaches the vacuum configuration. Formulated in a four-dimensional language, this is the case if the four-dimensional Hubble parameter $`H_4=\dot{a}_4/a_4`$ and $`\dot{\beta }`$ are small compared to $`e^\beta /2R`$. Furthermore, it is helpful to assume that the average $`<\gamma _{00}>`$ approaches one in this limit. This can always be achieved by a redefinition of time. We would like to explicitly work out the correspondence for the scalar sector in longitudinal gauge that we have focused on in this paper. The generalization to include vector and tensor perturbations is straightforward. Concretely, we apply the correspondence (102)–(105) to the five-dimensional quantities (48), (55), (56), (57) and (58) matching onto the four-dimensional quantities specified in (18), (19) and (20). Furthermore, we need to decompose the four-dimensional modulus $`\beta `$ as
$$\beta =\chi \mathrm{\Gamma }_4.$$
(106)
where $`\chi =\chi (t)`$ is the time-dependent background and $`\mathrm{\Gamma }_4=\mathrm{\Gamma }_4(t,x^i)`$ is the perturbation. The matching of background quantities leads to
$`e^{2\chi }`$ $`=`$ $`<a^2b^2>`$ (107)
$`a_4^2`$ $`=`$ $`e^\chi <a^2>`$ (108)
$`\rho _4`$ $`=`$ $`e^\chi <\rho >+{\displaystyle \frac{1}{2Re^{2\chi }}}{\displaystyle \underset{n=1}{\overset{2}{}}}\rho ^{(n)}`$ (109)
$`p_4`$ $`=`$ $`e^\chi <p>+{\displaystyle \frac{1}{2Re^{2\chi }}}{\displaystyle \underset{n=1}{\overset{2}{}}}p^{(n)}.`$ (110)
It is interesting to explicitly compute the background current $`J_4`$ in the equation of motion (101) for the modulus $`\chi `$. It is given by
$$J_4=\frac{1}{3}\left[\rho _43p_4+2e^\chi <q>\right].$$
(111)
With the above expression for $`\rho _4`$ and $`p_4`$, this implies that the modulus $`\chi `$ has a runaway potential leading to a growing size of the fifth dimension. Therefore, the theory, as stands, will not stay in the range of validity of the four-dimensional effective theory. As is well-known, it needs further stabilization of the modulus $`\beta `$ by means of a potential. In the context of string- or M-theory one expects such a potential to be generated by non-perturbative effects.
The correspondence for the scalar perturbations reads
$`\mathrm{\Gamma }_4`$ $`=`$ $`<\mathrm{\Gamma }>`$ (112)
$`\varphi _4`$ $`=`$ $`<\varphi \mathrm{\Gamma }/2>`$ (113)
$`\psi _4`$ $`=`$ $`<\psi +\mathrm{\Gamma }/2>`$ (114)
$`\delta \rho _4`$ $`=`$ $`e^\chi <\delta \rho +\mathrm{\Gamma }\rho >+{\displaystyle \frac{1}{2Re^{2\chi }}}{\displaystyle \underset{n=1}{\overset{2}{}}}(\delta \rho ^{(n)}+2\mathrm{\Gamma }_4\rho ^{(n)})`$ (115)
$`\delta p_4`$ $`=`$ $`e^\chi <\delta p+\mathrm{\Gamma }p>+{\displaystyle \frac{1}{2Re^{2\chi }}}{\displaystyle \underset{n=1}{\overset{2}{}}}(\delta p^{(n)}+2\mathrm{\Gamma }_4p^{(n)})`$ (116)
$`\sigma _4`$ $`=`$ $`e^\chi <\sigma >`$ (117)
$`v_4`$ $`=`$ $`{\displaystyle \frac{e^{3\chi }a_4^4}{\rho _4+p_4}}\left[<(\rho +p)v>+{\displaystyle \frac{1}{2Re^\beta }}{\displaystyle \underset{n=1}{\overset{2}{}}}(\rho ^{(n)}+p^{(n)})v^{(n)}\right].`$ (118)
In particular, we conclude that
$$\psi _4\varphi _4=<\psi \varphi +\mathrm{\Gamma }>$$
(119)
Hence, the four- and five-dimensional quantities that measure the presence of anisotropic stress are in direct correspondence to one another as they should.
## VI Conclusion
In this paper we have laid down a gauge-invariant formalism to describe metric fluctuations in brane-world theories. This formalism is a straightforward generalization of the well known formalism in four dimensions. It categorizes the perturbations according to their tensor properties with respect to the usual three-dimensional maximally symmetric space rather than a higher-dimensional space as one might have expected. This is a direct consequence of the brane-world nature of the theory which generally leads to cosmological backgrounds that are inhomogeneous in the additional dimensions. We have introduced a generalized longitudinal gauge in order to further study scalar perturbations. In the case of a five-dimensional model on the orbifold $`S^1/Z_2`$, on which we have focused, we have identified four scalar metric perturbations $`\varphi `$, $`\psi `$, $`\mathrm{\Gamma }`$ and $`W`$. This has to be contrasted to the four-dimensional case where one only has two such perturbations. We have presented the evolution equations for these scalar perturbations which, mainly due to the dependence of the background on the additional dimension, are significantly more complicated than the corresponding four-dimensional equations. It is those additional terms, related to the non-linearity of the background in the additional coordinates, that encode possible new and interesting information about the formation and evolution of perturbations. Furthermore, given the gauge choices and assumptions about the continuity of the metric that we have made, we have determined the resulting most general form of the stress–energy on the brane. In particular, we have found that the perturbed brane stress-energy has to have vanishing anisotropic stress and that its $`05`$ component is non-zero. We have applied our formalism to calculate the matching conditions (Israel conditions) for the five-dimensional metric restricted to the branes. Among other results we have derived the equations describing energy and momentum conservation for metric perturbations on the brane. As is expected on physical grounds, they illustrate that the brane cannot be viewed as an isolated object but is subject to energy and momentum flow between the bulk and the brane. Finally, we have shown how the five-dimensional formalism for metric fluctuations can be matched to the known four-dimensional one in the limit where the brane-world theory has an effective four-dimensional description. This allows one to extract predictions for structure in the late universe originating from brane-world theories. We hope to address this problem more explicitly in a future publication.
As this manuscript was prepared for submission, ref. \[VI\], \[VI\] and \[VI\] appeared which have some overlap with the present paper. Cosmological perturbations are also discussed in \[VI\].
Acknowledgments:
C. v. d. Bruck is grateful to E. Eyras, J. Martin, C. Martins, H. Reall and T. Shiromizu, and in particular D. Langlois, R. Maartens and D. Wands for useful discussions. C. v. d. Bruck is supported by Nato/DAAD (at Brown) and Deutsche Forschungsgemeinschaft (DFG, at Cambridge). Miquel Dorca is supported by the Fundación Ramón Areces. The research was supported in part (at Brown) by the U.S. Department of Energy under Contract DE-FG02-91ER40688, TASK A. A. L. is supported by the European Community under contract No. FMRXCT 960090.
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# 1 Introduction
## 1 Introduction
The AdS/CFT correspondence is arguably the best currently available way of getting nontrivial dynamical information for the strong coupling behavior of certain conformal field theories. In particular, the $`𝒩=4`$ supersymmetric $`SU(N)`$ Yang-Mills theory in four dimensions (SYM<sub>4</sub>) at large $`N`$ and at strong ‘t Hooft coupling $`\lambda =g_{YM}^2N`$ is dual to type IIB supergravity on the $`AdS_5\times S^5`$ background. The supergravity fields are dual to certain quasi-primary operators in SYM<sub>4</sub>. According to , the generating functional for the connected Green functions of these operators coincides with the on-shell value of type IIB supergravity action which has to be further modified by the addition of definite boundary terms . Thus, computing $`n`$-point correlation functions in the supergravity approximation is generally divided into two independent problems - finding first the supergravity action up to the $`n`$-th order and then evaluating its on-shell value. Although a covariant action for type IIB supergravity is unknown, one can use the covariant equations of motion and the quadratic action to find cubic actions for its physical fields and to compute the corresponding 3-point functions using the technique developed in .
Computing 4-point functions - in the supergravity approximation in general requires the derivation of the supergravity action up to fourth order. The part of the action relevant to the massless modes, corresponding to the dilaton and axion fields, was already known as pointed out in where the calculation of the corresponding 4-point functions was initiated. The complete expression for the 4-point functions was obtained in and was further analyzed in . Unfortunately, these modes are dual to the rather complicated operators $`\mathrm{tr}F^2+\mathrm{}`$ and $`\mathrm{tr}F\stackrel{~}{F}`$ and the analysis performed in was unavoidably incomplete.
It is known that all operators dual to the type IIB supergravity fields belong to short representations of the conformal superalgebra $`SU(2,2|4)`$ and are supersymmetric descendents of Chiral Primary Operators (CPOs) of the form $`O_k^I=\mathrm{tr}(\varphi ^{(i_1}\mathrm{}\varphi ^{i_k)})`$. CPOs are dual to scalar fields $`s^I`$ that are mixtures of the five form field strength on $`S^5`$ and the trace of the graviton on $`S^5`$. The relevant part of the quartic action of type IIB supergravity for the scalars $`s^I`$ was found in and was then used in to compute the 4-point functions of the simplest CPOs $`O^I=\mathrm{tr}(\varphi ^{(i}\varphi ^{j)})`$. In the present paper we use these 4-point functions to analyze in detail the Operator Product Expansion (OPE) of the lowest weight CPOs at strong coupling.
It is widely believed that the structure of a Conformal Field Theory (CFT) is encoded in the OPE since knowledge of the latter allows, in principle, the calculation of all $`n`$-point functions. Thus, in the context of AdS/CFT correspondence one would eventually like to prove that the 4-point functions, (and in general $`n`$-point functions) of CPOs in the boundary CFT computed in the supergravity approximation admit an OPE interpretation. This is a rather complicated problem because an infinite number of quasi-primary operators may in principle appear in the OPE of two CPOs. Therefore, the best one can presently do is to show that the leading terms in a double OPE expansion of the 4-point functions exactly match the contributions of the conformal blocks of the first few quasi-primary operators with the lowest conformal dimensions. This is the main line of investigation which we follow in the present work in our analysis of the 4-point function of the lowest weight CPOs in $`𝒩=4`$ SYM<sub>4</sub>.
Our study shows that there are four singular terms in the OPE of two lowest weight CPOs corresponding to the identity operator, the lowest weight CPO itself, the $`R`$-symmetry vector current and the stress tensor. These three nontrivial operators are dual to the scalars $`s^I`$, the vector fields $`A_\mu `$ and the graviton $`h_{\mu \nu }`$ that appear in the exchange Feynman diagrams of type IIB supergravity. The most singular terms in the 4-point functions computed in the supergravity approximation exactly coincide with the contributions coming from the conformal blocks of the above three operators.
We compare the strong coupling OPE with the free field theory OPE, and explicitly observe, at weak coupling, the splitting of the $`R`$-symmetry current and of the stress tensor into 2 and 3 terms respectively which belong to different supermultiplets. Only one term in each splitting is dual to a supergravity field and survives at strong coupling while the others acquire large dimensions and decouple. A similar type of splitting also occurs in the case of the double-trace operators transforming in the $`\mathrm{𝟖𝟒}`$ and $`\mathrm{𝟏𝟕𝟓}`$ irreps.
We also analyze the leading nonsingular terms in the OPE which are due to double-trace operators of the schematic form $`:^mO^I^nO^J:`$ with free field conformal dimensions $`4+m+n`$. A generic property of any correlation function computed in the supergravity approximation is the appearance of logarithmic terms. In an unitary CFT logarithmic terms have a natural interpretation in terms of anomalous dimensions of operators and such an interpretation was used in the past in studies of the $`O(N)`$ vector model . Since the operators dual to the supergravity fields have protected conformal dimensions, the logarithmic terms in the correlation functions of supergravity can only be attributed to anomalous dimensions of double-trace operators.
We show that among the scalar double-trace operators with free field conformal dimension 4, the only one acquiring an anomalous dimension is the operator $`:O^IO^I:`$, which transforms in the trivial representation of the $`R`$-symmetry group $`SO(6)SU(4)`$. The anomalous dimension of this operator is found to be $`16/N^2`$ and coincides with the anomalous dimension of the operator $`B`$ which was calculated in . This is consistent with the fact that $`B`$ is a supersymmetric descendent of $`:O^IO^I:`$. It is worth noting that among the non-renormalized operators we find a double-trace scalar operator in the $`\mathrm{𝟐𝟎}`$ irrep of $`SO(6)`$ whose non-renormalization property does not follow from the shortening condition discussed in .
Finally, we compute the anomalous dimensions of the double-trace vector operators with free field conformal dimension 5 transforming in the $`\mathrm{𝟏𝟓}`$ and $`\mathrm{𝟏𝟕𝟓}`$ irreps respectively. We show that there are several towers of traceless symmetric tensor operators in the 105, 84 and 175 irreps, whose anomalous conformal dimensions vanish. Some of these tensor operators are not subject to any known non-renormalization theorem.
The 4-point functions of CPOs also allow us to find the leading $`1/N^2`$ corrections to the normalization constants of the 3-point functions involving two CPOs and one double-trace operator with low conformal dimension. In the case when a double trace operator has protected dimension we interpret these corrections as manifestation of the splitting of the free field theory operator in two orthogonal parts carrying different representations of supersymmetry. The first one has protected both the dimension and the normalization constant, the other one acquires infinite anomalous dimension and disappears at strong coupling. To make this interpretation precise one should further show that the linear splitting arising due to the difference between normalization constants in free theory and at strong coupling is consistent with the fact that the split fields transform in different representations of supersymmetry. It would be quite interesting to investigate such a property in more detail.
The plan of the paper is as follows. In section 2 we recall how logarithmic terms are related to anomalous conformal dimensions in an unitary CFT and in the framework of the AdS/CFT correspondence. In section 3 we discuss the structure of the OPE of the lowest weight CPOs in free field theory and at strong coupling. In section 4 we compute anomalous dimensions and first corrections to the 2- and 3-point normalization constants of double-trace operators of approximate dimensions 4 and 5. A discussion of the results obtained and our conclusions are presented in section 5. Several technical issues are considered in five Appendices. In the Appendix A we discuss a decomposition of a bi-local operator which is a normal-ordered product of two quasi-primary scalar operators into a sum of conformal blocks of local tensor primary operators. In the Appendix B explicit formulae for conformal partial amplitudes of scalar, conserved vector current and stress tensor are derived. A convenient series representation used throughout the paper is obtained in the Appendix C. In the Appendix D we discuss the projectors which single out the contributions of irreps occurring in the decomposition $`\mathrm{𝟐𝟎}\times \mathrm{𝟐𝟎}`$ of $`SO(6)`$ from the 4-point function of CPOs. In the Appendix E an explicit formula for the conformal block of the stress tensor is derived.
## 2 Anomalous dimensions and logarithmic terms in CFT
An arbitrary unitary CFT is completely characterized by a set of quasi-primary operators $`O_i`$ of conformal dimensions $`\mathrm{\Delta }_i`$ and by their OPE
$`O_i(x)O_j(y)={\displaystyle \underset{k}{}}{\displaystyle \frac{1}{|xy|^{\mathrm{\Delta }_i+\mathrm{\Delta }_j\mathrm{\Delta }_k}}}C_{ij}^k(xy,_y)O_k(y).`$ (2.1)
Here the sum runs over the set of all the quasi-primary operators and $`i,j,k`$ are multi-indices which in general include the indices of the $`R`$-symmetry and of the Lorentz groups. The operator algebra structure constants $`C_{ij}^k(xy,_y)`$ can be decomposed in a power series in $`xy`$ and $`_y`$. Without loss of generality one can assume that the operators $`O_i`$ are orthogonal
$`O_i(x)O_j(0)=C_i{\displaystyle \frac{\delta _{ij}}{x^{2\mathrm{\Delta }_i}}},`$
where $`C_i`$ is a normalization constant of the 2-point function. Then the operator algebra structure constants are fixed by the conformal dimensions $`\mathrm{\Delta }_i,\mathrm{\Delta }_j,\mathrm{\Delta }_k`$, and by the ratio $`C_{ijk}/C_k`$, where the structure constants $`C_{ijk}`$ appear in the 3-point functions
$`O_i(x)O_j(y)O_k(z)={\displaystyle \frac{C_{ijk}}{|xy|^{\mathrm{\Delta }_i+\mathrm{\Delta }_j\mathrm{\Delta }_k}|xz|^{\mathrm{\Delta }_i+\mathrm{\Delta }_k\mathrm{\Delta }_j}|yz|^{\mathrm{\Delta }_j+\mathrm{\Delta }_k\mathrm{\Delta }_i}}}.`$ (2.2)
The conformal dimensions and the structure constants depend on the coupling constants of the CFT. In principle, the OPE (2.1) allows one to compute any correlation function in the CFT. In particular, 4-point functions are given by the following (schematic) double OPE expansion
$`O_i(x)O_j(y)O_k(z)O_l(w)`$ $`=`$ $`{\displaystyle \underset{m}{}}{\displaystyle \frac{1}{|xy|^{\mathrm{\Delta }_i+\mathrm{\Delta }_j\mathrm{\Delta }_m}|zw|^{\mathrm{\Delta }_k+\mathrm{\Delta }_l\mathrm{\Delta }_m}}}`$ (2.3)
$`\times C_{ij}^m(xy,_y)C_{kl}^m(zw,_w){\displaystyle \frac{C_m}{|yw|^{2\mathrm{\Delta }_m}}}.`$
Thus we see that the short distance expansion of exact CFT correlation functions does not contain logarithmic terms. Suppose, however, that one can only calculate correlation functions up to some order in the coupling constant or another small parameter of the CFT. Then it is clear from (2.3) that logarithmic terms would appear due to the nontrivial dependence of conformal dimensions on the coupling or on the small parameter. These terms can be easily found representing the conformal dimensions as $`\mathrm{\Delta }=\mathrm{\Delta }^{(0)}+\mathrm{\Delta }^{(1)}`$, where $`\mathrm{\Delta }^{(0)}`$ is the “canonical” part and $`\mathrm{\Delta }^{(1)}`$ is the “anomalous” coupling constant dependent part. Such a representation leads then to an expansion for the two-point functions of the form $`|x|^{\mathrm{\Delta }^{(1)}}=1+\mathrm{\Delta }^{(1)}\mathrm{log}|x|+\mathrm{}`$, connecting the logarithmic terms to the anomalous dimensions, that may be used to compute the latter. It is worthwhile to note that at the $`n`$-th order of perturbation theory one encounters terms of the form $`(\mathrm{log}|x|)^n`$.
The $`𝒩=4`$ SYM<sub>4</sub> theory provides an example of such a logarithmic behavior of correlation functions, both in the weak coupling standard perturbation expansion - and also in the supergravity approximation . Due to superconformal invariance all quasi-primary operators of SYM<sub>4</sub> belong either to short or long representations of the conformal superalgebra $`SU(2,2|4)`$ and in the framework of the AdS/CFT correspondence fall into three classes:
i) Chiral operators dual to the type IIB supergravity fields which belong to short representations and have protected conformal dimensions. The simplest operators in this class are the lowest weight CPOs $`O^I=\mathrm{tr}(\varphi ^{(i}\varphi ^{j)})`$.
ii) Operators dual to multi-particle supergravity states which are obtained as “normal-ordered” products of the chiral operators, e.g. the double-trace operators $`:O^IO^J:`$. They may belong either to short or long representations and have conformal dimensions restricted from above.
iii) Operators dual to string states (single- or multi-particle) which belong to long representations and whose conformal dimensions grow as $`\lambda ^{1/4}`$ in the strong coupling limit. The simplest example of such an operator is the Konishi operator $`\mathrm{tr}(\varphi ^i\varphi ^i)`$.
In the supergravity approximation to the AdS/CFT correspondence the operators dual to string states decouple from the spectrum and one can calculate the connected $`n`$-point functions of chiral operators dual to the supergravity fields to leading order which is $`1/N^{n2}`$. Since the expansion parameter is $`1/N^2`$, an $`n`$-point function contains logarithmic terms of the form $`(\mathrm{log}|x|)^{[(n2)/2]}`$. In particular, a 4-point function can have only $`\mathrm{log}|x|`$-dependent terms, and cannot have, say, terms of the form $`(\mathrm{log}|x|)^2`$. Moreover, since chiral operators have protected conformal dimensions only the operators dual to multi-particle supergravity states contribute to $`\mathrm{log}`$-dependent terms.
The AdS/CFT correspondence predicts a simple form of the OPE of chiral operators in the strong coupling limit. Let $`O_1`$ and $`O_2`$ be operators dual to the supergravity fields $`\phi _1`$ and $`\phi _2`$ respectively and let the supergravity action contain the non-vanishing cubic couplings $`\frac{1}{N}\lambda _{12k}\phi _1\phi _2\phi _k`$ with some fields $`\phi _k`$. Then, the OPE of $`O_1`$ and $`O_2`$ takes the form (suppressing the indices of the operators and structure constants)
$`O_1(x)O_2(y)={\displaystyle \frac{1}{N}}{\displaystyle \underset{k}{}}{\displaystyle \frac{1}{|xy|^{\mathrm{\Delta }_1+\mathrm{\Delta }_2\mathrm{\Delta }_k}}}C_{12}^k(xy,_y)O_k(y)+[:O_1(x)O_2(y):],`$ (2.4)
where $`O_k`$ is an operator dual to $`\phi _k`$. Here we denote by $`[:O_1(x)O_2(y):]`$ an infinite sum of tensor quasi-primary operators and their descendents, which are dual to multi-particle supergravity states. In general these operators acquire anomalous dimensions and are responsible for the appearance of logarithms in correlation functions. An important property of the operators dual to multi-particle supergravity states is that their structure constants are of order 1, while the structure constants of the operators dual to supergravity fields are of order $`1/N`$. Due to such a property, the sum of these operators coincides in the limit $`N\mathrm{}`$ with the corresponding free field theory normal-ordered operator $`:O_1^{fr}(x)O_2^{fr}(y):`$. This can be seen as follows. A 4-point function of chiral operators is given by a sum of a disconnected contribution which is of order 1 and a connected Green function which is of order $`1/N^2`$. Since the structure constants of the operators dual to supergravity fields are of order $`1/N`$, they do not contribute to the disconnected part of the 4-point function. Thus only the “normal-ordered” operators contribute. The disconnected part is given by a sum of products of 2-point functions of chiral operators, hence it does not depend on the coupling constant and $`N`$ (we assume that all the chiral operators are orthonormal) and coincides with the free field disconnected part. Therefore, in the limit $`N\mathrm{}`$ the sum $`[:O_1(x)O_2(y):]`$ has to coincide with the free field normal-ordered product $`:O_1^{fr}(x)O_2^{fr}(y):`$,<sup>4</sup><sup>4</sup>4 One can easily see that the normal-ordered product $`:O_1^{fr}(x)O_2^{fr}(y):`$ is the only term of order 1 in the free field OPE of chiral operators. that is decomposed into a sum of local tensor quasi-primary operators. However, at finite $`N`$ an infinite number of the tensor operators acquire anomalous dimensions and their structure constants get $`1/N^2`$ corrections to their free field values. For this reason it seems hardly possible to prove that a 4-point function computed in the supergravity approximation admits an OPE interpretation. This would require the knowledge of the conformal partial wave amplitude of an arbitrary tensor operator. Another reason that complicates the analysis of 4-point functions is that in general one should split the free field theory double-trace operators into a sum of operators with the same free field theory dimensions, each one transforming irreducibly under the superconformal group. In the context of the present work we are able to successfully deal with both the above problems.
## 3 OPE of the lowest weight CPOs
In this section we study the OPE of the lowest weight CPOs in free field theory and at strong coupling. Recall that the normalized lowest weight CPOs in $`𝒩=4`$ SYM<sub>4</sub> are operators of the form
$$O^I(x)=\frac{2^{3/2}\pi ^2}{\lambda }C_{ij}^I\mathrm{tr}(:\varphi ^i\varphi ^j:),$$
where the symmetric traceless tensors $`C_{ij}^I`$, $`i,j=1,2,..,6`$ form a basis of the $`\mathrm{𝟐𝟎}`$ of $`SO(6)`$ and satisfy the orthonormality condition
$$C_{ij}^IC_{ij}^J=\delta ^{IJ}.$$
Using for the Wick contractions the following propagator
$`\varphi _a^i\varphi _b^j={\displaystyle \frac{g_{YM}^2\delta _{ab}\delta ^{ij}}{(2\pi )^2x_{12}^2}},`$ (3.1)
where $`a,b`$ are color indices and $`x_{ij}=x_ix_j`$, one finds the following expressions for the free field theory 2-, 3- and 4-point functions of $`O^I`$:
$`O^{I_1}(x_1)O^{I_2}(x_2)_{fr}={\displaystyle \frac{\delta ^{I_1I_2}}{x_{12}^2}},`$
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)_{fr}={\displaystyle \frac{1}{N}}{\displaystyle \frac{2^{3/2}C^{I_1I_2I_3}}{x_{12}^2x_{13}^2x_{23}^2}},`$
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)_{fr}=[{\displaystyle \frac{\delta ^{I_1I_2}\delta ^{I_3I_4}}{x_{12}^4x_{34}^4}}+{\displaystyle \frac{\delta ^{I_1I_3}\delta ^{I_2I_4}}{x_{13}^4x_{24}^4}}+{\displaystyle \frac{\delta ^{I_1I_4}\delta ^{I_2I_3}}{x_{14}^4x_{23}^4}}]`$
$`+{\displaystyle \frac{4}{N^2}}[{\displaystyle \frac{C^{I_1I_2I_3I_4}}{x_{12}^2x_{23}^2x_{34}^2x_{41}^2}}+{\displaystyle \frac{C^{I_1I_3I_2I_4}}{x_{13}^2x_{32}^2x_{24}^2x_{41}^2}}+{\displaystyle \frac{C^{I_1I_3I_4I_2}}{x_{13}^2x_{34}^2x_{42}^2x_{21}^2}}],`$ (3.2)
where the first term in the 4-point function represents the contribution of disconnected diagrams. We have also introduced the shorthand notations $`C^{I_1I_2I_3}=C_{i_1i_2}^{I_1}C_{i_2i_3}^{I_2}C_{i_3i_1}^{I_3}`$ and $`C^{I_1I_2I_3I_4}=C_{i_1i_2}^{I_1}C_{i_2i_3}^{I_2}C_{i_3i_4}^{I_3}C_{i_4i_1}^{I_4}`$ for the trace products of matrices $`C^I`$.
### 3.1 Free field theory OPE
The simplest way to derive the OPE in free field theory is to apply Wick’s theorem. Using the propagator (3.1) we find the following formula for the product of two CPOs
$`O^{I_1}(x_1)O^{I_2}(x_2)`$ $`=`$ $`{\displaystyle \frac{\delta ^{I_1I_2}}{x_{12}^4}}+{\displaystyle \frac{2^3\pi ^2}{\lambda Nx_{12}^2}}C_{ik}^{I_1}C_{kj}^{I_2}:\mathrm{tr}(\varphi ^i(x_1)\varphi ^j(x_2)):`$ (3.3)
$`+`$ $`:O^{I_1}(x_1)O^{I_2}(x_2):`$
On the r.h.s. of (3.3) we have bi-local operators of the form $`:O^\alpha (x_1)O^\beta (x_2):`$, where $`O^\alpha `$ is either $`\varphi ^i`$ or $`O^{I_1}`$ and $`O^\beta `$ is either $`\varphi ^j`$ or $`O^{I_2}`$. To find the operator content of the r.h.s. of (3.3) one should perform the Taylor expansion of the operator $`O^\alpha `$ and rearrange the resulting series as a sum of conformal blocks of local quasi-primary operators. It is clear that in free field theory any bilocal operator $`:O^\alpha (x_1)O^\beta (x_2):`$ may be represented as an infinite sum of conformal blocks of symmetric traceless rank $`l`$ tensor operators with dimensions $`\mathrm{\Delta }_\alpha +\mathrm{\Delta }_\beta +l+2k`$,
$`:O^\alpha (x)O^\beta (0):={\displaystyle \underset{l,k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(l+2k)!}}x^{2k}x_1^\mu \mathrm{}x_l^\mu [O_{\mu _1\mathrm{}\mu _l}^{(k)}(0)],`$ (3.4)
where the square brackets $`[]`$ are used to denote the whole conformal block of a quasi-primary operator. In an interacting theory the tensor quasi-primary operators may acquire anomalous dimensions. Explicit expressions of the tensor operators through $`O^\alpha ,O^\beta `$ are unknown and the best we can do is to find the first few terms in the series. In particular, as shown in Appendix A, the terms up to two derivatives are given by the following formula
$`:O^\alpha (x)O^\beta (0):`$ $`=`$ $`:O^\alpha (0)O^\beta (0):+x^\mu :_\mu O^\alpha (0)O^\beta (0):+{\displaystyle \frac{1}{2}}x^\mu x^\nu :_\mu _\nu O^\alpha (0)O^\beta (0):`$ (3.5)
$`=`$ $`[O^{\alpha \beta }(0)]+x^\mu [O_\mu ^{\alpha \beta }(0)]{\displaystyle \frac{1}{2}}x^\mu x^\nu [T_{\mu \nu }^{\alpha \beta }(0)]+{\displaystyle \frac{1}{2}}x^2[T^{\alpha \beta }(0)].`$ (3.6)
Here the quasi-primary operators are given by
$`O^{\alpha \beta }`$ $`=`$ $`:O^\alpha O^\beta :,`$ (3.7)
$`O_\mu ^{\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{1}{2}}:(_\mu O^\alpha O^\beta O^\alpha _\mu O^\beta ):,`$
$`T_{\mu \nu }^{\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{1}{2}}:(_\mu O^\alpha _\nu O^\beta +_\nu O^\alpha _\mu O^\beta ):{\displaystyle \frac{\mathrm{\Delta }}{2(2\mathrm{\Delta }+1)}}_\mu _\nu (:O^\alpha O^\beta :)`$
$`+`$ $`{\displaystyle \frac{\delta _{\mu \nu }}{8}}({\displaystyle \frac{\mathrm{\Delta }+1}{2\mathrm{\Delta }+1}}^2(:O^\alpha O^\beta :)+:^2O^\alpha O^\beta :+:O^\alpha ^2O^\beta :),`$
$`T^{\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{1}{8}}({\displaystyle \frac{\mathrm{\Delta }1}{2\mathrm{\Delta }1}}^2(:O^\alpha O^\beta :)+:^2O^\alpha O^\beta :+:O^\alpha ^2O^\beta :),`$
where $`\mathrm{\Delta }`$ is the conformal dimension of the operators $`O^\alpha ,O^\beta `$ which takes the values 1 and 2 in the cases under consideration.
Obviously the conformal dimensions of the scalar operators $`O^{\alpha \beta }`$ and $`T^{\alpha \beta }`$ are equal to $`2\mathrm{\Delta }`$ and $`2\mathrm{\Delta }+2`$ respectively, the dimension of the vector operator is $`2\mathrm{\Delta }+1`$ and the dimension of the traceless symmetric tensor operator is $`2\mathrm{\Delta }+2`$. Consider first the case when $`\mathrm{\Delta }=1`$. The scalar operator $`\mathrm{tr}(\varphi ^i\varphi ^j)`$ is decomposed into a sum of the traceless part in the 20 \- which is a lowest weight CPO $`O^I`$ \- and the trace part. The trace part is the normalized Konishi scalar field $`𝒦=\frac{2\pi ^2}{3^{1/2}\lambda }\mathrm{tr}(\varphi ^2)`$. If $`\mathrm{\Delta }=1`$ the vector and tensor operators are conserved and the operator $`T^{ij}`$ vanishes because of the on-shell equation $`^2\varphi ^i=0`$. In fact the conserved current transforms in the $`\mathrm{𝟏𝟓}`$ irrep of $`SO(6)`$ and is the $`R`$-symmetry current of the free field theory of 6 scalars $`\varphi ^i`$. Decomposing the tensor operator $`T_{\mu \nu }^{ij}`$ into irreducible representations of the $`R`$-symmetry group $`SO(6)`$, i.e. into the traceless and trace parts with respect to the indices $`i,j`$, one sees that the trace part $`T_{\mu \nu }^{ii}`$ coincides with the stress tensor of the free field theory. The Konishi scalar and the traceless part of $`T_{\mu \nu }^{ij}`$ are dual to string modes and are expected to decouple in the strong coupling limit.
To complete the consideration of the free field theory OPE we have to decompose the remaining operators into irreducible representations of $`SO(6)SU(4)`$. One has the general decomposition of the $`\mathrm{𝟐𝟎}\times \mathrm{𝟐𝟎}`$ of $`SU(4)`$ as
$`\mathrm{𝟐𝟎}\times \mathrm{𝟐𝟎}`$ $`=`$ $`[0,0,0]+[0,2,0]+[0,4,0]+[2,0,2]`$ (3.8)
$`+`$ $`[1,0,1]+[1,2,1].`$
The representations in the first and the second lines of (3.8) are symmetric and antisymmetric in the indices of the $`\mathrm{𝟐𝟎}`$’s $`I_1,I_2`$, respectively. The dimensions of the representations are
$`D([0,0,0])=1,D([0,2,0])=20,D([0,4,0])=105,D([2,0,2])=84,`$
$`D([1,0,1])=15,D([1,2,1])=175.`$ (3.9)
Introducing the orthonormal Clebsh-Gordon coefficients $`C_{𝒥_D}^{I_1I_2}`$
$$C_{𝒥_D}^{I_1I_2}C_{𝒥_D^{}}^{I_1I_2}=\delta _{𝒥_D𝒥_D^{}},$$
where $`𝒥_D`$ is the index of an irrep of dimension $`D`$, as well as the operators
$`O^{𝒥_D}=C_{𝒥_D}^{I_1I_2}:O^{I_1}O^{I_2}:,O_\mu ^{𝒥_D}=C_{𝒥_D}^{I_1I_2}O_\mu ^{I_1I_2},`$ (3.10)
we can write
$`:O^{I_1}O^{I_2}:`$ $`=`$ $`\delta ^{I_1I_2}O_1+C_{𝒥_{20}}^{I_1I_2}O^{𝒥_{20}}+C_{𝒥_{105}}^{I_1I_2}O^{𝒥_{105}}+C_{𝒥_{84}}^{I_1I_2}O^{𝒥_{84}}`$
$`O_\mu ^{I_1I_2}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(:_\mu O^{I_1}O^{I_2}::O^{I_1}_\mu O^{I_2}:)=C_{𝒥_{15}}^{I_1I_2}O_\mu ^{𝒥_{15}}+C_{𝒥_{175}}^{I_1I_2}O_\mu ^{𝒥_{175}},`$
and a similar decomposition for $`T_{\mu \nu }^{I_1I_2}`$ and $`T^{I_1I_2}`$. Note that the operators have the following free field theory 2-point functions<sup>5</sup><sup>5</sup>5The only exception is the operator $`O_1=\frac{1}{20}:O^{I_1}O^{I_2}:`$ in the singlet representation, whose normalization constant is $`\frac{1}{10}+O(\frac{1}{N^2})`$.
$`O^{𝒥_1}(x_1)O^{𝒥_2}(x_2)=(2+O({\displaystyle \frac{1}{N^2}})){\displaystyle \frac{\delta ^{𝒥_1𝒥_2}}{x_{12}^8}},`$
$`O_\mu ^{𝒥_1}(x_1)O_\nu ^{𝒥_2}(x_2)=(4+O({\displaystyle \frac{1}{N^2}})){\displaystyle \frac{I_{\mu \nu }(x_{12})}{x_{12}^{10}}}\delta ^{𝒥_1𝒥_2},`$
where $`I_{\mu \nu }(x)=\delta _{\mu \nu }2\frac{x_\mu x_\nu }{x^2}`$. The precise values of the normalization constants will be determined in the next section. Due to the definition of the double-trace operators, the 3-point normalization constants which appear in the following 3-point functions
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{𝒥_D}(x_3)`$ $`=`$ $`C_{OOO_𝐃}{\displaystyle \frac{C_{𝒥_D}^{I_1I_2}}{|x_{12}|^{4\mathrm{\Delta }_D}|x_{13}|^{\mathrm{\Delta }_D}|x_{23}|^{\mathrm{\Delta }_D}}},`$
$`O^{I_1}(x_1)O^{I_2}(x_2)O_\mu ^{𝒥_D}(x_3)`$ $`=`$ $`C_{OOO_𝐃}{\displaystyle \frac{C_{𝒥_D}^{I_1I_2}\left(x_{23}^2x_{31}^\mu x_{31}^2x_{23}^\mu \right)}{|x_{12}|^{5\mathrm{\Delta }_D}|x_{13}|^{\mathrm{\Delta }_D+1}|x_{23}|^{\mathrm{\Delta }_D+1}}}`$
are equal to the 2-point normalization constants $`C_{O_𝐃}`$.
Combining all pieces together we obtain the first few terms in the free field OPE of the CPOs as
$`O^{I_1}(x_1)O^{I_2}(x_2)`$ $`=`$ $`{\displaystyle \frac{\delta ^{I_1I_2}}{x_{12}^4}}+{\displaystyle \frac{2^{3/2}}{N}}C^{I_1I_2I}{\displaystyle \frac{1}{x_{12}^2}}[O^I]+{\displaystyle \frac{2}{3^{1/2}N}}\delta ^{I_1I_2}{\displaystyle \frac{1}{x_{12}^2}}[𝒦]`$ (3.11)
$`+{\displaystyle \frac{2^{7/2}\pi ^2}{\lambda N}}{\displaystyle \frac{x_{12}^\mu }{x_{12}^2}}C_{𝒥_{15}}^{I_1I_2}[J_\mu ^{𝒥_{15}}]{\displaystyle \frac{2\pi ^2\delta ^{I_1I_2}}{3\lambda N}}{\displaystyle \frac{x_{12}^\mu x_{12}^\nu }{x_{12}^2}}[T_{\mu \nu }^{fr}]+{\displaystyle \frac{4\pi ^2}{\lambda N}}{\displaystyle \frac{x_{12}^\mu x_{12}^\nu }{x_{12}^2}}C^{I_1I_2I}[T_{\mu \nu }^I]`$
$`+\delta ^{I_1I_2}[O_1]+C_{𝒥_{20}}^{I_1I_2}[O^{𝒥_{20}}]+C_{𝒥_{105}}^{I_1I_2}[O^{𝒥_{105}}]+C_{𝒥_{84}}^{I_1I_2}[O^{𝒥_{84}}]`$
$`+C_{𝒥_{15}}^{I_1I_2}x_{12}^\mu [O_\mu ^{𝒥_{15}}]+C_{𝒥_{175}}^{I_1I_2}x_{12}^\mu [O_\mu ^{𝒥_{175}}]+\mathrm{}.`$
Here $`T_{\mu \nu }^{fr}`$ is the stress tensor of the free field theory of six scalar fields, while the normalized $`R`$-symmetry current $`J_\mu ^{𝒥_{15}}`$ is defined as follows
$$J_\mu ^{𝒥_{15}}=C_{ij}^{𝒥_{15}}\frac{1}{2}\mathrm{tr}(:_\mu \varphi ^i\varphi ^j::\varphi ^i_\mu \varphi ^j:),$$
where the antisymmetric tensors $`C_{ij}^{𝒥_{15}}`$ form a basis of the $`\mathrm{𝟏𝟓}`$ of $`SO(6)`$ and satisfy the orthogonality condition $`C_{ij}^{𝒥_{15}}C_{ij}^{𝒥_{15}^{}}=\delta ^{𝒥_{15}𝒥_{15}^{}}`$. The $`R`$-symmetry current has the following 2-point function
$$J_\mu ^{𝒥_{15}}(x_1)J_\nu ^{𝒥_{15}^{}}(x_2)=\frac{\lambda ^2}{8\pi ^4}\delta ^{𝒥_{15}𝒥_{15}^{}}\frac{I_{\mu \nu }(x_{12})}{x_{12}^6}.$$
We would like to stress that in addition to the above fields the OPE contains infinite towers of both single-trace as well as double-trace operators.
### 3.2 Strong coupling OPE
As was explained in the previous section, the strong coupling OPE of CPOs is easily determined from the cubic terms in the scalars $`s^I`$ dual to the lowest weight CPOs in the type IIB supergravity action. There are three different cubic vertices in the action describing the cubic couplings among the three scalars $`s^I`$, the interaction of the scalars with the graviton and the interaction with the $`SO(6)`$ vector fields. Thus, according to the discussion in the previous section the strong coupling OPE has the form
$`O^{I_1}(x_1)O^{I_2}(x_2)`$ $`=`$ $`{\displaystyle \frac{\delta ^{I_1I_2}}{x_{12}^4}}+{\displaystyle \frac{2^{3/2}}{N}}C^{I_1I_2I}{\displaystyle \frac{1}{x_{12}^2}}[O^I]+{\displaystyle \frac{2^{7/2}\pi ^2}{3\lambda N}}{\displaystyle \frac{x_{12}^\mu }{x_{12}^2}}C_{𝒥_{15}}^{I_1I_2}[R_\mu ^{𝒥_{15}}]`$ (3.12)
$``$ $`{\displaystyle \frac{2\pi ^2}{15\lambda N}}\delta ^{I_1I_2}{\displaystyle \frac{x_{12}^\mu x_{12}^\nu }{x_{12}^2}}[T_{\mu \nu }]+\delta ^{I_1I_2}x_{12}^{\mathrm{\Delta }_1^{(1)}}[O_1]`$
$`+`$ $`C_{𝒥_{20}}^{I_1I_2}x_{12}^{\mathrm{\Delta }_{20}^{(1)}}[O^{𝒥_{20}}]+C_{𝒥_{105}}^{I_1I_2}x_{12}^{\mathrm{\Delta }_{105}^{(1)}}[O^{𝒥_{105}}]+C_{𝒥_{84}}^{I_1I_2}x_{12}^{\mathrm{\Delta }_{84}^{(1)}}[O^{𝒥_{84}}]`$
$`+`$ $`C_{𝒥_{15}}^{I_1I_2}x_{12}^{\mathrm{\Delta }_{15}^{(1)}}x_{12}^\mu [O_\mu ^{𝒥_{15}}]+C_{𝒥_{175}}^{I_1I_2}x_{12}^{\mathrm{\Delta }_{175}^{(1)}}x_{12}^\mu [O_\mu ^{𝒥_{175}}]+\mathrm{}.`$
Here $`R_\mu ^{𝒥_{15}}`$ is the R-symmetry current and $`T_{\mu \nu }`$ is the stress tensor of $`𝒩=4`$ SYM<sub>4</sub>. The structure constants of the operators $`O^I`$, $`R_\mu ^{𝒥_{15}}`$, $`T_{\mu \nu }`$ are found by requiring that the above OPE reproduces the known 3-point functions of two CPOs with another CPO, the $`R`$-symmetry current and the stress tensor respectively, as the latter were computed in the supergravity approximation in . The operator algebra structure constants of the double-trace operators in (3.12) are chosen to be 1, which means that their 2- and 3-point normalization constants are kept equal. The anomalous dimensions $`\mathrm{\Delta }_1,\mathrm{\Delta }_{20},\mathrm{},\mathrm{\Delta }_{175}`$ of the double-trace operators will be determined in the next section by studying the 4-point functions of the CPOs.
Comparing (3.12) with (3.11), we see that the structure of the strong coupling OPE is simpler than the corresponding free field theory one. Instead of having an infinite number of single-trace operators as in (3.11), we find in (3.12) only three single-trace operators giving rise to the most singular terms. The coefficients in front of the $`R`$-symmetry current and the stress tensor are, however, different from the ones in (3.11). The reason is that the free field operators $`J_\mu ^{𝒥_{15}}`$ and $`T_{\mu \nu }^{fr}`$ receiving contribution only from bosons may be represented as
$`J_\mu ^{𝒥_{15}}={\displaystyle \frac{1}{3}}R_\mu ^{𝒥_{15}}+{\displaystyle \frac{2}{3}}𝒦_\mu ^{𝒥_{15}};T_{\mu \nu }^{fr}={\displaystyle \frac{1}{5}}T_{\mu \nu }+{\displaystyle \frac{10}{35}}𝒦_{\mu \nu }+{\displaystyle \frac{18}{35}}\mathrm{\Xi }_{\mu \nu },`$ (3.13)
where $`𝒦_\mu ^{𝒥_{15}}`$ and $`𝒦_{\mu \nu }`$ are vector and tensor operators from the Konishi supermultiplet which has as leading component that scalar $`𝒦`$, while $`\mathrm{\Xi }_{\mu \nu }`$ is the leading component of a new supersymmetry multiplet. The splitting (3.13) is explained by the fact that $`T_{\mu \nu }`$, $`𝒦_{\mu \nu }`$ and $`\mathrm{\Xi }_{\mu \nu }`$ have pairwise vanishing two-point functions and belong to different supersymmetry multiplets. The operators in the Konishi supermultiplet as well as $`\mathrm{\Xi }_{\mu \nu }`$ are dual to string modes and therefore decouple in the strong coupling limit.
A splitting analogous to (3.13) may also occur for the free field theory double-trace operators. However, there is an important difference. If we assume that all operators have free field theory 2-point normalization constants of order 1, then the splitting has the following schematic form
$$O^{fr}=O^{gr}+\frac{1}{N}O^{str},$$
where a free field theory double-trace operator $`O^{fr}`$ is split into a sum of operators $`O^{gr}`$ dual to supergravity multi-particle states, and operators $`O^{str}`$ dual to string states. As follows from the discussion in the previous section the coefficient in front of $`O^{str}`$ has to be of order $`1/N`$, because otherwise one would not reproduce the disconnected part of the 4-point function. Such a splitting manifests itself in the $`1/N^2`$ corrections to 2- and 3-point normalization constants of double-trace operators. In what follows we will be mostly interested in double-trace operators with free-field dimensions 4 and 5. We will see that such a splitting does occur for all the operators except the operators in the 20 and 105 irreps.
## 4 Anomalous dimensions of double-trace operators
In this section we determine the anomalous dimensions of double-trace operators and the leading $`1/N^2`$ corrections to their 2- and 3-point function normalization constants $`C_D(N)`$. To this end, we study the asymptotic behavior of the 4-point functions of CPOs in the direct channel $`x_{12}^2,x_{34}^20`$. Since we know all the 4-point functions, we do not need to consider the crossed channels. It is well-known that a conformally-invariant 4-point function is given as a general analytic function of two variables, which are here conveniently chosen to be the “biharmonic ratios”
$$u=\frac{x_{12}^2x_{34}^2}{x_{13}^2x_{24}^2},v=\frac{x_{12}^2x_{34}^2}{x_{14}^2x_{23}^2}.$$
We also use in the following the variable $`Y=1\frac{v}{u}`$. The biharmonic ratios above and the variable $`Y`$ have the property that $`u,v,Y0`$ as $`x_{12}^2,x_{34}^20`$.
To perform the computation we need to know the contributions of various quasi-primary operators and their descendents in the 4-point functions of CPOs, i.e. the conformal partial wave amplitudes of quasi-primary operators. We restrict ourselves mainly to the contributions of scalar, vector and second rank symmetric traceless tensor operators. Let the OPE of CPOs be of the form
$`O^{I_1}(x_1)O^{I_2}(x_2)`$ $`=`$ $`C_𝒥^{I_1I_2}({\displaystyle \frac{C_{OOS}}{C_S}}{\displaystyle \frac{1}{x_{12}^{4\mathrm{\Delta }_S}}}[S^𝒥]+{\displaystyle \frac{C_{OOT}}{C_T}}{\displaystyle \frac{x_{12}^\mu x_{12}^\nu }{x_{12}^{6\mathrm{\Delta }_T}}}[T_{\mu \nu }^𝒥]`$ (4.1)
$`+{\displaystyle \frac{C_{OOV}}{C_V}}{\displaystyle \frac{x_{12}^\mu }{x_{12}^{5\mathrm{\Delta }_V}}}[V_\mu ^𝒥]+\mathrm{}),`$
where $`𝒥`$ denotes an index of an irreducible representation of the $`R`$-symmetry group $`SO(6)`$, $`C_𝒥^{I_1I_2}`$ are the Clebsh-Gordon coefficients and $`\mathrm{\Delta }_S,\mathrm{\Delta }_T,\mathrm{\Delta }_V`$ are the conformal dimensions of the scalar, tensor and vector operators respectively. For any of the operators, $`C_𝒪`$ and $`C_{OO𝒪}`$ denote the normalization constant in the 2-point function $`𝒪(x_1)𝒪(x_2)`$ and the coupling constant in the three-point function $`O^I(x_1)O^J(x_2)𝒪(x_3)`$, respectively. Then, one can show that the short-distance expansion of the conformal partial amplitudes of the scalar S, tensor T and vector V operators can be written as
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)={\displaystyle \frac{C_𝒥^{I_1I_2}C_𝒥^{I_3I_4}}{x_{12}^4x_{34}^4}}`$
$`\times [{\displaystyle \frac{C_{OOS}^2}{C_S}}v^{\frac{\mathrm{\Delta }_S}{2}}(1+{\displaystyle \frac{\mathrm{\Delta }_S}{4}}Y+{\displaystyle \frac{\mathrm{\Delta }_S^3}{16(\mathrm{\Delta }_S1)(\mathrm{\Delta }_S+1)}}v(1+{\displaystyle \frac{\mathrm{\Delta }_S+2}{4}}Y)+\mathrm{})`$
$`+{\displaystyle \frac{C_{OOT}^2}{C_T}}v^{\frac{\mathrm{\Delta }_T}{2}1}({\displaystyle \frac{1}{4}}Y^2{\displaystyle \frac{1}{\mathrm{\Delta }_T}}v{\displaystyle \frac{1}{\mathrm{\Delta }_T}}vY\mathrm{})`$
$`+{\displaystyle \frac{C_{OOV}^2}{C_V}}v^{\frac{\mathrm{\Delta }_V1}{2}}({\displaystyle \frac{1}{2}}Y+\mathrm{})].`$ (4.2)
The formulas for the leading contributions of a rank-2 traceless symmetric tensor and a vector can be generalized to the case of a rank-$`l`$ traceless symmetric tensor of dimension $`\mathrm{\Delta }_l`$ and one gets a leading term of the form
$$v^{\frac{\mathrm{\Delta }_ll}{2}}Y^l.$$
For this reason a term of the form $`v^{\mathrm{\Delta }/2}F(Y)`$ in a 4-point function contains, in principle, the contributions not only from a scalar operator, but also from any symmetric tensor operator of rank $`l`$ and conformal dimension $`\mathrm{\Delta }+l`$. Moreover, (4.2) shows that the anomalous dimensions are related to terms of the type $`v^{\frac{\mathrm{\Delta }_S^{(0)}}{2}}\mathrm{log}v`$ for scalar operators, $`v^{\frac{\mathrm{\Delta }_V^{(0)}1}{2}}Y\mathrm{log}v`$ for vector operators and $`v^{\frac{\mathrm{\Delta }_T^{(0)}2}{2}}Y^2\mathrm{log}v`$ for rank-2 tensor operators.
The 4-point functions of CPOs were computed in the supergravity approximation in and can be written as follows
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)={\displaystyle \frac{\delta ^{I_1I_2}\delta ^{I_3I_4}}{x_{12}^4x_{34}^4}}+{\displaystyle \frac{\delta ^{I_1I_3}\delta ^{I_2I_4}}{x_{13}^4x_{24}^4}}+{\displaystyle \frac{\delta ^{I_1I_4}\delta ^{I_2I_3}}{x_{14}^4x_{23}^4}}`$ (4.3)
$`+{\displaystyle \frac{8}{N^2\pi ^2}}[{\displaystyle \frac{C_{I_1I_2I_3I_4}^{}}{x_{12}^2x_{34}^2}}(2(x_{13}^2x_{24}^2x_{14}^2x_{23}^2)D_{2222}`$
$`x_{24}^2D_{1212}x_{13}^2D_{2121}+x_{14}^2D_{2112}+x_{23}^2D_{1221})`$
$`+\delta ^{I_1I_2}\delta ^{I_3I_4}\left({\displaystyle \frac{1}{2x_{34}^2}}D_{2211}+{\displaystyle \frac{(x_{13}^2x_{24}^2+x_{14}^2x_{23}^2x_{12}^2x_{34}^2)}{x_{34}^2}}D_{3322}+{\displaystyle \frac{3}{2}}D_{2222}\right)`$
$`+2C_{I_1I_2I_3I_4}^+({\displaystyle \frac{1}{x_{34}^2}}D_{2211}+4x_{34}^2D_{2233}3D_{2222})+t+u],`$ (4.4)
where $`C_{I_1I_2I_3I_4}^\pm =\frac{1}{2}\left(C_{I_1I_2I_3I_4}\pm C_{I_2I_1I_3I_4}\right)`$ and $`t`$ and $`u`$ stand for the contributions of the $`t`$\- and $`u`$-channels obtained by the interchange $`14`$ and $`13`$, respectively. The $`D`$-functions are defined as
$`D_{\mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_4}(x_1,x_2,x_3,x_4)=`$ (4.5)
$`={\displaystyle \mathrm{d}^{d+1}\widehat{x}\frac{x_0^{d1+\mathrm{\Delta }_1+\mathrm{\Delta }_2+\mathrm{\Delta }_3+\mathrm{\Delta }_4}}{[x_0^2+(xx_1)^2]^{\mathrm{\Delta }_1}[x_0^2+(xx_2)^2]^{\mathrm{\Delta }_2}[x_0^2+(xx_3)^2]^{\mathrm{\Delta }_3}[x_0^2+(xx_4)^2]^{\mathrm{\Delta }_4}}}.`$
It is convenient to represent $`D`$-functions in the form
$`D_{\mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_4}(x_1,x_2,x_3,x_4)=`$
$`={\displaystyle \frac{\overline{D}_{\mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_4}(v,Y)}{(x_{12}^2)^{\frac{\mathrm{\Delta }_1+\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_4}{2}}(x_{13}^2)^{\frac{\mathrm{\Delta }_1+\mathrm{\Delta }_3\mathrm{\Delta }_2\mathrm{\Delta }_4}{2}}(x_{23}^2)^{\frac{\mathrm{\Delta }_2+\mathrm{\Delta }_3+\mathrm{\Delta }_4\mathrm{\Delta }_1}{2}}(x_{14}^2)^{\mathrm{\Delta }_4}}}.`$
As shown in Appendix C, a $`\overline{D}`$-function is given by a convergent series in $`v`$ and $`Y`$. In terms of the biharmonic ratios $`u`$ and $`v`$ the 4-point function acquires the form
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)={\displaystyle \frac{1}{x_{12}^4x_{34}^4}}[\delta ^{I_1I_2}\delta ^{I_3I_4}+u^2\delta ^{I_1I_3}\delta ^{I_2I_4}+v^2\delta ^{I_1I_4}\delta ^{I_2I_3}]`$ (4.6)
$`+{\displaystyle \frac{8}{\pi ^2N^2}}{\displaystyle \frac{1}{x_{12}^4x_{34}^4}}\{C_{I_1I_2I_3I_4}^{}[\overline{D}_{2222}(2v2{\displaystyle \frac{v^2}{u}}+vuv^2u{\displaystyle \frac{v^3}{u}}+v^3)`$
$`+\overline{D}_{1212}\left({\displaystyle \frac{2v^2}{u}}v^2+{\displaystyle \frac{v^3}{u}}\right)+\overline{D}_{2112}\left(2vvu+v^2\right)`$
$`+\overline{D}_{2211}(vuv^2)+\overline{D}_{2323}({\displaystyle \frac{4v^3}{u}})+4v^2\overline{D}_{3223}]`$
$`+C_{I_1I_2I_3I_4}^+[\overline{D}_{2222}(12v^2vu{\displaystyle \frac{v^3}{u}}+v^2u+v^3)`$ (4.7)
$`+\overline{D}_{1212}\left(v^2+{\displaystyle \frac{v^3}{u}}\right)+\overline{D}_{2112}\left(vu+v^2\right)`$
$`+\overline{D}_{2211}(2vvuv^2)+8v^2\overline{D}_{3322}+{\displaystyle \frac{4v^3}{u}}\overline{D}_{2323}+4v^2\overline{D}_{3223}]`$
$`+C_{I_1I_3I_2I_4}[\overline{D}_{2222}(6v^2+vu+{\displaystyle \frac{v^3}{u}}v^2uv^3)`$ (4.8)
$`+\overline{D}_{1212}\left(v^2{\displaystyle \frac{v^3}{u}}\right)+\overline{D}_{2112}\left(vu+v^2\right)`$
$`+\overline{D}_{2211}(vu+v^2)+{\displaystyle \frac{4v^3}{u}}\overline{D}_{2323}+4v^2\overline{D}_{3223}]`$
$`+\delta ^{I_1I_2}\delta ^{I_3I_4}\left[{\displaystyle \frac{3}{2}}v^2\overline{D}_{2222}{\displaystyle \frac{1}{2}}v\overline{D}_{2211}+\overline{D}_{3322}(v+{\displaystyle \frac{v^2}{u}}v^2)\right]`$
$`+\delta ^{I_1I_3}\delta ^{I_2I_4}\left[{\displaystyle \frac{3}{2}}v^2\overline{D}_{2222}{\displaystyle \frac{1}{2}}v^2\overline{D}_{1212}+\overline{D}_{2323}(v^2{\displaystyle \frac{v^3}{u}}+v^3)\right]`$
$`+\delta ^{I_1I_4}\delta ^{I_2I_3}[{\displaystyle \frac{3}{2}}v^2\overline{D}_{2222}{\displaystyle \frac{1}{2}}v^2\overline{D}_{2112}+\overline{D}_{3223}(v^2+{\displaystyle \frac{v^3}{u}}+v^3)]\}.`$
This 4-point function is given as a sum of contributions from quasi-primary operators transforming in the six irreducible representations (3.8) of $`SO(6)`$. It is clear that to obtain a contribution of operators belonging to a $`D`$-dimensional irrep one should multiply the 4-point function by a $`SO(6)`$ tensor $`C_{𝒥_D}^{I_1I_2}C_{𝒥_D}^{I_3I_4}`$ which is a projector onto the irrep.
In what follows it will be sometimes useful to compare the short-distance expansion of the 4-point function (4.7) with the one of the free field 4-point function (3.2), which in terms of the biharmonic ratios takes the form
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)_{fr}={\displaystyle \frac{1}{x_{12}^4x_{34}^4}}[\delta ^{I_1I_2}\delta ^{I_3I_4}+u^2\delta ^{I_1I_3}\delta ^{I_2I_4}+v^2\delta ^{I_1I_4}\delta ^{I_2I_3}`$ (4.9)
$`+{\displaystyle \frac{4}{N^2}}((u+v)C_{I_1I_2I_3I_4}^++(vu)C_{I_1I_2I_3I_4}^{}+uvC_{I_1I_3I_2I_4})].`$ (4.10)
### 4.1 Projection on the singlet
First we project the 4-point function on the singlet part that amounts to applying to it $`\frac{1}{400}\delta ^{I_1I_2}\delta ^{I_3I_4}`$. From the strong coupling OPE (3.12) we expect to find the stress tensor contribution and a contribution of the double-trace scalar operator $`O_1`$ of approximate dimension 4.
The result for the connected part is
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)|_\mathrm{𝟏}={\displaystyle \frac{8}{20\pi ^2N^2}}{\displaystyle \frac{\delta ^{I_1I_2}\delta ^{I_3I_4}}{x_{12}^4x_{34}^4}}[`$
$`\overline{D}_{2222}\left(9v^2{\displaystyle \frac{3v^3}{u}}3vu+3v^2u+3v^3\right)`$
$`+\overline{D}_{1212}\left({\displaystyle \frac{19}{6}}v^2+3{\displaystyle \frac{v^3}{u}}\right)+\overline{D}_{2112}\left(3vu+{\displaystyle \frac{19}{6}}v^2\right)`$ (4.11)
$`+\overline{D}_{2211}\left({\displaystyle \frac{10}{3}}v3vu3v^2\right)+\overline{D}_{3322}\left(20{\displaystyle \frac{v^2}{u}}+20v+{\displaystyle \frac{20}{3}}v^2\right)`$
$`+\overline{D}_{2323}(v^2+{\displaystyle \frac{41v^3}{3u}}+v^3)+\overline{D}_{3223}({\displaystyle \frac{41}{3}}v^2+{\displaystyle \frac{v^3}{u}}+v^3)]`$
Using the formulas for the $`\overline{D}`$-functions from the Appendix C, we can find that the most singular terms of the $`v`$-expansion are
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)|_\mathrm{𝟏}={\displaystyle \frac{\delta ^{I_1I_2}\delta ^{I_3I_4}}{N^2x_{12}^4x_{34}^4}}\left[vF_1(Y)+v^2F_2(Y)+v^2\mathrm{log}vG_2(Y)\right],`$ (4.12)
where
$`F_1(Y)`$ $`=`$ $`{\displaystyle \frac{4Y^28Y}{Y^3}}+{\displaystyle \frac{4(6+6YY^2)\mathrm{log}(1Y)}{3Y^3}},`$
$`F_2(Y)`$ $`=`$ $`{\displaystyle \frac{1680+3360Y2108Y^2+428Y^321Y^4}{15(1Y)Y^4}},`$
$``$ $`{\displaystyle \frac{4\left(11401890Y+962Y^2151Y^3+5Y^4\right)}{15Y^5}}\mathrm{log}(1Y)`$
$`+`$ $`{\displaystyle \frac{16(Y2)(66Y+Y^2)}{Y^5}}\mathrm{Li}_2(Y),`$
$`G_2(Y)`$ $`=`$ $`{\displaystyle \frac{4(66Y+Y^2)}{3Y^4}}\left({\displaystyle \frac{1212Y+Y^2}{Y1}}+{\displaystyle \frac{6\left(Y2\right)\mathrm{log}(1Y)}{Y}}\right).`$
Expanding the functions in powers of $`Y`$ we then obtain
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)|_\mathrm{𝟏}={\displaystyle \frac{1}{N^2}}{\displaystyle \frac{\delta ^{I_1I_2}\delta ^{I_3I_4}}{x_{12}^4x_{34}^4}}[{\displaystyle \frac{2}{45}}vY^2+v^2({\displaystyle \frac{47}{225}}{\displaystyle \frac{4}{5}}\mathrm{log}v)`$
$`{\displaystyle \frac{43}{225}}v^2Y].`$ (4.13)
Comparing this asymptotics with (4.2), we see that the contribution from a scalar field of dimension 2 is absent, as it should be, since the Konishi field acquires large anomalous dimension and decouples in the strong coupling limit. We also get the relation:
$`{\displaystyle \frac{C_{OOT}^2}{4C_T}}={\displaystyle \frac{2}{45N^2}}.`$
Since for $`C_{OOT}`$ one has $`C_{OOT}=\frac{4}{3\pi ^2}\frac{\lambda }{N}`$ <sup>6</sup><sup>6</sup>6This value of the coupling constant is fixed by a conformal Ward identity , the same value was also obtained in the supergravity approximation in . one finds at strong coupling
$`C_T={\displaystyle \frac{10\lambda ^2}{\pi ^4}},`$
which represents the normalization of the complete stress tensor of the $`𝒩=4`$ SYM<sub>4</sub> .
As it was discussed above, a term of the form $`vF(Y)`$ contains, in general, contributions from all traceless symmetric tensor operators of rank $`l`$ and dimension $`2+l`$. However, comparing $`F_1(Y)`$ in (4.12) with the corresponding term in the conformal partial wave amplitude of the stress tensor (7.14) we see that they coincide. Thus, the strong coupling OPE does not contain single-trace rank-$`l`$ traceless symmetric tensors with dimension $`2+l`$ in its singlet part. Nevertheless, it may in principle contain tensors of dimension $`4+l`$ or higher. However, as it was shown in section 3 a possible single-trace scalar operator of dimension 4 vanishes. Thus the only scalar operator of approximate dimension 4 is the double-trace operator $`O_\mathrm{𝟏}`$.<sup>7</sup><sup>7</sup>7The free field theory operator $`O_\mathrm{𝟏}^{fr}`$ probably splits into a linear combination of $`O_\mathrm{𝟏}`$ and an operator $`O_\mathrm{𝟏}^{str}`$ dual to a string mode. However, the coefficient in front of $`O_\mathrm{𝟏}^{str}`$ is of order $`1/N`$, and even if the latter operator does not decouple in the strong coupling limit it cannot contribute to $`\mathrm{log}`$-dependent terms in 4-point functions. In the following, when discussing double-trace operators in other irreps we simply assume that operators such as $`O_\mathrm{𝟏}^{str}`$ above do decouple, making at the same time a consistency check to confirm our assumption.
The formula (4.13) also allows us to determine the anomalous dimension of $`O_\mathrm{𝟏}`$. Assuming the existence at strong coupling of a scalar field with dimension $`\mathrm{\Delta }=\mathrm{\Delta }^{(0)}+\mathrm{\Delta }^{(1)}`$, where $`\mathrm{\Delta }^{(0)}=4`$ and $`\mathrm{\Delta }^{(1)}`$ is the anomalous dimension, we find that
$$v^{\frac{\mathrm{\Delta }}{2}}=v^2+\frac{1}{2}\mathrm{\Delta }^{(1)}v^2\mathrm{log}v+\mathrm{}$$
Since there is only one operator of approximate dimension 4, we do not face the problem of operator mixing and from (4.2) we get
$$\frac{1}{2}\frac{C_{OOO_\mathrm{𝟏}}^2}{C_{O_\mathrm{𝟏}}}\mathrm{\Delta }^{(1)}=\frac{4}{5N^2}.$$
Since $`\mathrm{\Delta }^{(1)}`$ is of order $`1/N^2`$ we use for $`\frac{C_{OOO_\mathrm{𝟏}}^2}{C_{O_\mathrm{𝟏}}}`$ the $`O(1)`$ result which is $`1/10`$. In this way we obtain
$`\mathrm{\Delta }^{(1)}={\displaystyle \frac{16}{N^2}},`$ (4.14)
for the anomalous dimension of $`O_\mathrm{𝟏}`$. This coincides with the anomalous dimension of the operator $`B`$ considered in , as it should be, since $`B`$ is a descendent operator of $`O_\mathrm{𝟏}`$.
We can also find the leading $`1/N^2`$ correction to the 2- and 3-point normalization constant $`C_{O_\mathrm{𝟏}}`$. Writing as
$$C_{O_\mathrm{𝟏}}=\frac{1}{10}\left(1+\frac{1}{N^2}C_{O_\mathrm{𝟏}}^{(1)}\right),$$
and taking into account that $`C_{O_\mathrm{𝟏}}=C_{OOO_\mathrm{𝟏}}`$, we find from the term of order $`v^2`$
$$C_{O_\mathrm{𝟏}}^{(1)}=\frac{38}{15}.$$
Finally, we can make a consistency check of our computation. Namely, since we know corrections to the conformal dimension, $`\mathrm{\Delta }^{(1)}=16/N^2`$ and to the structure constant we can compute the term of order $`v^2Y`$ by using (4.2), in order to compare it with the corresponding value obtained from our 4-point function. Taking into account the contribution of the stress tensor we get from (4.2) and from the expansion of our 4-point function the same number $`\frac{43}{225}`$. This also confirms that there is only one operator of approximate dimension 4 in the strong coupling OPE, and that the operator $`O_\mathrm{𝟏}^{str}`$ decouples in the strong coupling limit.
We can also compute the 2-point normalization constant in free field theory by using (3.1) and the definition of the operator. A simple calculation gives
$$C_{O_\mathrm{𝟏}}=\frac{1}{10}\left(1+\frac{2}{3N^2}\right).$$
Thus, not only the conformal dimension but also the 2- and 3-point normalization constants get $`\frac{1}{N^2}`$ corrections in the strong coupling limit.
### 4.2 Projection on 20
According to (4.2), to obtain the contribution of the operators transforming in a D-dimensional irrep, we should multiply the 4-point function by the projector onto the representation
$`(P_𝐃)_{I_1I_2I_3I_4}={\displaystyle \frac{1}{\nu _D}}C_{𝒥_D}^{I_1I_2}C_{𝒥_D}^{I_3I_4},`$ (4.15)
where
$$\nu _D=\underset{I_i}{}C_{𝒥_D}^{I_1I_2}C_{𝒥_D}^{I_3I_4}C_{𝒥_D^{}}^{I_1I_2}C_{𝒥_D^{}}^{I_3I_4},$$
is the dimension of the irrep so that $`P_𝐃^2=\frac{1}{\nu _D}`$.
The projector on the $`\mathrm{𝟐𝟎}`$ can be easily found by taking into account that the Clebsh-Gordon coefficient $`C_{𝒥_{20}}^{I_1I_2}`$ is proportional to the $`SO(6)`$ tensor $`C^{I_1I_2I_3}`$. Then, one can show that
$`(P_{\mathrm{𝟐𝟎}})_{I_1I_2I_3I_4}={\displaystyle \frac{3}{100}}\left(C_{I_1I_2I_3I_4}^+{\displaystyle \frac{1}{6}}\delta _{I_1I_2}\delta _{I_3I_4}\right).`$ (4.16)
Using the Table 1 from Appendix D for the contractions of the projector with the $`SO(6)`$ tensors appearing in the 4-point function, we find the contribution of the operators in the $`\mathrm{𝟐𝟎}`$ to the connected part of the 4-point function
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)|_{\mathrm{𝟐𝟎}}={\displaystyle \frac{8}{\pi ^2N^2}}{\displaystyle \frac{C_{𝒥_{20}}^{I_1I_2}C_{𝒥_{20}}^{I_3I_4}}{x_{12}^4x_{34}^4}}[`$
$`\overline{D}_{2222}\left(18v^2{\displaystyle \frac{3}{2}}vu{\displaystyle \frac{3v^3}{2u}}+{\displaystyle \frac{3}{2}}v^2u+{\displaystyle \frac{3}{2}}v^3\right)`$
$`+\overline{D}_{1212}\left({\displaystyle \frac{4}{3}}v^2+{\displaystyle \frac{3v^3}{2u}}\right)+\overline{D}_{2112}\left({\displaystyle \frac{3}{2}}vu+{\displaystyle \frac{4}{3}}v^2\right)+\overline{D}_{2211}\left({\displaystyle \frac{10}{3}}v{\displaystyle \frac{3}{2}}vu{\displaystyle \frac{3}{2}}v^2\right)`$
$`+{\displaystyle \frac{40}{3}}v^2\overline{D}_{3322}+\overline{D}_{2323}(v^2+{\displaystyle \frac{19v^3}{3u}}+v^3)+\overline{D}_{3223}({\displaystyle \frac{19}{3}}v^2+{\displaystyle \frac{v^3}{u}}+v^3)].`$ (4.17)
Expanding the $`\overline{D}`$-functions in powers of $`v`$, we obtain
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)|_{\mathrm{𝟐𝟎}}={\displaystyle \frac{1}{N^2}}{\displaystyle \frac{C_{𝒥_{20}}^{I_1I_2}C_{𝒥_{20}}^{I_3I_4}}{x_{12}^4x_{34}^4}}[`$
$`vF_1(Y)+v^2F_2(Y)+v^2\mathrm{log}vG_2(Y)],`$ (4.18)
where
$`F_1(Y)`$ $`=`$ $`{\displaystyle \frac{40\mathrm{log}(1Y)}{3Y}},`$
$`F_2(Y)`$ $`=`$ $`{\displaystyle \frac{8\left(6565Y+6Y^2\right)}{3(1Y)Y^2}}{\displaystyle \frac{20\left(7449Y+2Y^2\right)}{3Y^3}}\mathrm{log}(1Y)`$
$`+`$ $`{\displaystyle \frac{160(Y2)}{Y^3}}\mathrm{Li}_2(Y),`$
$`G_2(Y)`$ $`=`$ $`{\displaystyle \frac{40}{3Y^2}}\left({\displaystyle \frac{1212Y+Y^2}{Y1}}+{\displaystyle \frac{6\left(Y2\right)\mathrm{log}(1Y)}{Y}}\right).`$
Expanding the above functions in powers of $`Y`$ we finally obtain
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)|_{\mathrm{𝟐𝟎}}`$ $`=`$ $`{\displaystyle \frac{1}{N^2}}{\displaystyle \frac{C_{𝒥_{20}}^{I_1I_2}C_{𝒥_{20}}^{I_3I_4}}{x_{12}^4x_{34}^4}}[{\displaystyle \frac{40}{3}}v+{\displaystyle \frac{26}{9}}v^2(1+Y)`$ (4.19)
$`{\displaystyle \frac{4}{3}}v^2Y^2\mathrm{log}v].`$
The analysis of the results obtained follows the one in the previous subsection. Firstly, comparing $`F_1(Y)`$ in (4.18) with the corresponding term of the conformal partial amplitude of a scalar operator of dimension 2 (7.2), we see that they coincide.<sup>8</sup><sup>8</sup>8Recall that for the lowest weight CPOs one has $`\frac{C_{OOO}^2}{C_O}=\frac{40}{3N^2}`$. Therefore, all single-trace rank-$`l`$ traceless tensors of dimension $`2+l`$ transforming in the $`\mathrm{𝟐𝟎}`$ are absent in the OPE. Then, the only scalar operator of approximate dimension 4 is the double-trace operator $`O_{\mathrm{𝟐𝟎}}`$. Moreover, we see that $`\mathrm{log}v`$-dependent terms appear starting from the term $`v^2Y^2\mathrm{log}v`$. Thus we conclude from (4.2) that the double-trace operator $`O_{\mathrm{𝟐𝟎}}`$ has protected conformal dimension. It is worth noting that the non-renormalization of the conformal dimension of this operator is not related to the shortening condition discussed in and is a prediction of the AdS/CFT correspondence. The first operators which acquire anomalous dimensions are scalar and tensor operators of approximate dimension 6.
The first $`1/N^2`$ correction to the 2- and 3-point normalization constant $`C_{O_{\mathrm{𝟐𝟎}}}`$ can also be easily found. Writing the constant as
$$C_{O_{\mathrm{𝟐𝟎}}}=2\left(1+\frac{1}{N^2}C_{O_{\mathrm{𝟐𝟎}}}^{(1)}\right),$$
and taking into account the contribution of the single-trace operator $`O^I`$ and that $`C_{O_{\mathrm{𝟐𝟎}}}=C_{OOO_{\mathrm{𝟐𝟎}}}`$, we find from the term of order $`v^2`$
$$C_{O_{\mathrm{𝟐𝟎}}}^{(1)}=\frac{1}{3}.$$
The 2-point normalization constant can be also computed in free field theory by using (3.1) and the definition of the operator (3.10) and appears to coincide with the value obtained in the strong coupling limit
$$C_{O_{\mathrm{𝟐𝟎}}}=2\left(1+\frac{1}{3N^2}\right).$$
Thus, both the conformal dimension and the 2-point function normalization constant of the double-trace operator in the $`\mathrm{𝟐𝟎}`$ are non-renormalized in the strong coupling limit. This also shows that in this case there is no splitting, and the free field theory double-trace operator coincides with $`O_{\mathrm{𝟐𝟎}}`$.
### 4.3 Projection on 105
The free field theory OPE (3.11) and the strong coupling OPE (3.12) do not contain single-trace operators transforming in the 105 irrep. Thus, only double-trace operators contribute to this part of the 4-point function. The corresponding connected contribution can be easily found using the Table 1 from Appendix D and is given by
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)|_{\mathrm{𝟏𝟎𝟓}}={\displaystyle \frac{8}{\pi ^2N^2}}{\displaystyle \frac{C_{𝒥_{105}}^{I_1I_2}C_{𝒥_{105}}^{I_3I_4}}{x_{12}^4x_{34}^4}}[`$
$`\overline{D}_{2222}\left(3v^2+vu+{\displaystyle \frac{v^3}{u}}v^2uv^3\right)`$
$`+\overline{D}_{1212}\left({\displaystyle \frac{1}{2}}v^2{\displaystyle \frac{v^3}{u}}\right)+\overline{D}_{2112}\left(vu+{\displaystyle \frac{1}{2}}v^2\right)+\overline{D}_{2211}\left(vu+v^2\right)`$
$`+\overline{D}_{2323}({\displaystyle \frac{3v^3}{u}}+v^3+v^2)+\overline{D}_{3223}(3v^2+v^3+{\displaystyle \frac{v^3}{3}})].`$ (4.20)
Expanding the $`\overline{D}`$-functions in powers of $`v`$, we obtain
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)|_{\mathrm{𝟏𝟎𝟓}}`$ $`=`$ $`{\displaystyle \frac{1}{N^2}}{\displaystyle \frac{C_{𝒥_{105}}^{I_1I_2}C_{𝒥_{105}}^{I_3I_4}}{x_{12}^4x_{34}^4}}[v^2F_2(Y)+v^3F_3(Y)`$ (4.21)
$`+`$ $`v^4F_4(Y)+v^4\mathrm{log}vG_4(Y)],`$
where
$`F_2(Y)`$ $`=`$ $`{\displaystyle \frac{4}{1Y}},`$
$`F_3(Y)`$ $`=`$ $`{\displaystyle \frac{4\left(Y2\right)}{(1Y)Y^2}}{\displaystyle \frac{8}{Y^3}}\mathrm{log}(1Y),`$
$`F_4(Y)`$ $`=`$ $`{\displaystyle \frac{4\left(2828Y+3Y^2\right)}{(1Y)Y^4}}{\displaystyle \frac{8\left(3825Y+Y^2\right)}{Y^5}}\mathrm{log}(1Y)`$
$`+`$ $`{\displaystyle \frac{96(Y2)}{Y^5}}\mathrm{Li}_2(Y),`$
$`G_4(Y)`$ $`=`$ $`{\displaystyle \frac{8}{Y^4}}\left({\displaystyle \frac{1212Y+Y^2}{Y1}}+{\displaystyle \frac{6\left(Y2\right)\mathrm{log}(1Y)}{Y}}\right).`$
Since only double-trace operators contribute, it is useful to compare (4.21) with the corresponding part of the free field theory 4-point function (4.10)
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)_{fr}|_{\mathrm{𝟏𝟎𝟓}}={\displaystyle \frac{C_{𝒥_{105}}^{I_1I_2}C_{𝒥_{105}}^{I_3I_4}}{x_{12}^4x_{34}^4}}[`$
$`v^2(1+{\displaystyle \frac{1}{(1Y)^2}})+{\displaystyle \frac{v^2}{N^2}}{\displaystyle \frac{4}{1Y}}].`$ (4.22)
The first term on the r.h.s. of this equation shows the disconnected part of the free field theory 4-point function. Comparing the term of order $`1/N^2`$ in (4.22) with the term $`v^2F_2(Y)`$ in (4.21), we see that they coincide. This means that the conformal dimensions and the leading corrections in $`1/N^2`$ to 2- and 3-point functions normalization constants of any symmetric traceless rank-$`2k`$ tensor operator of dimension $`4+2k`$ transforming in the $`\mathrm{𝟏𝟎𝟓}`$ coincide with the ones computed in free field theory. Thus, all these operators are non-renormalized in the strong coupling limit. The first correction to the 2- and 3-point functions normalization constant of the double-trace operator $`O_{\mathrm{𝟏𝟎𝟓}}`$ can be easily found from (4.22) and is given by
$$C_{O_{\mathrm{𝟏𝟎𝟓}}}=2\left(1+\frac{2}{N^2}\right).$$
The non-renormalization of the double-trace operator $`O_{\mathrm{𝟏𝟎𝟓}}`$ follows from the shortening conditions derived in , and was also checked in perturbation theory at small YM coupling in .
The expansion (4.21) also shows that the first $`\mathrm{log}v`$-term appears at order $`v^4`$. Therefore, all symmetric traceless rank-$`2k`$ tensor operators of dimension $`6+2k`$ transforming in the $`\mathrm{𝟏𝟎𝟓}`$ have protected conformal dimensions. Note, however, that the normalization constants of their 2- and 3-point functions certainly receive corrections at strong coupling, which are encoded in the function $`F_3(Y)`$. The vanishing of anomalous dimensions of these tensor operators does not seem to follow from any known non-renormalization theorem. These results also demonstrate that the free field theory symmetric traceless rank-$`2k`$ tensor operators of dimension $`4+2k`$ do not split, while the ones with dimension $`6+2k`$ do.
Since $`G_4(Y)=\frac{4}{5}+\mathrm{}`$ the first double-trace operator in the $`\mathrm{𝟏𝟎𝟓}`$ which acquires anomalous dimension is the scalar operator with approximate dimension 8.
### 4.4 Projection on $`\mathrm{𝟖𝟒}`$
Just as it was the case for the operators in the 105, only double-trace operators transforming in the 84 irrep can contribute to this part of the 4-point function. The corresponding connected contribution is again found by using the Table 1 from Appendix D:
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)|_{\mathrm{𝟖𝟒}}={\displaystyle \frac{8}{\pi ^2N^2}}{\displaystyle \frac{C_{𝒥_{84}}^{I_1I_2}C_{𝒥_{84}}^{I_3I_4}}{x_{12}^4x_{34}^4}}[`$
$`\overline{D}_{2222}\left(6v^2{\displaystyle \frac{uv}{2}}+{\displaystyle \frac{uv^2}{2}}+{\displaystyle \frac{v^3}{2}}{\displaystyle \frac{v^3}{2u}}\right)`$
$`+\overline{D}_{1212}\left(v^2+{\displaystyle \frac{v^3}{2u}}\right)+\overline{D}_{2112}\left({\displaystyle \frac{uv}{2}}v^2\right)\overline{D}_{2211}\left({\displaystyle \frac{uv}{2}}+{\displaystyle \frac{v^2}{2}}\right)`$
$`+\overline{D}_{2323}(v^2+v^3{\displaystyle \frac{3v^3}{u}})+\overline{D}_{3223}(3v^2+v^3+{\displaystyle \frac{v^3}{u}})].`$ (4.23)
Expanding the $`\overline{D}`$-functions in powers of $`v`$, we obtain
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)|_{\mathrm{𝟖𝟒}}={\displaystyle \frac{1}{N^2}}{\displaystyle \frac{C_{𝒥_{84}}^{I_1I_2}C_{𝒥_{84}}^{I_3I_4}}{x_{12}^4x_{34}^4}}[v^2F_2(Y)+v^3F_3(Y)`$
$`+v^3\mathrm{log}vG_3(Y)],`$ (4.24)
where
$`F_2(Y)`$ $`=`$ $`{\displaystyle \frac{8\left(33Y+Y^2\right)}{(1Y)Y^2}}+{\displaystyle \frac{12(Y2)}{Y^3}}\mathrm{log}(1Y),`$
$`F_3(Y)`$ $`=`$ $`{\displaystyle \frac{8(Y2)\left(2121Y+2Y^2\right)}{(1Y)Y^4}}+{\displaystyle \frac{4\left(228+264Y80Y^2+3Y^3\right)}{Y^5}}\mathrm{log}(1Y)`$
$``$ $`{\displaystyle \frac{144(Y2)^2}{Y^5}}\mathrm{Li}_2(Y),`$
$`G_3(Y)`$ $`=`$ $`{\displaystyle \frac{12(Y2)}{Y^4}}\left({\displaystyle \frac{1212Y+Y^2}{Y1}}+{\displaystyle \frac{6\left(Y2\right)\mathrm{log}(1Y)}{Y}}\right).`$
Since the first $`\mathrm{log}v`$-term appears at order $`v^3`$, all symmetric traceless rank-$`2k`$ tensor operators of dimension $`4+2k`$ transforming in the $`\mathrm{𝟖𝟒}`$ have protected conformal dimensions. The first double-trace operator in the $`\mathrm{𝟖𝟒}`$ which acquires an anomalous dimension is the scalar operator with approximate dimension 6. However, contrary to the case of the 105 irrep, the leading $`1/N^2`$ corrections to the normalization constants of the 2- and 3-point functions of these operators differ from their free field theory values. To see this we compare (4.24) with the corresponding part of the free field theory 4-point function (4.10)
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)_{fr}|_{\mathrm{𝟖𝟒}}={\displaystyle \frac{C_{𝒥_{84}}^{I_1I_2}C_{𝒥_{84}}^{I_3I_4}}{x_{12}^4x_{34}^4}}[v^2(1+{\displaystyle \frac{1}{(1Y)^2}})`$
$`{\displaystyle \frac{v^2}{N^2}}{\displaystyle \frac{2}{1Y}}].`$ (4.25)
Expanding (4.24) and (4.25) in powers of $`Y`$, we obtain the normalization constants of 2- and 3-point functions of the operator $`O_{\mathrm{𝟖𝟒}}`$ at strong coupling and in free field theory correspondingly as
$`C_{O_{\mathrm{𝟖𝟒}}}^{str}`$ $`=`$ $`2\left(1{\displaystyle \frac{3}{N^2}}\right),`$
$`C_{O_{\mathrm{𝟖𝟒}}}^{fr}`$ $`=`$ $`2\left(1{\displaystyle \frac{1}{N^2}}\right).`$
The vanishing of the anomalous dimensions of the double-trace operator $`O_{\mathrm{𝟖𝟒}}`$ follows from the shortening conditions discussed in and was also shown in perturbation theory at small YM coupling in . The difference between $`C_{O_{\mathrm{𝟖𝟒}}}^{str}`$ and $`C_{O_{\mathrm{𝟖𝟒}}}^{fr}`$ again may find a natural explanation in the fact that the corresponding free field theory operator undergoes a linear splitting on $`O_{\mathrm{𝟖𝟒}}`$ and $`𝒦_{\mathrm{𝟖𝟒}}`$, where $`O_{\mathrm{𝟖𝟒}}`$ has protected both its dimension and the normalization constants of the 2- and 3-point functions, while the operator $`𝒦_{\mathrm{𝟖𝟒}}`$ belongs to the Konishi multiplet and, therefore, decouples at strong coupling.
### 4.5 Projection on $`\mathrm{𝟏𝟓}`$
By using the projector $`(P_{\mathrm{𝟏𝟓}})_{I_1I_2I_3I_4}`$ constructed in the Appendix D and the results of Table 1 we find the following contribution of the operators in $`\mathrm{𝟏𝟓}`$ to the connected part of the 4-point function
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)|_{\mathrm{𝟏𝟓}}={\displaystyle \frac{8}{\pi ^2N^2}}{\displaystyle \frac{C_{𝒥_{15}}^{I_1I_2}C_{𝒥_{15}}^{I_3I_4}}{x_{12}^4x_{34}^4}}[`$
$`\overline{D}_{2222}\left(4v+2uv{\displaystyle \frac{4v^2}{u}}2uv^2+2v^3{\displaystyle \frac{2v^3}{u}}\right)`$
$`+\overline{D}_{1212}\left({\displaystyle \frac{3v^2}{2}}+{\displaystyle \frac{4v^2}{u}}+{\displaystyle \frac{2v^3}{u}}\right)+\overline{D}_{2112}\left(4v2uv+{\displaystyle \frac{3v^2}{2}}\right)`$
$`+\overline{D}_{2211}(2vu2v^2)+(v^2v^3{\displaystyle \frac{7v^3}{u}})\overline{D}_{2323}+(7v^2+v^3+{\displaystyle \frac{v^3}{2u}})\overline{D}_{3223}].`$
Expansion of the $`D`$-functions in powers of $`v`$ produces now the following expression for leading terms
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)|_{\mathrm{𝟏𝟓}}`$ $`=`$ $`{\displaystyle \frac{1}{N^2}}{\displaystyle \frac{C_{𝒥_{15}}^{I_1I_2}C_{𝒥_{15}}^{I_3I_4}}{x_{12}^4x_{34}^4}}[vF_1(Y)+v^2F_2(Y)`$ (4.26)
$`+v^2\mathrm{log}vG_2(Y)].`$
Here the functions $`F_1,F_2`$ and $`G_2`$ are given by
$`F_1(Y)`$ $`=`$ $`{\displaystyle \frac{16}{Y^2}}\left(2Y+(Y2)\mathrm{log}(1Y)\right),`$
$`F_2(Y)`$ $`=`$ $`{\displaystyle \frac{4(Y2)(5656Y+5Y^2)}{(Y1)Y^3}}+{\displaystyle \frac{8(152+176Y53Y^2+2Y^3)}{Y^4}}\mathrm{log}(1Y)`$
$``$ $`{\displaystyle \frac{192(Y2)^2}{Y^4}}\mathrm{Li}_2(Y),`$
$`G_2(Y)`$ $`=`$ $`{\displaystyle \frac{16(Y2)}{Y^3}}\left({\displaystyle \frac{1212Y+Y^2}{Y1}}+{\displaystyle \frac{6(Y2)\mathrm{log}(1Y)}{Y}}\right).`$
Expansion in powers of $`Y`$ produces the following leading terms
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)|_{\mathrm{𝟏𝟓}}={\displaystyle \frac{C_{𝒥_{15}}^{I_1I_2}C_{𝒥_{15}}^{I_3I_4}}{N^2x_{12}^4x_{34}^4}}[{\displaystyle \frac{8}{3}}vY+{\displaystyle \frac{12}{25}}v^2Y{\displaystyle \frac{16}{5}}v^2Y\mathrm{log}v].`$ (4.27)
The absence of the terms $`vY\mathrm{log}v`$ shows that the vector operator of the dimension 3, which is the $`R`$-symmetry current $`R_\mu ^{𝒥_{15}}`$, has protected conformal dimension. According to the discussion above, the function $`F_1(Y)`$ may receive contributions from single-trace rank $`2k+1`$ traceless tensors of dimension $`2k+3`$ transforming in $`\mathrm{𝟏𝟓}`$, which is what indeed happens in the free field theory limit. However, comparing the function $`F_1(Y)`$ with the relevant part of the conformal partial amplitude of the conserved vector current of dimension 3 (7.9) one concludes that they coincide, therefore, the corresponding tensors are absent in the strong-coupling OPE. Next, comparing (4.27) with eq.(4.2) we read off the value of the ratio
$$\frac{C_{OOR}^2}{2C_R}=\frac{8}{3N^2}.$$
Since the value of $`C_{OOR}`$ is fixed by the conformal Ward identity to be $`C_{OOR}=\frac{2^{1/2}}{\pi ^2}\frac{\lambda }{N}`$ one finds
$$C_R=\frac{3\lambda ^2}{8\pi ^4}$$
which corresponds to the normalization of the two-point function of the complete $`R`$-symmetry current of the $`𝒩=4`$ SYM<sub>4</sub> .
The function $`F_2(Y)`$ receives contributions both from the $`R`$-symmetry current and from traceless symmetric rank $`2k+1`$ tensors with approximate dimension $`2k+5`$. Since $`R_\mu ^{𝒥_{15}}`$ is non-renormalized, the presence of the function $`G_2`$ shows that operators from the above tensor tower acquire anomalous dimensions. We can find the anomalous dimension of the lowest current $`O_{15}`$ in this tower whose free field theory counterpart $`O_\mu ^{𝒥_{15}}`$ with conformal dimension $`\mathrm{\Delta }^{(0)}=5`$ was discussed in section 3. In fact in perturbation theory the free-field operator $`O_\mu ^{𝒥_{15}}`$ contains in the split a descendent of $`O_1`$ and currents from the Konishi and the $`\mathrm{\Xi }`$-multiplets. It is a descendent of $`O_\mathrm{𝟏}`$ that is responsible for the logarithmic term in (4.27) and, therefore, its anomalous dimension at strong coupling is $`\frac{16}{N^2}`$. Comparing the coefficient in front of $`v^2Y\mathrm{log}v`$ in (4.27) with the asymptotic (4.2) one gets
$$\frac{1}{4}\frac{C_{OOO_{15}}^2}{C_{O_{15}}}\mathrm{\Delta }^{(1)}=\frac{16}{5N^2}$$
and substituting $`\mathrm{\Delta }^{(1)}=\frac{16}{N^2}`$ one obtains $`\frac{C_{OOO_{15}}^2}{C_{O_{15}}}=\frac{4}{5}`$ that is different from the free-field ratio $`\frac{C_{OOO_{15}}^2}{C_{O_{15}}}=4`$.
### 4.6 Projection on $`\mathrm{𝟏𝟕𝟓}`$
Only double-trace operators transforming in the 175 appear in the free field theory OPE (3.11) and in the strong coupling OPE (3.12). Applying the projector $`(P_{\mathrm{𝟏𝟕𝟓}})_{I_1I_2I_3I_4}`$ constructed in Appendix D to the 4-point function we find the following expression for the contribution of the operators in the $`\mathrm{𝟏𝟕𝟓}`$:
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)|_{\mathrm{𝟏𝟕𝟓}}={\displaystyle \frac{8}{\pi ^2N^2}}{\displaystyle \frac{C_{𝒥_{175}}^{I_1I_2}C_{𝒥_{175}}^{I_3I_4}}{x_{12}^4x_{34}^4}}[{\displaystyle \frac{v^2}{2}}\overline{D}_{1212}+{\displaystyle \frac{v^2}{2}}\overline{D}_{2112}`$
$`+(v^2+v^3{\displaystyle \frac{v^3}{u}})\overline{D}_{2323}+(v^2v^3{\displaystyle \frac{v^3}{u}})\overline{D}_{3223}].`$
Expanding $`\overline{D}`$ functions in $`v`$ we keep the leading terms $`v^2`$ and $`v^3`$
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)|_{\mathrm{𝟏𝟕𝟓}}`$ $`=`$ $`{\displaystyle \frac{1}{N^2}}{\displaystyle \frac{C_{𝒥_{175}}^{I_1I_2}C_{𝒥_{175}}^{I_3I_4}}{x_{12}^4x_{34}^4}}[v^2F_2(Y)+v^3F_3(Y)`$
$`+v^3\mathrm{log}vG_3(Y)]`$
with
$`F_2(Y)`$ $`=`$ $`{\displaystyle \frac{4(Y(Y2)+2(Y1)\mathrm{log}(1Y))}{Y^2(Y1)}},`$
$`F_3(Y)`$ $`=`$ $`{\displaystyle \frac{4(2828Y+3Y^2)}{Y^3(Y1)}}{\displaystyle \frac{8(3825Y+Y^2)\mathrm{log}(1Y)}{Y^4}},`$
$`+`$ $`{\displaystyle \frac{96(Y2)}{Y^4}}\mathrm{Li}_2(Y),`$
$`G_3(Y)`$ $`=`$ $`{\displaystyle \frac{8}{Y^3}}\left({\displaystyle \frac{1212Y+Y^2}{Y1}}+{\displaystyle \frac{6(Y2)\mathrm{log}(1Y)}{Y}}\right).`$
The function $`F_2`$ receives contributions from tensor operators of rank $`2k+1`$ with approximate dimensions $`2k+5`$. Since the term proportional to $`v^2\mathrm{log}v`$ is absent, we conclude that these tensor operators have protected conformal dimensions. The lowest current $`O_\mu ^{𝒥_{175}}`$ among them, with dimension 5, was discussed in section 3. Note that these operators also contribute to $`F_3`$ together with operators of rank $`2k+1`$ and approximate dimensions $`2k+7`$. For the two terms of the $`Y`$-expansion one finds
$`O^{I_1}(x_1)O^{I_2}(x_2)O^{I_3}(x_3)O^{I_4}(x_4)|_{\mathrm{𝟏𝟕𝟓}}={\displaystyle \frac{1}{N^2}}{\displaystyle \frac{C_{𝒥_{175}}^{I_1I_2}C_{𝒥_{175}}^{I_3I_4}}{x_{12}^4x_{34}^4}}\left[{\displaystyle \frac{4}{3}}v^2Y2v^2Y^2\right].`$ (4.28)
This allows us to determine the $`1/N^2`$ correction to the 2- and 3-point normalization constant $`C_{O_{175}}`$ of the operator $`O_\mu ^{𝒥_{175}}`$. Taking into account that in free field theory $`C_{OOO_{175}}^{fr}=C_{O_{175}}^{fr}=4`$ as can be easily seen from the free field theory 4-point function (4.10), we write as
$$C_{O_{175}}=4\left(1+\frac{1}{N^2}C_{O_{175}}^{(1)}\right).$$
Then from the first term of order $`v^2`$ in (4.28) one finds
$$C_{O_{175}}^{(1)}=\frac{2}{3}.$$
Apparently, the splitting mechanism is again at work, i.e. the corresponding free field theory operator is split in two orthogonal parts carrying different representation of the supersymmetry; one has protected both its dimension and the normalization constants, while the other one is dual to a string mode and decouples at strong coupling.
## 5 Conclusions
We studied in detail the 4-point functions of the lowest weight CPOs and we showed that they have a structure compatible with the OPE of CPOs predicted by the AdS/CFT correspondence. We demonstrated that all power-singular terms in the 4-point functions exactly match the corresponding terms in the conformal partial wave amplitudes of the CPOs, of the $`R`$-symmetry current and of the stress tensor. As these operators are dual to type IIB supergravity fields, we concluded that the operators dual to string modes, which appear in the free field theory OPE, decouple in the strong coupling limit.
We also computed the anomalous dimensions and the leading $`1/N^2`$ corrections to the normalization constants of the 2- and 3-point functions of the scalar double-trace operators with approximate dimension 4 and of vector operators with approximate dimension 5. The only scalar double-trace operator that acquires an anomalous dimension appears to be the operator in the singlet of the $`R`$-symmetry group $`SO(6)`$. The double-trace operator in the 20 seems to be protected, however as this does not follow from the shortening condition discussed in we do not have a satisfactory explanation for such a non-renormalization property.
The anomalous dimension of the singlet operator is negative, hence this operator is relevant and can be used to study non-conformal deformations of the $`𝒩=4`$ SYM<sub>4</sub>. All other scalar double-trace operators have protected dimension 4 and are marginal. They can be added to the Lagrangian in order to study conformal deformations. Nevertheless, it is unclear at present how dual deformations of type IIB supergravity (or string theory) can be described.
We have also found several towers of traceless symmetric double-trace operators in the 105, 84 and 175 irreps, whose anomalous conformal dimensions vanish. The rank-$`2k`$ tensor operators of dimension $`6+2k`$ satisfy the shortening condition A’) of . However, even if they contain the highest weight states of the $`SU(2,2|4)`$ superalgebra the shortening condition A’) does not imply non-renormalization of the corresponding multiplets. On the other hand operators from other towers are certainly not the highest weight states, and at present we are not aware if the lowest weight states of their supermultiplets satisfy the shortening condition responsible for non-renormalization.
There are two interesting facts related to the structure of the leading log-dependent terms in the 4-point functions. Namely, all the functions $`G(Y)`$ which appear in (4.12), (4.18) and so on, differ from each other by some simple rational factors. We expect that this is an indication that the anomalous dimensions of all double-trace operators may be related by some relatively simple formula. Then, the leading $`\mathrm{log}v`$-dependent terms appear in the 4-point functions exactly at the same order of $`v`$ where the dilogarithm $`\mathrm{Li}_2`$ appears for the first time.
ACKNOWLEDGMENTS We would like to thank A. Tseytlin for valuable comments. G.A. is grateful to S. Theisen and, especially, to S. Kuzenko for discussions of the structure of the Konishi multiplet. S.F. is grateful to S. Mathur and, especially, to A. Tseytlin for valuable discussions. A.C.P. wishes to thank W. Rühl for sharing with him his insight on CFT. The work of G.A. was supported by the Alexander von Humboldt Foundation and in part by the RFBI grant N99-01-00166. The work of S.F. was supported by the U.S. Department of Energy under grant No. DE-FG02-96ER40967 and in part by RFBI grant N99-01-00190. The work of A. C. P. was supported by the Alexander von Humboldt Foundation. G.A. and A.C.P. wish to acknowledge the warm hospitality and financial support of the E.S.I. in Vienna where part of the work was done.
## 6 Appendix A. Free field OPE and conformal blocks
A quasi-primary field of the CFT appearing in the OPE together with all its derivative descendents is known as a conformal block. If two fields $`O^\alpha `$ and $`O^\beta `$ transforming in some representation of an $`R`$-symmetry group have the one and the same conformal dimension $`\mathrm{\Delta }`$ then their OPE has the following structure
$`O^\alpha (x)O^\beta (0)`$ $`=`$ $`{\displaystyle \frac{1}{(x^2)^{\frac{1}{2}(2\mathrm{\Delta }\mathrm{\Delta }_O)}}}C(x,)O^{\alpha \beta }(0)`$
$`+`$ $`{\displaystyle \frac{1}{(x^2)^{\frac{1}{2}(2\mathrm{\Delta }\mathrm{\Delta }_J+1)}}}C_\mu (x,)J_\mu ^{\alpha \beta }(0)`$
$`+`$ $`{\displaystyle \frac{1}{(x^2)^{\frac{1}{2}(2\mathrm{\Delta }\mathrm{\Delta }_T+2)}}}C_{\mu \nu }(x,)T_{\mu \nu }^{\alpha \beta }(0)+\mathrm{}.`$
Here we identify the leading quasi-primary fields with conformal dimensions $`\mathrm{\Delta }_O`$, $`\mathrm{\Delta }_J`$ and $`\mathrm{\Delta }_T`$ as a scalar $`O^{\alpha \beta }`$, a vector current $`J_\mu ^{\alpha \beta }`$ and a symmetric traceless second rank tensor $`T_{\mu \nu }^{\alpha \beta }`$ respectively. The OPE coefficient $`C(x,)`$ denotes a power series in derivatives generating the conformal block $`[O^{\alpha \beta }]`$ of the scalar $`O^{\alpha \beta }`$. Similarly we denote the OPE coefficients for for the other fields.
The structure of the conformal blocks is uniquely fixed by the conformal symmetry and it may be found by requiring consistency of the OPE with 2- and 3-point functions of the fields involved. In particularly, the conformal block of a scalar field with dimension $`\mathrm{\Delta }`$ is given by the following differential operator :
$`C(x,_y)`$ $`=`$ $`{\displaystyle \frac{1}{B(\frac{1}{2}\mathrm{\Delta }_O,\frac{1}{2}\mathrm{\Delta }_O)}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m!(\mathrm{\Delta }_O\eta +1)_m}}`$
$`\times `$ $`{\displaystyle _0^1}dt[t(1t)]^{\frac{1}{2}\mathrm{\Delta }_O1}\left({\displaystyle \frac{1}{4}}t(1t)x^2\mathrm{\Delta }_y\right)^me^{tx_y},`$
where the Euclidean space-time dimension $`d`$ enters as $`d=2\eta `$, $`x_y=x^\mu _{y,\mu }`$, $`\mathrm{\Delta }_y=_y^2`$ and we use the Pochhammer symbol $`(a)_n=\mathrm{\Gamma }(a+n)/\mathrm{\Gamma }(a)`$. In what follows we need to specify explicitly the first three terms of $`C(x,_y)`$ in the derivative expansion:
$`C(x,_y)=1+{\displaystyle \frac{1}{2}}(x_y)+{\displaystyle \frac{\mathrm{\Delta }_O+2}{8(\mathrm{\Delta }_O+1)}}(x_y)^2{\displaystyle \frac{\mathrm{\Delta }_O}{16(\mathrm{\Delta }_O+1)(\mathrm{\Delta }_O+1\eta )}}x^2\mathrm{\Delta }_y+\mathrm{}`$ (6.3)
The conformal blocks of a conserved vector current and a conserved second rank tensor with canonical dimensions $`2\eta 1`$ and $`\eta `$ respectively are also available. For a vector current one has
$`C_\mu (x,_y)`$ $`=`$ $`{\displaystyle \frac{x_\mu }{B(\eta ,\eta )}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m!(\eta )_m}}{\displaystyle _0^1}dt[t(1t)]^{\eta 1}\left({\displaystyle \frac{1}{4}}t(1t)x^2\mathrm{\Delta }_y\right)^me^{tx_y}`$
$`=`$ $`x_\mu +{\displaystyle \frac{1}{2}}x_\mu (x_y)+{\displaystyle \frac{x_\mu }{4(2\eta +1)}}\left((\eta +1)(x_y)^2{\displaystyle \frac{1}{2}}x^2\mathrm{\Delta }_y\right)+\mathrm{}`$
and for a conserved symmetric traceless tensor one finds (see Appendix E)
$`C_{\mu \nu }(x,_y)`$ $`=`$ $`{\displaystyle \frac{x_\mu x_\nu }{B(\eta +1,\eta +1)}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m!(\eta +1)_m}}{\displaystyle _0^1}dt[t(1t)]^\eta \left({\displaystyle \frac{1}{4}}t(1t)x^2\mathrm{\Delta }_y\right)^me^{tx_y}`$
$`=`$ $`x_\mu x_\nu +\mathrm{}`$
Using the above formulae, one can now consider the operator product $`:O^\alpha (x)O^\beta (0):`$ in a free field theory and find explicit expressions for $`J_\mu ^{\alpha \beta }`$ and $`T_{\mu \nu }^{\alpha \beta }`$. Indeed, from the Taylor expansion one sees that the leading component is a quasi-primary field $`O^{\alpha \beta }=:O^\alpha (x)O^\beta (0):`$ with conformal dimension $`\mathrm{\Delta }_O=2\mathrm{\Delta }`$, therefore it should appear in the OPE with its whole conformal block. Subtracting from the Taylor expansion the first three terms of the conformal block of the scalar with dimension $`2\mathrm{\Delta }`$ we find at the next level another quasi-primary operator $`O_\mu ^{\alpha \beta }`$ that turns out to be a vector current $`J_\mu ^{\alpha \beta }=\frac{1}{2}:(_\mu O^\alpha O^\beta O^\alpha _\mu O^\beta ):`$ with dimension $`\mathrm{\Delta }_J=2\mathrm{\Delta }+1`$. Now subtracting from what we get the first two terms of the conformal block of the vector current<sup>9</sup><sup>9</sup>9We do not assume here that $`J_\mu ^{\alpha \beta }`$ is conserved, however the first two terms in the conformal blocks of the conserved and non-conserved vector currents are the same. and decomposing the resulting second rank tensor on the traceless and trace parts we are left with two new fields, one is a tensor and another one is a new scalar, which are given by
$`T_{\mu \nu }^{\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{1}{2}}:(_\mu O^\alpha _\nu O^\beta +_\nu O^\alpha _\mu O^\beta ):{\displaystyle \frac{\mathrm{\Delta }}{2(2\mathrm{\Delta }+1)}}_\mu _\nu (:O^\alpha O^\beta :)`$
$`+`$ $`{\displaystyle \frac{\delta _{\mu \nu }}{4\eta }}({\displaystyle \frac{\mathrm{\Delta }+1}{2\mathrm{\Delta }+1}}^2(:O^\alpha O^\beta :)+:^2O^\alpha O^\beta :+:O^\alpha ^2O^\beta :),`$
$`T^{\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{1}{4\eta }}({\displaystyle \frac{\mathrm{\Delta }\eta +1}{2\mathrm{\Delta }+1\eta }}^2(:O^\alpha O^\beta :)+:^2O^\alpha O^\beta :+:O^\alpha ^2O^\beta :).`$
The transformation properties of these fields under the conformal group show that the are both quasi-primary. Thus, for $`\eta =2`$ we get the desired result (3.6). Note that $`T_{\mu \nu }^{\alpha \beta }`$ is conserved while $`T^{\alpha \beta }`$ vanishes on-shell as soon as $`\eta =\mathrm{\Delta }+1`$. Clearly with the knowledge of the conformal blocks of the higher rank tensor operators the procedure of identifying the quasi-primary operators on the r.h.s of (3.4) may be extended to any desired order.
## 7 Appendix B. Conformal partial wave amplitudes of a scalar, a conserved vector current and the stress tensor
The full contribution of the conformal block of an operator carrying and irreducible representation of the conformal group into the 4-point function is known as the conformal partial wave amplitude (CPWA). The scalar CPWA was computed in by evaluating the corresponding scalar exchange diagram. If we consider operators with the same conformal dimension, then the CPWA of a scalar operator with dimension $`\mathrm{\Delta }_S`$ contributes to its 4-point function as :
$`_S(v,Y)=v^{\frac{\mathrm{\Delta }_S}{2}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{v^n}{n!}}{\displaystyle \frac{\left(\frac{1}{2}\mathrm{\Delta }_S\right)_n^4}{(\mathrm{\Delta }_S)_{2n}(\mathrm{\Delta }_S+1\eta )_n}}{}_{2}{}^{}F_{1}^{}({\displaystyle \frac{1}{2}}\mathrm{\Delta }_S+n,{\displaystyle \frac{1}{2}}\mathrm{\Delta }_S+n;\mathrm{\Delta }_S+2n;Y),`$ (7.1)
where we have represented the result as the convergent series in conformal variables $`v`$ and $`Y`$. The first few terms of the $`v,Y`$ expansion of $`_S(v,Y)`$ are given in (4.2). In particular, for $`\mathrm{\Delta }_S=2`$ the first term of $`v`$-expansion reads as
$`_S(v,Y)={\displaystyle \frac{3}{40}}vF_1(Y)+\mathrm{},`$ (7.2)
where $`F_1(Y)`$ is defined in section 4.2.
The CPWA of traceless symmetric tensors of dimension $`\mathrm{\Delta }`$ and rank $`l`$, corresponding to irreducible representations of dimension $`\mathrm{\Delta }`$ and spin $`l`$ of $`SO(d,2)`$, can be also calculated in CFT as the relevant graphs reduce to sums of scalar exchanges. Using the following normalization prescriptions for the 2- and 3-point functions of the exchanged tensor fields
$`M_{\mu _1,..,\mu _l}(x_1)M_{\nu _1,..,\nu _l}(x_2)`$ $`=`$ $`C_{\mathrm{\Delta },l}{\displaystyle \frac{𝒩(\mathrm{\Delta },l)}{x_{12}^{2\mathrm{\Delta }}}}[\{I_{\mu _1\nu _1}(x_{12})\mathrm{}I_{\mu _l\nu _l}(x_{12})\}_{sym}\mathrm{traces}],`$
$`O(x_1)O(x_3)M_{\mu _1,\mu _2,..,\mu _l}(x_5)`$ $`=`$ $`{\displaystyle \frac{g_{\stackrel{~}{\mathrm{\Delta }}\stackrel{~}{\mathrm{\Delta }}\mathrm{\Delta },l}𝒩(\stackrel{~}{\mathrm{\Delta }};\mathrm{\Delta },l)}{(x_{13}^2)^{\stackrel{~}{\mathrm{\Delta }}\frac{1}{2}\mathrm{\Delta }}(x_{15}^2x_{35}^2)^{\frac{1}{2}\mathrm{\Delta }}}}\left[{\displaystyle \frac{\xi _{\mu _1}\xi _{\mu _2}\mathrm{}\xi _{\mu _l}}{(\xi ^2)^{\frac{1}{2}l}}}\mathrm{trace}\mathrm{terms}\right],`$
where the normalization constants are taken to be
$`𝒩(\mathrm{\Delta },l)={\displaystyle \frac{2^\mathrm{\Delta }\mathrm{\Gamma }(\mathrm{\Delta }+l)\mathrm{\Gamma }(d\mathrm{\Delta }1)}{(2\pi )^{\frac{1}{2}d}\mathrm{\Gamma }(\frac{1}{2}d\mathrm{\Delta })\mathrm{\Gamma }(d\mathrm{\Delta }+l1)}},`$
$`𝒩(\stackrel{~}{\mathrm{\Delta }};\mathrm{\Delta },l)={\displaystyle \frac{2^{\stackrel{~}{\mathrm{\Delta }}+\frac{1}{2}\mathrm{\Delta }+\frac{1}{2}l}}{(2\pi )^{\frac{1}{2}d}}}\left({\displaystyle \frac{\mathrm{\Gamma }(\stackrel{~}{\mathrm{\Delta }}+\frac{1}{2}\mathrm{\Delta }+\frac{1}{2}l\frac{1}{2}d)\mathrm{\Gamma }(\stackrel{~}{\mathrm{\Delta }}\frac{1}{2}\mathrm{\Delta }+\frac{1}{2}l)\mathrm{\Gamma }^2(\frac{1}{2}\mathrm{\Delta }+\frac{1}{2}l)}{\mathrm{\Gamma }(d\stackrel{~}{\mathrm{\Delta }}\frac{1}{2}\mathrm{\Delta }+\frac{1}{2}l)\mathrm{\Gamma }(\frac{1}{2}d\stackrel{~}{\mathrm{\Delta }}+\frac{1}{2}\mathrm{\Delta }+\frac{1}{2}l)\mathrm{\Gamma }^2(\frac{1}{2}d\frac{1}{2}\mathrm{\Delta }+\frac{1}{2}l)}}\right)^{\frac{1}{2}}`$
and
$`\xi _\mu (1,2;3)={\displaystyle \frac{(x_{13})_\mu }{x_{13}^2}}{\displaystyle \frac{(x_{23})_\mu }{x_{23}^2}},\xi ^2(1,2;3)={\displaystyle \frac{x_{12}^2}{x_{13}^2x_{23}^2}},`$
the contribution of the tensor field to the 4-point function of a scalar operator with dimension $`\stackrel{~}{\mathrm{\Delta }}`$ takes the form
$`\beta _{\stackrel{~}{\mathrm{\Delta }}}(x_1,x_3;x_2,x_4;\mathrm{\Delta },l)`$ $`=`$ $`\beta _{\stackrel{~}{\mathrm{\Delta }};\mathrm{\Delta },l}{\displaystyle \frac{1}{(x_{13}^2)^{\stackrel{~}{\mathrm{\Delta }}\frac{1}{2}\mathrm{\Delta }}(x_{24}^2)^{\stackrel{~}{\mathrm{\Delta }}\frac{1}{2}d+\frac{1}{2}\mathrm{\Delta }}}}`$ (7.3)
$`\times {\displaystyle }\mathrm{d}^dx_{\stackrel{~}{5}}{\displaystyle \frac{\{e_{\mu _1}\mathrm{}e_{\mu _l}\mathrm{traces}\}\{e_{\mu _1}^{}\mathrm{}e_{\mu _l}^{}\mathrm{traces}\}}{(x_{1\stackrel{~}{5}}^2x_{3\stackrel{~}{5}}^2)^{\frac{1}{2}\mathrm{\Delta }}(x_{2\stackrel{~}{5}}^2x_{4\stackrel{~}{5}}^2)^{\frac{1}{2}d\frac{1}{2}\mathrm{\Delta }}}}.`$
The constant $`\beta _{\stackrel{~}{\mathrm{\Delta }};\mathrm{\Delta },l}`$ is then given by
$`\beta _{\stackrel{~}{\mathrm{\Delta }};\mathrm{\Delta },l}`$ $`=`$ $`{\displaystyle \frac{g_{\stackrel{~}{\mathrm{\Delta }}\stackrel{~}{\mathrm{\Delta }}\mathrm{\Delta },l}^2}{C_{\mathrm{\Delta },l}}}{\displaystyle \frac{2^{2\stackrel{~}{\mathrm{\Delta }}+\frac{1}{2}d+\frac{1}{2}l}\mathrm{\Gamma }(\stackrel{~}{\mathrm{\Delta }}\frac{1}{2}\mathrm{\Delta }+\frac{1}{2}l)\mathrm{\Gamma }(\stackrel{~}{\mathrm{\Delta }}+\frac{1}{2}\mathrm{\Delta }+\frac{1}{2}l\frac{1}{2}d)}{(2\pi )^{\frac{d}{2}}\mathrm{\Gamma }(\frac{1}{2}d\stackrel{~}{\mathrm{\Delta }}+\frac{1}{2}\mathrm{\Delta }+\frac{1}{2}l)\mathrm{\Gamma }(d\stackrel{~}{\mathrm{\Delta }}\frac{1}{2}\mathrm{\Delta }+\frac{1}{2}l)}},`$
where we have introduced the concise notation
$`e_\mu `$ $`=`$ $`{\displaystyle \frac{\xi _\mu (1,3;\stackrel{~}{5})}{|\xi ^2(1,3;\stackrel{~}{5})|^{\frac{1}{2}}}},e_\mu ^{}={\displaystyle \frac{\xi _\mu (2,4;\stackrel{~}{5})}{|\xi ^2(2,4;\stackrel{~}{5})|^{\frac{1}{2}}}},ee=e^{}e^{}=\mathrm{\hspace{0.17em}\hspace{0.17em}1}.`$
One can show that for the general tensor exchange (7.3) is reduced to a finite sum of four-star integrals $`S(a_1,a_2;a_3,a_4)`$:
$$S(a_1,a_2;a_3,a_4)=\mathrm{d}^4x_{\stackrel{~}{5}}\frac{1}{x_{1\stackrel{~}{5}}^{2a_1}x_{2\stackrel{~}{5}}^{2a_2}x_{3\stackrel{~}{5}}^{2a_3}x_{4\stackrel{~}{5}}^{2a_4}},$$
(7.4)
which can be directly evaluated. The final result is obtained after dropping the “shadow series” of the four-star integral, as the latter corresponds to the exchange of the “shadow tensor” field with dimension $`d\mathrm{\Delta }`$.
Here, we apply the general formula (7.3) to the two cases we are interested in the paper; the case of the conserved vector current with $`\mathrm{\Delta }=d1`$ and $`l=1`$ and the stress tensor with $`\mathrm{\Delta }=d`$ and $`l=2`$. Choosing to work directly in $`d=4`$, the contribution of a conserved vector field in the scalar four-point function is given by
$`\beta _2(x_1,x_2;x_3,x_4;3,1)`$ $`=`$ $`\beta _{2;3,1}{\displaystyle \frac{1}{(x_{12}^2)^{2\frac{3}{2}}(x_{34}^2)^{2\frac{1}{2}}}}{\displaystyle \mathrm{d}^4x_{\stackrel{~}{5}}\frac{ee^{}}{(x_{1\stackrel{~}{5}}^2x_{2\stackrel{~}{5}}^2)^{\frac{3}{2}}(x_{3\stackrel{~}{5}}^2x_{4\stackrel{~}{5}}^2)^{\frac{1}{2}}}}.`$ (7.5)
The inner product $`ee^{}`$ can be written as
$$ee^{}=\frac{1}{2}\left(\frac{x_{1\stackrel{~}{5}}^2x_{2\stackrel{~}{5}}^2x_{3\stackrel{~}{5}}^2x_{4\stackrel{~}{5}}^2}{x_{12}^2x_{34}^2}\right)^{\frac{1}{2}}\left[\frac{x_{24}^2}{x_{2\stackrel{~}{5}}^2x_{4\stackrel{~}{5}}^2}\frac{x_{14}^2}{x_{1\stackrel{~}{5}}^2x_{4\stackrel{~}{5}}^2}+\frac{x_{13}^2}{x_{1\stackrel{~}{5}}^2x_{3\stackrel{~}{5}}^2}\frac{x_{23}^2}{x_{2\stackrel{~}{5}}^2x_{3\stackrel{~}{5}}^2}\right].$$
Substituting (7) into (7.5) we obtain four 4-star functions as
$`\beta _2(x_1,x_2;x_3,x_4;3,1)={\displaystyle \frac{1}{2}}\beta _{2;3,1}{\displaystyle \frac{1}{(x_{12}^2)^{21}(x_{34}^2)^2}}\times `$ (7.6)
$`[x_{24}^2S(1+ϵ,2+ϵ;ϵ,1ϵ)x_{14}^2S(2+ϵ,1+ϵ;ϵ,1ϵ)`$
$`+x_{13}^2S(2+ϵ,1+ϵ;1ϵ,ϵ)x_{23}^2S(1+ϵ,2+ϵ;1ϵ,ϵ)].`$
Note that we have also regularized the dimension of the vector field as $`\mathrm{\Delta }=3+2ϵ`$ to deal with the singularities contained in the four-star functions involved into (7.6). The singularities are avoided by keeping the regulating parameter $`ϵ`$ non-zero in the intermediate stages of the calculation. The analyticity of the exchange graph then ensures that taking the limit $`ϵ0`$ at the end of the calculation one recovers the correct result. Using the expression for the four-star function derived in we then obtain (here we present the formula for general $`d`$ and $`\mathrm{\Delta }`$ to ensure a wider applicability of our result)
$`\beta _{\stackrel{~}{\mathrm{\Delta }}}(x_1,x_2;x_3,x_4;\mathrm{\Delta },1)`$ $`=`$ $`{\displaystyle \frac{2}{3}}\beta _{\stackrel{~}{\mathrm{\Delta }};\mathrm{\Delta },1}{\displaystyle \frac{\pi ^\eta \mathrm{\Gamma }(\frac{\mathrm{\Delta }}{2}\frac{1}{2})\mathrm{\Gamma }(\frac{\mathrm{\Delta }}{2}+\frac{1}{2})\mathrm{\Gamma }(\eta \mathrm{\Delta })}{\mathrm{\Gamma }(\eta \frac{\mathrm{\Delta }}{2}\frac{1}{2})\mathrm{\Gamma }(\eta \frac{\mathrm{\Delta }}{2}+\frac{1}{2})\mathrm{\Gamma }(\mathrm{\Delta })}}{\displaystyle \frac{1}{(x_{12}^2x_{34}^2)^{\stackrel{~}{\mathrm{\Delta }}}}}`$
$`\times `$ $`\left(_V(v,Y)+\text{shadow part}\right),`$
where the function $`_V(v,Y)`$ represents the CPWA of the vector current
$`_V(v,Y)`$ $`=`$ $`{\displaystyle \frac{3}{4}}v^{\frac{\mathrm{\Delta }1}{2}}{\displaystyle \underset{n,m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{v^nY^m}{n!m!}}{\displaystyle \frac{1}{(1\eta +\mathrm{\Delta })_n(\mathrm{\Delta })_{2n+m}}}`$
$`\times `$ $`[\left({\displaystyle \frac{\mathrm{\Delta }+1}{2}}\right)_n^2\left({\displaystyle \frac{\mathrm{\Delta }1}{2}}\right)_{n+m}^2+\left({\displaystyle \frac{\mathrm{\Delta }1}{2}}\right)_n^2\left({\displaystyle \frac{\mathrm{\Delta }+1}{2}}\right)_{n+m}^2`$
$`2\left({\displaystyle \frac{\mathrm{\Delta }1}{2}}\right)_n\left({\displaystyle \frac{\mathrm{\Delta }+1}{2}}\right)_n\left({\displaystyle \frac{\mathrm{\Delta }+1}{2}}\right)_{n+m}\left({\displaystyle \frac{\mathrm{\Delta }1}{2}}\right)_{n+m}`$
$`Y\left({\displaystyle \frac{\mathrm{\Delta }1}{2}}\right)_n^2\left({\displaystyle \frac{\mathrm{\Delta }+1}{2}}\right)_{n+m}^2].`$
For $`\mathrm{\Delta }=d1=3`$ the CPWA of the vector current simplifies to give
$`_V(v,Y)`$ $`=`$ $`{\displaystyle \frac{3}{4}}v{\displaystyle \underset{n,m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{v^nY^m}{n!m!}}{\displaystyle \frac{1}{(2)_n(3)_{2n+m}}}`$
$`\times `$ $`[(2)_n^2(1)_{n+m}^2+(1Y)(1)_n^2(2)_{n+m}^22(1)_n(2)_n(2)_{n+m}(1)_{n+m}],`$
and it is normalized to start as $`_V(v,Y)=\frac{1}{2}Y+\mathrm{}`$ (cf. (4.2)). To make a comparison with the supergravity results in section 4.5 we need to single out in eq.(7) the leading-$`v`$ contribution. Putting in the previous formula $`n=0`$ and performing the summation in $`m`$ we obtain
$`_V(v,Y)`$ $`=`$ $`{\displaystyle \frac{3}{16}}vF_1(Y)+\mathrm{},`$ (7.9)
where $`F_1(Y)`$ is defined in section 4.5.
Analogously, the contribution of the stress tensor is given by
$`\beta _2(x_1,x_2;x_3,x_4;4,2)`$ $`=`$ $`\beta _{2;4,2}{\displaystyle \frac{1}{(x_{34}^2)^2}}{\displaystyle \mathrm{d}^4x_{\stackrel{~}{5}}\frac{(e_\mu e_\nu \frac{1}{4}\delta _{\mu \nu })(e_\mu ^{}e_\nu ^{}\frac{1}{4}\delta _{\mu \nu })}{(x_{1\stackrel{~}{5}}^2x_{2\stackrel{~}{5}}^2)^{2+ϵ}(x_{3\stackrel{~}{5}}^2x_{4\stackrel{~}{5}}^2)^ϵ}}.`$ (7.10)
Using then (7) and regularizing the tensor dimension as $`\mathrm{\Delta }=4+2ϵ`$ we obtain
$`\beta _2(x_1,x_2;x_3,x_4;4,2)`$ $`=`$ $`{\displaystyle \frac{1}{12}}\beta _{2;4,2}{\displaystyle \frac{1}{(x_{34}^2)^2}}[{\displaystyle \frac{1}{x_{12}^2x_{34}^2}}(x_{24}^4S(1+ϵ,3+ϵ;1ϵ,1ϵ)`$ (7.11)
$`+x_{14}^4S(3+ϵ,1+ϵ;1ϵ,1ϵ)+x_{13}^4S(3+ϵ,1+ϵ;1ϵ,1ϵ)`$
$`+x_{23}^4S(1+ϵ,3+ϵ;1ϵ,1ϵ)2x_{24}^2x_{14}^2S(2+ϵ,2+ϵ;1ϵ,1ϵ)`$
$`2x_{13}^2x_{23}^2S(2+ϵ,2+ϵ;1ϵ,1ϵ)+2x_{24}^2x_{13}^2S(2+ϵ,2+ϵ;ϵ,ϵ)`$
$`2x_{24}^2x_{23}^2S(1+ϵ,3+ϵ;ϵ,ϵ)2x_{14}^2x_{13}^2S(3+ϵ,1+ϵ;ϵ,ϵ)`$
$`+2x_{14}^2x_{23}^2S(2+ϵ,2+ϵ;ϵ,ϵ))S(2+ϵ,2+ϵ;ϵ,ϵ)].`$
One then observes that (7.11) contains a number of four-star functions which are $`O(ϵ)`$ and therefore vanish in the $`ϵ0`$ limit. These are all the four-star functions with $`ϵ`$ in the last two positions. Then, by virtue of
$$\frac{\mathrm{\Gamma }(22ϵ)}{\mathrm{\Gamma }(1ϵ)}=\frac{1}{4}+O(ϵ),$$
(7.12)
the remaining four-star functions give a finite result which reads
$`\beta _2(x_1,x_2;x_3,x_4;4,2)`$ $``$ $`{\displaystyle \frac{1}{x_{12}^4x_{34}^4}}\left(_T(v,Y)+\text{shadow part}\right),`$
where $`_T(v,Y)`$ represents the CPWA of the stress tensor:
$`_T(v,Y)`$ $`=`$ $`{\displaystyle \frac{5}{4}}v{\displaystyle \underset{n,m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{v^nY^m}{n!m!}}{\displaystyle \frac{1}{(3)_n(4)_{2n+m}}}`$
$`\times `$ $`[(3)_n^2(1)_{n+m}^2+2(3)_n(1)_n(3)_{n+m}(1)_{n+m}+(1Y)^2(1)_n^2(3)_{n+m}^2`$
$`2(3)_n(2)_n(1)_{n+m}(2)_{n+m}2(1Y)(1)_n(2)_n(3)_{n+m}(2)_{n+m}].`$
The normalization of $`_T(v,Y)`$ is fixed such that its $`v,Y`$ expansion reproduces the corresponding terms in (4.2). Again to establish a link with supergravity results in section 4.1 we single out the $`v`$ term in eq.(7) and, performing the summation in $`m`$, get
$`_T(v,Y)={\displaystyle \frac{45}{8}}vF_1(Y)+\mathrm{},`$ (7.14)
where $`F_1(Y)`$ is the function defined in section 4.1. This completes the construction of the CPWA for conserved vector and tensor currents.
## 8 Appendix C. Series representation for $`\overline{D}`$-functions
Here we derive a representation for the $`\overline{D}_{\mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_4}`$-functions in a form of a convergent series in $`v`$ and $`Y`$ variables by using a technique similar to .
We start with the definition (4.5). Standard Feynman parameter manipulations based on the formula
$$\frac{1}{z^\lambda }=\frac{1}{\mathrm{\Gamma }(\lambda )}_0^{\mathrm{}}dtt^{\lambda 1}e^{tz},$$
and two integrals
$`{\displaystyle e^{_it_ix_0^2}x_0^{d1+{\scriptscriptstyle \mathrm{\Delta }_i}}dx_0}={\displaystyle \frac{1}{2}}(S_t)^{\frac{d_i\mathrm{\Delta }_i}{2}}\mathrm{\Gamma }\left({\displaystyle \frac{\mathrm{\Delta }_i}{2}}{\displaystyle \frac{d}{2}}\right),`$
$`{\displaystyle \mathrm{d}^d\stackrel{}{x}e^{t_i|\stackrel{}{x}\stackrel{}{x}_i|^2}}={\displaystyle \frac{\pi ^{d/2}}{S_t^{d/2}}}e^{\frac{1}{S_t}_{i<j}t_it_jx_{ij}^2},`$
lead to
$`D_{\mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_4}(x_1,x_2,x_3,x_4)`$
$`=K{\displaystyle _0^{\mathrm{}}}\mathrm{d}t_1..\mathrm{d}t_4t_1^{\mathrm{\Delta }_11}\mathrm{}t_4^{\mathrm{\Delta }_41}\left(S_t\right)^{\frac{\mathrm{\Delta }_1+\mathrm{\Delta }_2+\mathrm{\Delta }_3+\mathrm{\Delta }_4}{2}}\mathrm{exp}[{\displaystyle \frac{1}{S_t}}(t_1t_2x_{12}^2+..+t_3t_4x_{34}^2)],`$
where the short-hand notations
$`S_t=t_1+t_2+t_3+t_4,K={\displaystyle \frac{\pi ^{\frac{d}{2}}\mathrm{\Gamma }(\frac{\mathrm{\Delta }_1+\mathrm{\Delta }_2+\mathrm{\Delta }_3+\mathrm{\Delta }_4}{2}\frac{d}{2})}{2\mathrm{\Gamma }(\mathrm{\Delta }_1)\mathrm{}\mathrm{\Gamma }(\mathrm{\Delta }_4)}},`$
were introduced. Performing the change of variables
$$t_i=S_t^{1/2}t_i^{}=(\underset{i}{}t_j^{})t_i^{}ut_i^{},det\left(\frac{t_i}{t_j^{}}\right)=2u^4,$$
one obtains
$`D_{\mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_4}(x_1,x_2,x_3,x_4)=2K{\displaystyle _0^{\mathrm{}}}\mathrm{d}t_1..\mathrm{d}t_4t_1^{\mathrm{\Delta }_11}\mathrm{}t_4^{\mathrm{\Delta }_41}\mathrm{exp}[{\displaystyle \underset{i<j}{}}t_it_jx_{ij}^2].`$
Now we rescale the variables $`t_i`$: $`t_i\lambda _it_i`$, where the constant parameters $`\lambda _i`$ are chosen to induce the following scale transformations
$`t_1t_2{\displaystyle \frac{1}{x_{12}^2}}t_1t_2,t_1t_3{\displaystyle \frac{1}{x_{13}^2}}t_1t_3,t_1t_4{\displaystyle \frac{1}{x_{14}^2}}t_1t_4,t_2t_3{\displaystyle \frac{1}{x_{23}^2}}t_2t_3,`$
and as the consequence
$`t_2t_4={\displaystyle \frac{t_2t_3t_1t_4}{t_1t_3}}{\displaystyle \frac{x_{13}^2}{x_{14}^2x_{23}^2}}t_2t_4,t_3t_4={\displaystyle \frac{t_2t_3t_1t_4}{t_1t_2}}{\displaystyle \frac{x_{12}^2}{x_{14}^2x_{23}^2}}t_3t_4.`$
Under this rescaling the integral transforms into
$`D_{\mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_4}(x_1,x_2,x_3,x_4)=`$
$`={\displaystyle \frac{\overline{D}_{\mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_4}(v,Y)}{(x_{12}^2)^{\frac{\mathrm{\Delta }_1+\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_4}{2}}(x_{13}^2)^{\frac{\mathrm{\Delta }_1+\mathrm{\Delta }_3\mathrm{\Delta }_2\mathrm{\Delta }_4}{2}}(x_{23}^2)^{\frac{\mathrm{\Delta }_2+\mathrm{\Delta }_3+\mathrm{\Delta }_4\mathrm{\Delta }_1}{2}}(x_{14}^2)^{\mathrm{\Delta }_4}}},`$ (8.1)
where
$`\overline{D}_{\mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_4}(v,Y)=`$
$`2K{\displaystyle }\mathrm{d}t_1\mathrm{}\mathrm{d}t_4t_1^{\mathrm{\Delta }_11}t_2^{\mathrm{\Delta }_21}t_3^{\mathrm{\Delta }_31}t_4^{\mathrm{\Delta }_41}\mathrm{exp}[t_1t_2t_1t_3t_1t_4t_2t_3{\displaystyle \frac{v}{u}}t_2t_4vt_3t_4],`$
and the integral is understood as a function of the conformal variables $`v`$ and $`Y`$.
Next, using the Mellin-Barnes integral representation
$`\mathrm{exp}\left[z\right]={\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle _{r\mathrm{i}\mathrm{}}^{r+\mathrm{i}\mathrm{}}}ds\mathrm{\Gamma }(s)z^s,r<0,|\mathrm{arg}z|<{\displaystyle \frac{1}{2}}\pi ,`$
for the two exponentials in the last formula which involve $`\frac{v}{u}`$ and $`v`$ the integral reduces to
$`\overline{D}_{\mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_4}(v,Y)=2K{\displaystyle \frac{\mathrm{d}\lambda \mathrm{d}s}{(2\pi \mathrm{i})^2}\mathrm{\Gamma }(s)\mathrm{\Gamma }(\lambda )v^\lambda \left(\frac{v}{u}\right)^s}`$
$`\times {\displaystyle }\mathrm{d}t_1\mathrm{}\mathrm{d}t_4t_1^{\mathrm{\Delta }_11}t_2^{\mathrm{\Delta }_2+s1}t_3^{\mathrm{\Delta }_3+\lambda 1}t_4^{\mathrm{\Delta }_4+s+\lambda 1}\mathrm{exp}[t_1t_2t_1t_3t_1t_4t_2t_3].`$
The following change of variables:
$$t_1t_2=u_1,t_1t_3=u_2,t_1t_4=u_3,t_2t_3=u_4,det\left(\frac{t_i}{u_j}\right)=\frac{1}{2u_1u_2},$$
allows one to perform the $`t`$-integration with the result
$`\overline{D}_{\mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_4}(v,Y)=`$
$`=K{\displaystyle }{\displaystyle \frac{\mathrm{d}\lambda \mathrm{d}s}{(2\pi \mathrm{i})^2}}[\mathrm{\Gamma }(s)\mathrm{\Gamma }(\lambda )\mathrm{\Gamma }({\displaystyle \frac{\mathrm{\Delta }_1+\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_4}{2}}\lambda )\mathrm{\Gamma }({\displaystyle \frac{\mathrm{\Delta }_1+\mathrm{\Delta }_3\mathrm{\Delta }_2\mathrm{\Delta }_4}{2}}s)`$
$`\times \mathrm{\Gamma }({\displaystyle \frac{\mathrm{\Delta }_2+\mathrm{\Delta }_3+\mathrm{\Delta }_4\mathrm{\Delta }_1}{2}}+s+\lambda )\mathrm{\Gamma }(\mathrm{\Delta }_4+s+\lambda )v^\lambda \left({\displaystyle \frac{v}{u}}\right)^s].`$
The $`s`$-integration is then performed by using the integral and series representations for the hypergeometric function $`F(a,b,c;1z)`$:
$`F(a,b,c;1z)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(c)}{\mathrm{\Gamma }(a)\mathrm{\Gamma }(b)\mathrm{\Gamma }(ca)\mathrm{\Gamma }(cb)}}`$
$`\times `$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _i\mathrm{}^i\mathrm{}}dsz^s\mathrm{\Gamma }(s)\mathrm{\Gamma }(cabs)\mathrm{\Gamma }(a+s)\mathrm{\Gamma }(b+s),`$
and
$`F(a,b,c;1z)={\displaystyle \frac{\mathrm{\Gamma }(c)}{\mathrm{\Gamma }(a)\mathrm{\Gamma }(b)}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }(a+m)\mathrm{\Gamma }(b+m)}{\mathrm{\Gamma }(c+m)m!}}(1z)^m,`$
where one needs to substitute
$$a=\frac{\mathrm{\Delta }_2+\mathrm{\Delta }_3+\mathrm{\Delta }_4\mathrm{\Delta }_1}{2}+\lambda ,b=\mathrm{\Delta }_4+\lambda ,c=\mathrm{\Delta }_3+\mathrm{\Delta }_4+2\lambda .$$
Thus one arrives at the convergent hypergeometric series in the variable $`Y`$:
$`\overline{D}_{\mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_4}(v,Y)=`$
$`=K{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{Y^m}{m!}}\{{\displaystyle }{\displaystyle \frac{\mathrm{d}\lambda }{2\pi \mathrm{i}}}[\mathrm{\Gamma }(\lambda ){\displaystyle \frac{\mathrm{\Gamma }(\frac{\mathrm{\Delta }_1+\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_4}{2}\lambda )\mathrm{\Gamma }(\frac{\mathrm{\Delta }_3+\mathrm{\Delta }_4+\mathrm{\Delta }_1\mathrm{\Delta }_2}{2}+\lambda )\mathrm{\Gamma }(\mathrm{\Delta }_3+\lambda )}{\mathrm{\Gamma }(\mathrm{\Delta }_3+\mathrm{\Delta }_4+2\lambda +m)}}`$
$`\mathrm{\Gamma }({\displaystyle \frac{\mathrm{\Delta }_2+\mathrm{\Delta }_3+\mathrm{\Delta }_4\mathrm{\Delta }_1}{2}}+\lambda +m)\mathrm{\Gamma }(\mathrm{\Delta }_4+\lambda +m)v^\lambda ]\}.`$ (8.2)
Since for any $`\overline{D}`$-function occurring in the 4-point function of CPOs the quantity $`\mathrm{\Delta }_1+\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_4`$ is an integer, the final Mellin-Barnes integral receives a contribution from double poles and, therefore, the integration can be done by using the general formula
$$_𝒞\frac{\mathrm{d}s}{2\pi \mathrm{i}}\mathrm{\Gamma }^2(s)g(s)v^s=\underset{n=0}{\overset{\mathrm{}}{}}\frac{v^n}{(n!)^2}\left[2\psi (n+1)g(n)g(n)\mathrm{ln}v\frac{\mathrm{d}}{\mathrm{d}\xi }[g(\xi )]_{\xi =n}\right],$$
(8.3)
valid for any function $`g(s)`$ regular at $`s=0`$. In this way we arrive at the representation for $`\overline{D}`$-functions in terms of double convergent series in $`v`$ and $`Y`$ variables.
Below we list explicitly the series representations for $`\overline{D}`$-functions we used in the paper
$`\overline{D}_{2222}(v,Y)`$ $`=`$ $`\pi ^2{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{Y^m}{m!}}{\displaystyle \frac{v^n}{(n!)^2}}{\displaystyle \frac{\mathrm{\Gamma }(n+2)^2\mathrm{\Gamma }(2+n+m)^2}{\mathrm{\Gamma }(4+2n+m)}}`$
$`\times `$ $`\left({\displaystyle \frac{1}{n+1}}+\psi (4+2n+m)\psi (n+m+2){\displaystyle \frac{1}{2}}\mathrm{ln}v\right),`$
$`\overline{D}_{2112}(v,Y)`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{2}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{Y^m}{m!}}{\displaystyle \frac{v^n}{(n!)^2}}{\displaystyle \frac{\mathrm{\Gamma }(n+2)\mathrm{\Gamma }(n+1)\mathrm{\Gamma }(n+m+1)\mathrm{\Gamma }(n+m+2)}{\mathrm{\Gamma }(3+2n+m)}}`$
$`\times `$ $`\left({\displaystyle \frac{1}{n+1}}+2\psi (3+2n+m)\psi (n+m+1)\psi (n+m+2)\mathrm{ln}v\right),`$
$`\overline{D}_{1212}(v,Y)`$ $`=`$ $`\pi ^2{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{Y^m}{m!}}{\displaystyle \frac{v^n}{(n!)^2}}{\displaystyle \frac{\mathrm{\Gamma }(n+1)^2\mathrm{\Gamma }(n+m+2)^2}{\mathrm{\Gamma }(3+2n+m)}}`$
$`\times `$ $`\left(\psi (3+2n+m)\psi (n+m+2){\displaystyle \frac{1}{2}}\mathrm{ln}v\right),`$
$`\overline{D}_{2211}(v,Y)`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{2}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{Y^m}{m!}}{\displaystyle \frac{v^n}{(n!)^2}}{\displaystyle \frac{n\mathrm{\Gamma }(n+1)^2\mathrm{\Gamma }(n+m+1)^2}{\mathrm{\Gamma }(2+2n+m)}}`$
$`\times `$ $`\left({\displaystyle \frac{1}{n}}2\psi (n+m+1)+2\psi (2+2n+m)\mathrm{ln}v\right),`$
$`\overline{D}_{3322}(v,Y)`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{4}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{Y^m}{m!}}{\displaystyle \frac{v^n}{(n!)^2}}{\displaystyle \frac{n\mathrm{\Gamma }(n+2)^2\mathrm{\Gamma }(2+n+m)^2}{\mathrm{\Gamma }(4+2n+m)}}`$
$`\times `$ $`\left({\displaystyle \frac{3n+1}{n(n+1)}}+2\psi (4+2n+m)2\psi (2+n+m)\mathrm{ln}v\right),`$
$`\overline{D}_{2323}(v,Y)`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{2}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{Y^m}{m!}}{\displaystyle \frac{v^n}{(n!)^2}}{\displaystyle \frac{\mathrm{\Gamma }(n+2)^2\mathrm{\Gamma }(3+n+m)^2}{\mathrm{\Gamma }(5+2n+m)}}`$
$`\times `$ $`\left({\displaystyle \frac{1}{n+1}}+\psi (5+2n+m)\psi (3+n+m){\displaystyle \frac{1}{2}}\mathrm{ln}v\right),`$
$`\overline{D}_{3223}(v,Y)`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{4}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{Y^m}{m!}}{\displaystyle \frac{v^n}{(n!)^2}}{\displaystyle \frac{\mathrm{\Gamma }(n+2)\mathrm{\Gamma }(n+3)\mathrm{\Gamma }(2+n+m)\mathrm{\Gamma }(3+n+m)}{\mathrm{\Gamma }(5+2n+m)}}`$ (8.4)
$`\times `$ $`({\displaystyle \frac{3n+5}{(n+1)(n+2)}}+2\psi (5+2n+m)`$
$``$ $`\psi (2+n+m)\psi (3+n+m)\mathrm{ln}v).`$
## 9 Appendix D. Projectors
Here we give an explicit construction of the projectors that single out the contributions of irreps occurring in the decomposition $`\mathrm{𝟐𝟎}\times \mathrm{𝟐𝟎}`$ of $`SO(6)`$ from the 4-point function of the lowest weight CPOs.
Matrices $`C_{ij}^I`$ and $`C_{ij}^{𝒥_{15}}`$ introduced in section 3 obey to the following summation formulae :
$`{\displaystyle \underset{I}{}}C_{ij}^IC_{kl}^I={\displaystyle \frac{1}{2}}\delta _{ik}\delta _{jl}+{\displaystyle \frac{1}{2}}\delta _{il}\delta _{jk}{\displaystyle \frac{1}{6}}\delta _{ij}\delta _{kl},{\displaystyle \underset{𝒥_{15}}{}}C_{ij}^{𝒥_{15}}C_{kl}^{𝒥_{15}}={\displaystyle \frac{1}{2}}(\delta _{ik}\delta _{jl}\delta _{il}\delta _{jk}).`$
It is then easy to check that the orthonormal Clebsh-Gordon coefficients $`C_{𝒥_{20}}^{I_1I_2}`$ and $`C_{𝒥_{15}}^{I_1I_2}`$ are given by
$`C_{𝒥_{20}}^{I_1I_2}={\displaystyle \frac{3^{1/2}}{5^{1/2}}}C^{I_1I_2I},C_{𝒥_{15}}^{I_1I_2}={\displaystyle \frac{1}{2^{1/2}}}C_{ij}^{I_1}C_{jk}^{I_2}C_{ik}^{𝒥_{15}}.`$ (9.1)
The other coefficients are constructed in a similar manner. Irreps $`\mathrm{𝟖𝟒}`$, $`\mathrm{𝟏𝟎𝟓}`$ and $`\mathrm{𝟏𝟕𝟓}`$ are described by traceless rank 4 tensors $`C_{ijkl}^{𝒥_D}`$ with the normalization condition
$$C_{ijkl}^{𝒥_D}C_{ijkl}^{𝒥_D^{}}=\delta ^{𝒥_D𝒥_D^{}}.$$
Tensor $`C_{ijkl}^{𝒥_{84}}`$ is antisymmetric in $`i,k`$ and in $`j,l`$ and symmetric under permutation of the pairs $`ij`$ and $`kl`$. It is also required to obey the condition $`\epsilon _{ijklmn}C_{klmn}^{𝒥_{84}}=0`$. Then $`C_{ijkl}^{𝒥_{105}}`$ is a totally symmetric and, finally, $`C_{ijkl}^{𝒥_{175}}`$ is symmetric in $`i,k`$ and in $`j,l`$ and antisymmetric under permutation of the pairs $`ij`$ and $`kl`$.
A projector on the contribution of irrep $`𝐃`$ into the 4-point function is defined by (4.15) with $`\nu _D`$ being the dimension of the representation. The sums $`C_{𝒥_{20}}^{I_1I_2}C_{𝒥_{20}}^{I_1I_2}`$ and $`C_{𝒥_{15}}^{I_1I_2}C_{𝒥_{15}}^{I_1I_2}`$ are computed straightforwardly by using eqs.(9.1). To find the other projectors we introduced the following three tensors $`Q_𝐃^{I_1I_2}`$ being elements of the corresponding representations:
$`(Q_{\mathrm{𝟖𝟒}}^{I_1I_2})_{ijkl}`$ $`=`$ $`C_{ij}^{I_1}C_{kl}^{I_2}C_{kj}^{I_1}C_{il}^{I_2}+C_{kl}^{I_1}C_{ij}^{I_2}C_{il}^{I_1}C_{kj}^{I_2}`$
$`+`$ $`{\displaystyle \frac{1}{4}}(C_{im}^{I_1}C_{mj}^{I_2}\delta _{kl}C_{km}^{I_1}C_{mj}^{I_2}\delta _{il}+C_{km}^{I_1}C_{ml}^{I_2}\delta _{ij}C_{im}^{I_1}C_{ml}^{I_2}\delta _{kj})`$
$`+`$ $`{\displaystyle \frac{1}{4}}(C_{jm}^{I_1}C_{mi}^{I_2}\delta _{kl}C_{jm}^{I_1}C_{mk}^{I_2}\delta _{il}+C_{lm}^{I_1}C_{mk}^{I_2}\delta _{ij}C_{lm}^{I_1}C_{mi}^{I_2}\delta _{kj})`$
$``$ $`{\displaystyle \frac{1}{10}}\delta ^{I_1I_2}(\delta _{ij}\delta _{kl}\delta _{il}\delta _{kj}),`$
$`(Q_{\mathrm{𝟏𝟎𝟓}}^{I_1I_2})_{ijkl}`$ $`=`$ $`C_{ij}^{I_1}C_{kl}^{I_2}+C_{ik}^{I_1}C_{jl}^{I_2}+C_{il}^{I_1}C_{jk}^{I_2}+C_{kl}^{I_1}C_{ij}^{I_2}+C_{kj}^{I_1}C_{il}^{I_2}+C_{jl}^{I_1}C_{ik}^{I_2}`$
$``$ $`{\displaystyle \frac{1}{5}}\delta _{ij}(C_{km}^{I_1}C_{ml}^{I_2}+C_{lm}^{I_1}C_{mk}^{I_2}){\displaystyle \frac{1}{5}}\delta _{kl}(C_{im}^{I_1}C_{mj}^{I_2}+C_{jm}^{I_1}C_{mi}^{I_2})`$
$``$ $`{\displaystyle \frac{1}{5}}\delta _{ik}(C_{jm}^{I_1}C_{ml}^{I_2}+C_{lm}^{I_1}C_{mj}^{I_2}){\displaystyle \frac{1}{5}}\delta _{il}(C_{jm}^{I_1}C_{mk}^{I_2}`$
$`+`$ $`C_{km}^{I_1}C_{mj}^{I_2}){\displaystyle \frac{1}{5}}\delta _{jk}(C_{im}^{I_1}C_{ml}^{I_2}+C_{lm}^{I_1}C_{mi}^{I_2}){\displaystyle \frac{1}{5}}\delta _{jl}(C_{im}^{I_1}C_{mk}^{I_2}+C_{km}^{I_1}C_{mi}^{I_2})`$
$`+`$ $`{\displaystyle \frac{1}{20}}(\delta _{ij}\delta _{kl}+\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk})\delta ^{I_1I_2},`$
$`(Q_{\mathrm{𝟏𝟕𝟓}}^{I_1I_2})_{ijkl}`$ $`=`$ $`C_{ik}^{I_1}C_{jl}^{I_2}C_{jl}^{I_1}C_{ik}^{I_2}{\displaystyle \frac{1}{8}}\delta _{ij}(C_{km}^{I_1}C_{ml}^{I_2}C_{lm}^{I_1}C_{mk}^{I_2}){\displaystyle \frac{1}{8}}\delta _{kj}(C_{im}^{I_1}C_{ml}^{I_2}C_{lm}^{I_1}C_{mi}^{I_2})`$
$``$ $`{\displaystyle \frac{1}{8}}\delta _{il}(C_{km}^{I_1}C_{mj}^{I_2}C_{jm}^{I_1}C_{mk}^{I_2}){\displaystyle \frac{1}{8}}\delta _{kl}(C_{im}^{I_1}C_{mj}^{I_2}C_{jm}^{I_1}C_{mi}^{I_2}).`$
Clearly one may write
$`C_{𝒥_D}^{I_1I_2}C_{ijkl}^{𝒥_D}=\gamma _D(Q_𝐃^{I_1I_2})_{ijkl},`$
where $`\gamma _D`$ is a normalization constant. Then one finds
$`C_{𝒥_D}^{I_1I_2}C_{𝒥_D}^{I_3I_4}=C_{𝒥_D}^{I_1I_2}C_{ijkl}^{𝒥_D}C_{𝒥_D^{}}^{I_3I_4}C_{ijkl}^{𝒥_D^{}}=\gamma _D^2(Q_𝐃^{I_1I_2})_{ijkl}(Q_𝐃^{I_3I_4})_{ijkl}`$
with the normalization constant $`\gamma _D`$ following from
$$\nu _D=\gamma _D^2(Q_𝐃^{I_1I_2})_{ijkl}(Q_𝐃^{I_1I_2})_{ijkl}$$
and, therefore,
$$(P_𝐃)_{I_1I_2I_3I_4}=\frac{(Q_𝐃^{I_1I_2})_{ijkl}(Q_𝐃^{I_3I_4})_{ijkl}}{(Q_𝐃^{IJ})_{mnsp}(Q_𝐃^{IJ})_{mnsp}}.$$
In this way one obtains the following explicit expressions for projectors singling out the contributions of the irreps:
$`(P_{\mathrm{𝟏𝟓}})_{I_1I_2I_3I_4}`$ $`=`$ $`{\displaystyle \frac{1}{30}}C_{I_1I_2I_3I_4}^{},`$
$`(P_{\mathrm{𝟐𝟎}})_{I_1I_2I_3I_4}`$ $`=`$ $`{\displaystyle \frac{3}{100}}\left(C_{I_1I_2I_3I_4}^+{\displaystyle \frac{1}{6}}\delta _{I_1I_2}\delta _{I_3I_4}\right),`$
$`(P_{\mathrm{𝟖𝟒}})_{I_1I_2I_3I_4}`$ $`=`$ $`{\displaystyle \frac{1}{504}}\left(2\delta _{I_1I_3}\delta _{I_2I_4}+2\delta _{I_1I_4}\delta _{I_2I_3}+{\displaystyle \frac{1}{5}}\delta _{I_1I_2}\delta _{I_3I_4}4C_{I_1I_3I_2I_4}2C_{I_1I_2I_3I_4}^+\right),`$
$`(P_{\mathrm{𝟏𝟎𝟓}})_{I_1I_2I_3I_4}`$ $`=`$ $`{\displaystyle \frac{1}{1260}}\left(2\delta _{I_1I_3}\delta _{I_2I_4}+2\delta _{I_1I_4}\delta _{I_2I_3}+{\displaystyle \frac{1}{5}}\delta _{I_1I_2}\delta _{I_3I_4}+8C_{I_1I_3I_2I_4}{\displaystyle \frac{16}{5}}C_{I_1I_2I_3I_4}^+\right),`$
$`(P_{\mathrm{𝟏𝟕𝟓}})_{I_1I_2I_3I_4}`$ $`=`$ $`{\displaystyle \frac{1}{350}}\left(\delta _{I_1I_3}\delta _{I_2I_4}\delta _{I_1I_4}\delta _{I_2I_3}+C_{I_1I_2I_3I_4}^{}\right).`$
Defining also $`P_\mathrm{𝟏}=\frac{1}{400}\delta _{I_1I_2}\delta _{I_3I_4}`$, one can verify that the tensors $`P_\mathrm{𝟏}`$,…,$`P_{\mathrm{𝟏𝟕𝟓}}`$ project an arbitrary 4-point function of operators in the irrep 20 into the corresponding six subspaces of the direct sum decomposition $`\mathrm{𝟐𝟎}\times \mathrm{𝟐𝟎}=\mathrm{𝟏}+\mathrm{𝟐𝟎}+\mathrm{𝟖𝟒}+\mathrm{𝟏𝟎𝟓}+\mathrm{𝟏𝟓}+\mathrm{𝟏𝟕𝟓}`$.
The following formulae
$`C_{I_1I_2I_3I_4}C_{I_1I_2I_3I_4}={\displaystyle \frac{380}{3}},C_{I_1I_2I_3I_4}C_{I_2I_1I_3I_4}={\displaystyle \frac{20}{3}},C_{I_1I_2I_3I_4}^{}C_{I_1I_3I_2I_4}=0,`$
$`C_{I_1I_2I_3I_4}^+C_{I_1I_2I_3I_4}^+={\displaystyle \frac{200}{3}},C_{I_1I_2I_3I_4}^+C_{I_1I_2I_3I_4}^{}=0,C_{I_1I_2I_3I_4}^+C_{I_1I_3I_2I_4}={\displaystyle \frac{20}{3}}.`$
are helpful to find the contractions of the projectors with tensors describing the 4-point function. The results for contractions are summarized in the Table 1.
| Tensor | $`C_{I_1I_2I_3I_4}^+`$ | $`C_{I_1I_2I_3I_4}^{}`$ | $`C_{I_1I_3I_2I_4}`$ | $`\delta _{I_1I_3}\delta _{I_2I_4}`$ | $`\delta _{I_1I_4}\delta _{I_2I_3}`$ |
| --- | --- | --- | --- | --- | --- |
| $`\frac{1}{400}\delta _{I_1I_2}\delta _{I_3I_4}`$ | $`\frac{1}{6}`$ | 0 | $`\frac{1}{60}`$ | $`\frac{1}{20}`$ | $`\frac{1}{20}`$ |
| $`(P_{\mathrm{𝟏𝟓}})_{I_1I_2I_3I_4}`$ | 0 | -2 | 0 | 1 | -1 |
| $`(P_{\mathrm{𝟐𝟎}})_{I_1I_2I_3I_4}`$ | $`\frac{5}{3}`$ | 0 | $`\frac{1}{6}`$ | 1 | 1 |
| $`(P_{\mathrm{𝟖𝟒}})_{I_1I_2I_3I_4}`$ | 0 | 0 | $`\frac{1}{2}`$ | 1 | 1 |
| $`(P_{\mathrm{𝟏𝟎𝟓}})_{I_1I_2I_3I_4}`$ | 0 | 0 | 1 | 1 | 1 |
| $`(P_{\mathrm{𝟏𝟕𝟓}})_{I_1I_2I_3I_4}`$ | 0 | 0 | 0 | 1 | -1 |
Table 1. The values of contractions of the projectors with tensors describing the structure of the 4-point function of the lowest weight CPOs.
## 10 Appendix E. Conformal block of the conserved 2nd rank tensor
Here we sketch the derivation of the conformal block of the conserved second rank tensor. We do not use this result in the paper, however, we feel that it might be useful for subsequent studies of the OPE.
We start with (6) and suppress the unessential indices $`\alpha `$ and $`\beta `$. The 3-point function is given by the following expression
$`O(x)O(0)T_{\mu \nu }(y)={\displaystyle \frac{1}{(x^2)^{\frac{1}{2}(2\mathrm{\Delta }\mathrm{\Delta }_T+2)}(y^2)^{\frac{1}{2}(\mathrm{\Delta }_T+2)}((yx)^2)^{\frac{1}{2}(\mathrm{\Delta }_T+2)}}}[{\displaystyle \frac{\delta _{\mu \nu }}{d}}x^2y^2(yx)^2`$
$`y^4(yx)_\mu (yx)_\nu (xy)^4y_\mu y_\nu +y^2(yx)^2((yx)_\mu y_\nu +y_\mu (yx)_\nu )],`$
where for simplicity we choose the constant $`C_{OOT}`$ to be equal to unity. Compatibility of the 3-point function with the conservation law requires the dimension of the tensor to be canonical, i.e. $`\mathrm{\Delta }_T=d`$, where $`d=2\eta `$ is a space-time dimension. However, in what follows we meet certain divergences and that is why we keep in some places $`\epsilon =d\mathrm{\Delta }_T`$ as a regularization parameter. Substituting eq.(6) into the 3-point function, we get an equation defining the conformal block
$`{\displaystyle \frac{1}{\mathrm{\Delta }_T(\mathrm{\Delta }_T2)}}\left({\displaystyle \frac{1}{(y^2)^{\frac{1}{2}(\mathrm{\Delta }_T2)}}}e^{x_y}_\mu _\nu {\displaystyle \frac{1}{(y^2)^{\frac{1}{2}(\mathrm{\Delta }_T2)}}}+_\mu _\nu {\displaystyle \frac{1}{(y^2)^{\frac{1}{2}(\mathrm{\Delta }_T2)}}}e^{x_y}{\displaystyle \frac{1}{(y^2)^{\frac{1}{2}(\mathrm{\Delta }_T2)}}}\right)`$
$`+{\displaystyle \frac{1}{(\mathrm{\Delta }_T2)^2}}\left(_\mu {\displaystyle \frac{1}{(y^2)^{\frac{1}{2}(\mathrm{\Delta }_T2)}}}e^{x_y}_\nu {\displaystyle \frac{1}{(y^2)^{\frac{1}{2}(\mathrm{\Delta }_T2)}}}+_\nu {\displaystyle \frac{1}{(y^2)^{\frac{1}{2}(\mathrm{\Delta }_T2)}}}e^{x_y}_\mu {\displaystyle \frac{1}{(y^2)^{\frac{1}{2}(\mathrm{\Delta }_T2)}}}\right)`$
$`{\displaystyle \frac{d\mathrm{\Delta }_T}{d\mathrm{\Delta }_T}}\delta _{\mu \nu }\left({\displaystyle \frac{1}{(y^2)^{\frac{1}{2}(\mathrm{\Delta }_T2)}}}e^{x_y}{\displaystyle \frac{1}{(y^2)^{\frac{1}{2}\mathrm{\Delta }_T}}}+{\displaystyle \frac{1}{(y^2)^{\frac{1}{2}\mathrm{\Delta }_T}}}e^{x_y}{\displaystyle \frac{1}{(y^2)^{\frac{1}{2}(\mathrm{\Delta }_T2)}}}\right)`$
$`+{\displaystyle \frac{2\delta _{\mu \nu }}{d(\mathrm{\Delta }_T2)^2}}_\lambda {\displaystyle \frac{1}{(y^2)^{\frac{1}{2}(\mathrm{\Delta }_T2)}}}e^{x_y}^\lambda {\displaystyle \frac{1}{(y^2)^{\frac{1}{2}(\mathrm{\Delta }_T2)}}}={\displaystyle \underset{k=2}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{k!}}\mathrm{\Delta }_{\rho \lambda }^k(x,_y)T_{\rho \lambda }(0)T_{\mu \nu }(y),`$
where we have introduced the following representation
$$C_{\mu \nu }(x,_y)=\underset{k=2}{\overset{\mathrm{}}{}}\frac{1}{k!}\mathrm{\Delta }_{\rho \lambda }^k(x,_y).$$
Using the series representation for the exponentials on the l.h.s. of the equation defining the conformal block, one then obtains an equation for $`\mathrm{\Delta }_{\rho \lambda }^k(x,_y)`$.
The 2-point function of the second rank tensor can be written in the form
$`T_{\rho \lambda }(0)T_{\mu \nu }(y)`$ $`=`$ $`{\displaystyle \frac{1}{(\mathrm{\Delta }_T3)(\mathrm{\Delta }_T2)\mathrm{\Delta }_T(\mathrm{\Delta }_T+1)}}_{\mu \nu \rho \lambda ;\alpha \beta \gamma \delta }_\alpha _\beta _\gamma _\delta {\displaystyle \frac{1}{(y^2)^{\mathrm{\Delta }_T2}}},`$
where again for simplicity we choose the constant $`C_T`$ to be the unity. Here $`_{\mu \nu \rho \lambda ;\alpha \beta \gamma \delta }`$ is a tensor with the following structure
$$_{\mu \nu \rho \lambda ;\alpha \beta \gamma \delta }_\alpha _\beta _\gamma _\delta =\frac{\mathrm{\Delta }_T3}{4(\mathrm{\Delta }_T2)}\left(\frac{1}{2}\text{ }\text{ }\text{ }\text{ }^2(\delta _{\mu \rho }\delta _{\nu \lambda }+\delta _{\mu \rho }\delta _{\nu \lambda })+\delta _{\rho \lambda }(\mathrm{})+\text{longitudinal}\right).$$
If we suppose that the conformal block acting on the 2-point function is symmetric traceless and transversal, then only the first term here is of importance and one gets
$`\mathrm{\Delta }_{\rho \lambda }^k(x,_y)T_{\rho \lambda }(0)T_{\mu \nu }(y)={\displaystyle \frac{4(\mathrm{\Delta }_T1)(\mathrm{\Delta }_T1\eta )(\mathrm{\Delta }_T\eta )}{(\mathrm{\Delta }_T2)\mathrm{\Delta }_T(\mathrm{\Delta }_T+1)}}\mathrm{\Delta }_{\mu \nu }^k(x,_y){\displaystyle \frac{1}{(y^2)^{\mathrm{\Delta }_T}}}.`$
Now we substitute every function $`\frac{1}{(y^2)^a}`$ appearing in the equation defining the conformal block for its Fourier transform
$`{\displaystyle \frac{1}{(y^2)^a}}=2^{2(\eta a)}\pi ^\eta {\displaystyle \frac{\mathrm{\Gamma }(\eta a)}{\mathrm{\Gamma }(a)}}{\displaystyle \frac{1}{(2\pi )^{2\eta }}}{\displaystyle dp\frac{e^{ipy}}{(p^2)^{\eta a}}}`$
and find<sup>10</sup><sup>10</sup>10We keep $`\epsilon `$ only there where it is actually needed to compute the limit $`\epsilon 0`$.
$`\mathrm{\Delta }_{\mu \nu }^k(x,ip)`$ $`=`$ $`{\displaystyle \frac{2^4}{\pi ^\eta }}{\displaystyle \frac{\mathrm{\Gamma }(2\eta )(2\eta +1)}{\mathrm{\Gamma }^2(\eta 1)\mathrm{\Gamma }(2\eta )}}(p^2)^\eta {\displaystyle dq\frac{(p_\mu q_\mu )(p_\nu q_\mu )}{((pq)^2q^2)^{1+\epsilon /2}}(ixq)^k}.`$
Since the conformal block is applied to the traceless transversal operator (2-point function) in the last expression we have omitted all trace and longitudinal terms proportional to $`\delta _{\mu \nu }`$ and to $`p_\mu `$ respectively. The equation can be then brought to the form
$`\mathrm{\Delta }_{\mu \nu }^k(x,ip)`$ $`=`$ $`{\displaystyle \frac{4}{\pi ^\eta }}{\displaystyle \frac{\mathrm{\Gamma }(2\eta )(2\eta +1)}{\mathrm{\Gamma }^2(\eta 1)\mathrm{\Gamma }(2\eta )}}{\displaystyle \frac{(p^2)^\eta }{\frac{\epsilon }{2}\left(\frac{\epsilon }{2}1\right)}}_\mu _\nu I_k({\displaystyle \frac{\epsilon }{2}}1;{\displaystyle \frac{\epsilon }{2}}+1),`$ (10.1)
again modulo unessential trace and longitudinal terms. Here we introduced the following integral
$`I_k(\alpha _1;\alpha _2)={\displaystyle dq\frac{(ixq)^k}{((pq)^2)^{\alpha _1}(q^2)^{\alpha _2}}}`$
that is explicitly evaluated to give
$`I_k(\alpha _1;\alpha _2)`$ $`=`$ $`{\displaystyle \frac{\pi ^\eta }{\mathrm{\Gamma }(\alpha _1)\mathrm{\Gamma }(\alpha _2)}}{\displaystyle \underset{m=0}{\overset{[k/2]}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{k}{2m}}\right){\displaystyle \frac{(2m)!}{m!}}(ixp)^{k2m}\left({\displaystyle \frac{1}{4}}x^2p^2\right)^m(p^2)^{\eta \alpha _1\alpha _2}`$
$`\times `$ $`{\displaystyle \frac{\mathrm{\Gamma }(km+\eta \alpha _2)\mathrm{\Gamma }(m+\eta \alpha _1)}{\mathrm{\Gamma }(k+2\eta \alpha _1\alpha _2)}}\mathrm{\Gamma }(\alpha _1+\alpha _2m\eta ).`$
Again neglecting the trace and longitudinal contributions, we evaluate the limit $`\epsilon 0`$ and normalize the resulting expression such that the first nontrivial term $`\mathrm{\Delta }_{\mu \nu }^2`$ starts as $`\mathrm{\Delta }_{\mu \nu }^2(x,ip)=2x_\mu x_\nu +\mathrm{}`$. In this way we find the following expression
$`\mathrm{\Delta }_{\mu \nu }^k(x,_y)`$ $`=`$ $`x_\mu x_\nu {\displaystyle \frac{\mathrm{\Gamma }(2\eta +2)}{\mathrm{\Gamma }(\eta +1)}}`$
$`\times `$ $`{\displaystyle \underset{m=0}{\overset{[k/2]1}{}}}{\displaystyle \frac{k!}{(k2m2)!m!}}{\displaystyle \frac{\mathrm{\Gamma }(km+\eta 1)}{\mathrm{\Gamma }(k+2\eta )}}(x_y)^{k2m2}\left({\displaystyle \frac{1}{4}}x^2\mathrm{\Delta }_y\right)^m.`$
Finally, performing the summation in $`k`$ we recover the expression $`C_{\mu \nu }(x,_y)`$ for the conformal block of the conserved second rank tensor given in Appendix A.
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# 1 Introduction
## 1 Introduction
The existence of a layered phase in gauge theories with anisotropic coupling has been first conjectured by Fu and Nielsen (, see also ) in the early eighties. If the couplings are equal in all directions, the Wilson loops are governed by the area law (if the couplings are strong) or by the perimeter law (if they are weak). If now the couplings are weak in d dimensions and strong in the remaining one dimension, then the model will exhibit Coulomb forces in the d-dimensional subspace and will confine in the last direction; this means that charged particles will be localized within the d–dimensional subspace. In the extreme case, where the coupling in the extra dimension is infinite, charged particles from each subspace will remain strictly confined to the d–dimensional space–time; in other words, the d-dimensional hyperplanes(layers) will decouple from each other and this is at the origin of the characterization of this phase as “layered”.
Monte Carlo evidence for a layered structure in pure U(1) has been given later on . There exists a lowest dimension for which the layered phase exists: it is the dimension above which the system has two phases. For the abelian gauge theory it is at five dimensions that one may first discriminate between the Confinement and the Coulomb phases; for the non-abelian case the critical dimension is six.
The initial motivation for the study of a model with anisotropic couplings has been dimensional reduction with cosmological implications: if someone lives in one of the d-dimensional hyperplanes of a system in the layered phase, he has no way of communicating with the remaining layers and a kind of dimensional reduction is achieved. These considerations have also been given a modern version in the context of membrane theories and confining (or localization) of particles within the branes . On the other hand, the construction of the chiral layer is an important ingredient of the models of chiral fermions on the lattice .
In this paper we try to perform an exploratory study of a richer model. We treat a (2+1)-dimensional Abelian Higgs model. This has the advantage that it is a low-dimensional model with a phase transition (at large gauge couplings), so it may have a layered phase. In addition, this system has a direct relation to the ceramics exhibiting high–$`T_c`$ superconductivity . These materials have a layered structure, so they are effectively planar and high–$`T_c`$ superconductivity is considered as a two dimensional phenomenon. In this sense the $`(2+1)`$ – dimensional theory acquires physical relevance. The isotropic version of the model has been extensively studied recently , in connection with cosmological considerations and Condensed Matter Physics, in which external magnetic fields play a predominant role, such as penetration and vortex formation.
Another feature of the various phase transitions associated with the model is that they may be strengthened in the presence of an external magnetic field; in addition, the phases of the model with compact action may be characterized by the behaviour of the system under the influence of the magnetic field.
Both the compact and non–compact versions of the model have been examined in this paper and they both have characteristics that one can attribute to layer decoupling.
## 2 Formulation of the model
The model under study is the Abelian Higgs model in the three-dimensional space. Direction $`\widehat{3},`$ corresponding to $`z,`$ will be singled out by couplings that will differ from the rest; later on an external magnetic field will also point in this direction.
We first discuss the non-compact version of the model; we proceed with writing down its lattice action.
$$S=\frac{1}{2}\beta _{gs}\underset{x}{}F_{12}^2(x)+\frac{1}{2}\beta _{gt}\underset{x}{}[F_{13}^2(x)+F_{23}^2(x)]$$
$$+\frac{1}{2}\beta _{hs}(\underset{x}{}[2\phi ^{}(x)\phi (x)\phi ^{}(x)U_{\widehat{1}}(x)\phi (x+\widehat{1})\phi ^{}(x)U_{\widehat{2}}(x)\phi (x+\widehat{2})]+h.c.)$$
$$+\frac{1}{2}\beta _{ht}(\underset{x}{}[\phi ^{}(x)\phi (x)\phi ^{}(x)U_{\widehat{3}}(x)\phi (x+\widehat{3})]+h.c.)$$
$$+\underset{x}{}[(12\beta _R2\beta _{hs}\beta _{ht})\phi ^{}(x)\phi (x)+\beta _R(\phi ^{}(x)\phi (x))^2]$$
(1)
where $`F_{ij}A_j(x+\widehat{i})A_j(x)A_i(x+\widehat{j})+A_j(x).`$
We have allowed for different couplings in the various directions: the ones pertaining to directions $`\widehat{1},\widehat{2}`$ are given the subscript “s” for “space”, while direction $`\widehat{3}`$ gets couplings with the subscript “t” for “time”. We stress that the third direction is not really “time”; it is just the direction singled out as explained previously. However, we use this notation for convenience; in particular, we refer below to “time-like” plaquettes, links or couplings with this understanding.
For future use we mention here some notations. The link variables $`U_{\widehat{k}}(x)`$ are defined as $`e^{ig_s\alpha _sA_s}`$ or $`e^{ig_t\alpha _tA_t}`$ respectively. They are also written in the form $`U_{\widehat{k}}(x)=e^{i\theta _{\widehat{k}}(x)},`$ since they are complex phases. In addition, the scalar fields are also written in the polar form $`\phi (x)=R(x)e^{i\chi (x)}.`$ The order parameters that we will use are the following:
$$\mathrm{Space}\mathrm{like}\mathrm{Plaquette}:P_s<\frac{1}{N^3}\underset{x}{}F_{12}^2(x)>$$
(2)
$$\mathrm{Time}\mathrm{like}\mathrm{Plaquette}:P_t<\frac{1}{2N^3}\underset{x}{}(F_{13}^2(x)+F_{23}^2(x))>$$
(3)
$$\mathrm{Space}\mathrm{like}\mathrm{Link}:L_s<\frac{1}{2N^3}\underset{x}{}[\mathrm{cos}(\chi (x+\widehat{1})+\theta _{\widehat{1}}(x)\chi (x))$$
(4)
$$+\mathrm{cos}(\chi (x+\widehat{2})+\theta _{\widehat{2}}(x)\chi (x))]>$$
(5)
$$\mathrm{Time}\mathrm{like}\mathrm{Link}:L_t<\frac{1}{N^3}\underset{x}{}\mathrm{cos}(\chi (x+\widehat{3})+\theta _{\widehat{3}}(x)\chi (x))>$$
(6)
$$\mathrm{Higgs}\mathrm{field}\mathrm{measure}\mathrm{squared}:\rho ^2\frac{1}{N^3}\underset{x}{}R^2(x)$$
(7)
We now proceed to the naïve continuum limit (at tree level) of the lattice action . The first step is to rewrite the action in terms of the continuum fields (denoted by a bar):
$$\phi =\overline{\phi }(\frac{2a_s}{\beta _{hs}})^{1/2},$$
$$A_1=a_s\overline{A}_1,A_2=a_s\overline{A}_2,A_3=a_t\overline{A}_3.$$
This means that the time-like part of the pure gauge action is rewritten in the form: $`\beta _{gt}a_ta_s^2a_t[(\overline{F}_{13})^2+(\overline{F}_{23})^2]\beta _{gt}a_td^3x[(\overline{F}_{13})^2+(\overline{F}_{23})^2]`$. On the other hand the space–like part is: $`\beta _{gs}\frac{a_s^2}{a_t}a_s^2a_t(\overline{F}_{12})^2\beta _{gs}\frac{a_s^2}{a_t}d^3x(\overline{F}_{12})^2`$. If we define
$$\beta _{gs}\frac{a_t}{g_s^2a_s^2},\beta _{gt}\frac{1}{g_t^2a_t},$$
(8)
the resulting continuum action reads:
$$\frac{1}{2}d^3x\left[\frac{1}{g_s^2}(\overline{F}_{12})^2+\frac{1}{g_t^2}[(\overline{F}_{13})^2+(\overline{F}_{23})^2]\right]$$
Defining $`\gamma _g(\frac{\beta _{gt}}{\beta _{gs}})^{1/2}`$ and using the definitions of $`\beta _{gs}`$, $`\beta _{gt}`$ we find that $`\gamma _g=\frac{g_s}{g_t}\frac{a_s}{a_t}`$. We denote by $`\xi `$ the important ratio $`\frac{a_s}{a_t}`$ of the two lattice spacings (the correlation anisotropy parameter) and finally derive the relation:
$$\gamma _g=\left(\frac{\beta _{gt}}{\beta _{gs}}\right)^{(1/2)}=\frac{g_s}{g_t}\xi .$$
Along the same lines, one may rewrite the scalar sector of the action in the form:
$$d^3x[|D_1\overline{\phi }|^2+|D_2\overline{\phi }|^2+\frac{\gamma _\phi ^2}{\xi ^2}|D_3\overline{\phi }|^2+m^2\overline{\phi }^{}\overline{\phi }+\lambda (\overline{\phi }^{}\overline{\phi })^2]$$
We have used the notations:
$$\gamma _\phi \left(\frac{\beta _{ht}}{\beta _{hs}}\right)^{1/2},$$
$$m^2a{}_{s}{}^{}{}_{}{}^{2}\frac{2}{\beta _{hs}}(12\beta _R2\beta _{hs}\beta _{ht}),\lambda a_s=\frac{4\beta _R}{\beta _{hs}^2\xi }.$$
Of course, the dependence of the space-like and time-like couplings on the lattice spacings $`a_s\mathrm{and}a_t`$ is a very interesting issue, which has been pursued in the context of QCD . Neglecting quantum corrections we have $`\gamma _\phi =\gamma _g=\xi .`$ In this paper we would like to explore the decoupling of the space-like planes, so in principle we should have a relation for the variation of $`\xi `$ as a function of $`\gamma _\phi `$ and $`\gamma _g`$, with the quantum corrections properly taken into account. We defer this study to a later stage and for the time being we adopt the tree level scaling refered to previously. In particular, we will vary the quantities $`\beta _{gs},\beta _{gt},\beta _{hs},\beta _{ht}`$ in such a way that we have
$$\gamma _\phi =\gamma _g\zeta .$$
We have used the letter $`\zeta `$ for the common ratio of the couplings, so our assumption reads: $`\xi =\zeta `$ and $`g_s=g_tg.`$
The parameter $`\beta _R`$ is found from the equation $`\beta _R=\frac{x\beta _{hs}^2}{4\beta _{gs}},`$ using the fixed value $`x=2`$ for the parameter $`x\frac{\lambda }{g^2}.`$
The compact action is defined as:
$$S_{compact}=\beta _{gs}\underset{x}{}(1\mathrm{cos}F_{12}(x))+\beta _{gt}\underset{x}{}\left[(1\mathrm{cos}F_{13}(x))+(1\mathrm{cos}F_{23}(x))\right]$$
$$+\beta _{hs}\underset{x}{}[2\phi ^{}(x)\phi (x)\phi ^{}(x)U_{\widehat{1}}(x)e^{iA_{\widehat{1},ext}(x)}\phi (x+\widehat{1})\phi ^{}(x)U_{\widehat{2}}(x)e^{iA_{\widehat{2},ext}(x)}\phi (x+\widehat{2})]$$
$$+\beta _{ht}\underset{x}{}[\phi ^{}(x)\phi (x)\phi ^{}(x)U_{\widehat{3}}(x)e^{iA_{\widehat{3},ext}(x)}\phi (x+\widehat{3})]$$
$$+\underset{x}{}[(12\beta _R2\beta _{hs}\beta _{ht})\phi ^{}(x)\phi (x)+\beta _R(\phi ^{}(x)\phi (x))^2]$$
(9)
The difference from the non-compact version lies in the gauge kinetic terms; however, since we will treat the compact system in an external magnetic field, we have taken the opportunity to also include in the action the background potential $`A_{\widehat{k},ext}(x),`$ which will generate the external field. There are many ways to impose such an external magnetic field ; for example, one might include it only in the interaction term with the matter fields; that is, the matter fields interact with both the quantum and the background fields. This has been our choice, as can be seen from equation (9).
The order parameters in this case are defined in a very similar way as in the non-compact case. We state the expressions which differ:
$$\mathrm{Space}\mathrm{like}\mathrm{Plaquette}:P_s<\frac{1}{N^3}\underset{x}{}cosF_{12}(x)>$$
(10)
$$\mathrm{Time}\mathrm{like}\mathrm{Plaquette}:P_t<\frac{1}{2N^3}\underset{x}{}[cosF_{13}(x)+cosF_{23}(x)]>$$
(11)
## 3 Algorithms
We used the Metropolis algorithm for the updating of both the gauge and the Higgs field. It is known that the scalar fields have much longer autocorrelation times than the gauge fields. Thus, special care must be taken to increase the efficiency of the updating for the Higgs field. We made the following additions to the Metropolis updating procedure :
a) Global radial update: We update the radial part of the Higgs field by multiplying it by the same factor at all sites: $`R(x)e^\xi R(x),`$ where $`\xi [\epsilon ,\epsilon ]`$ is randomly chosen. The quantity $`\epsilon `$ is adjusted such that the acceptance rate is kept between 0.6 and 0.7. The probability for the updating is $`P(\xi )=`$ min$`\{1,\mathrm{exp}(2V\xi \mathrm{\Delta }S(\xi ))\}`$ where $`\mathrm{\Delta }S(\xi )`$ is the change in action, while the $`2V\xi `$ term comes from the change in the measure.
b) Higgs field overrelaxation: We write the Higgs potential at $`x`$ in the form:
$$V(\phi (x))=𝐚𝐅+R^2(x)+\beta _R(R^2(x)1)^2$$
(12)
where
$$𝐚\left(\begin{array}{c}R(x)\mathrm{cos}\chi (x)\\ R(x)\mathrm{sin}\chi (x)\end{array}\right),$$
$$𝐅\left(\begin{array}{c}F_1\\ F_2\end{array}\right),$$
$$F_1\beta _{hs}[\underset{i=1,2}{}[R(x+\widehat{i})\mathrm{cos}(\chi (x+\widehat{i})+\theta _{\widehat{i}}(x))]+\beta _{ht}R(x+\widehat{3})\mathrm{cos}(\chi (x+\widehat{3})+\theta _{\widehat{3}}(x)),$$
$$F_2\beta _{hs}\left[\underset{i=1,2}{}R(x+\widehat{i})\mathrm{sin}(\chi (x+\widehat{i})+\theta _{\widehat{i}}(x))\right]+\beta _{ht}R(x+\widehat{3})\mathrm{sin}(\chi (x+\widehat{3})+\theta _{\widehat{3}}(x)).$$
We can perform the change of variables: $`(𝐚,𝐅)(X,F,𝐘)`$ ,where
$$F|𝐅|,𝐟\frac{𝐅}{\sqrt{F_1^2+F_2^2}},X𝐚𝐟,𝐘𝐚X𝐟.$$
(13)
The potential may be rewritten in terms of the new variables:
$$\overline{V}(X,F,𝐘)=XF+(1+2\beta _R(𝐘^21))X^2+𝐘^2(12\beta _R)+\beta _R(X^4+𝐘^4).$$
(14)
The updating of $`𝐘`$ is done simply by the reflection:
$$𝐘𝐘^{}=𝐘.$$
(15)
The updating of X is performed by solving the equation:
$$\overline{V}(X^{},F,𝐘)=\overline{V}(X,F,𝐘)$$
(16)
with respect to $`X^{}.`$ Noting that $`X^{}=X`$ is obviously a solution, we may factor out the quantity $`X^{}X`$ and reduce the quartic equation into a cubic one, which may be solved. The change $`XX^{}`$ is accepted with probability: $`P(X^{})=`$ min$`\{P_0,1\},`$ where $`P_0\frac{\overline{V}(X,F,𝐘)}{X}/\frac{\overline{V}(X^{},F^{},𝐘^{})}{X^{}}`$.
## 4 Results
### 4.1 Runs with fixed $`\zeta `$
In this set of measurements we set $`\beta _{gs}=6.0`$,$`x=2`$, and let $`\beta _{hs}`$ run. The remaining coupling constants vary according to the value of $`\zeta ,`$ that is:
$$\beta _{ht}=\zeta ^2\beta _{hs},\beta _{gt}=\zeta ^2\beta _{gs},\beta _R=\frac{x\beta _{hs}^2}{4\beta _{gs}}.$$
(17)
We recall that $`\zeta `$ is the ratio: $`\zeta =\sqrt{\frac{\beta _{gt}}{\beta _{gs}}}=\sqrt{\frac{\beta _{ht}}{\beta _{hs}}}`$. In figure 1 we show a comparison of the behaviour of the measure of the scalar field $`\rho ^2`$ for the symmetric $`(\zeta =1.0)`$ and a highly asymmetric model $`(\zeta =0.1)`$, that we call “S(ymmetric) model” and “A(symmetric) model” respectively. We have used the compact action in both models. The S model ($`\zeta =1.0`$) shows that the system undergoes a phase transition starting at $`\beta _{hs}0.35`$ from the 3D Coulomb phase to the 3D Higgs phase. (We note here that $`\rho ^2`$ is not going to infinity for big $`\beta _{hs}`$, which is the case in several treatments; the reason is that $`\beta _R`$ is not fixed, but increases with increasing $`\beta _{hs}`$). Then there is a transition region from $`\beta _{hs}0.35`$ to $`\beta _{hs}0.50`$; for $`\beta _{hs}0.50`$ the system is in the Higgs phase, characterized by a large value of $`\rho ^2`$.
The corresponding transition for the A model ($`\zeta =0.1`$) moves towards bigger values of $`\beta _{hs}`$ and it seems smoother. In the latter case the “time–like” coupling constants are $`\zeta ^2=0.01`$ times smaller than their “space–like” partners. Let us note that for $`\zeta =0`$ there will be no communication at all among the planes. This lends support to the assumption that for $`\zeta =0.1`$, that is $`\zeta ^2=0.01`$, we may have reached the limit of decoupled layers; indeed, we have checked that for an even smaller ratio of couplings ($`\zeta =0.01`$) the curve for the A model does not change much. This presumably means that the time–like separation $`a_t`$ of the spatial planes is already much bigger than the spatial lattice spacing $`a_s`$. This is further supported by the measurement of other quantities, which depend on the direction (see, for example, figure 2 and especially figure 3 below). It turns out that the quantities $`P_s`$, $`L_s`$ are much bigger than $`P_t`$, $`L_t`$ for all the range of $`\beta _{hs}`$, indicating that the quantities related to the communication of the planes are negligible in this case, as compared against the similar quantities within the layers. It seems safe to assume that the region $`\beta _{hs}0.55`$ corresponds to a Coulomb phase, where the layers are decoupled: this is equivalent to the Coulomb phase of the corresponding two–dimensional model. The transition region extends from $`\beta _{hs}0.55`$ to $`\beta _{hs}0.65`$. and then we have a Higgs phase for the A model: it has smaller $`\rho ^2`$ than the one of the S model and the quantities related to the third dimension are very small. This presumably means that the layers are decoupled also in this phase, so the picture is that we have moved effectively to the Higgs phase of the corresponding two–dimensional model. Thus we may say in brief that in the A model we see a transition from the 2D Coulomb phase to the 2D Higgs phase. Let us stress here that we have been using the words Higgs “phase” and “phase” transition in a rather loose meaning when refering to the asymmetric model. We have not checked that the passage from one region to the other is actually a phase transition: it may be a crossover, so the exact characterization is open for the time being. We should remark that the critical parameter for the Higgs phase transition is of order $`\frac{1}{d}`$, where d is the space–time dimension. Thus, it is not accidental that the isotropic model has a phase transition at $`\beta _{hs}0.35`$; this is close to the expected value, since d=3. On the other hand, the anisptropic model is effectively two–dimensional, so one should expect a value approximately equal to $`\frac{1}{2}`$, in semi–quantitative agreement with our results.
We note that the relative position of the transitions means that for $`\beta _{hs}0.35`$ the system lies in the Coulomb phase for either value of $`\zeta ;`$ for $`\beta _{hs}=0.50`$ the S model lies in its 3D Higgs phase, while the A model is still in its 2D Coulomb region; for $`\beta _{hs}=0.70`$ both models lie in their respective Higgs phases. The region between the two curves of figure 1 is presumably full of other curves corresponding to all values of $`\zeta `$ between 0.1 and 1. Later on we will consider the behaviour of several quantities with fixed $`\beta _{hs}`$ but varying $`\zeta `$; figure 1 suggests that we may use $`\beta _{hs}=0.5`$ and $`\beta _{hs}=0.7`$, based on the comments we made above.
Figure 2 shows the same models as above, but the space-like link is depicted rather than $`\rho ^2;`$ in addition, we have shown the results for both the compact and the non-compact version of the model. The changes take place at the same values of $`\beta _{hs}`$, as in figure 1. The results of the non-compact model are very close to the ones of the compact model, the most sizable differences being observed in model S, where the discrepancy reaches the value 0.2 in the region of the phase transition. In model A the differences are really small.
In figure 3 we display the time-like link and plaquette only for the compact version of model A $`(\zeta =0.1),`$ to get a more concrete idea of the layered phase. The parameters vary in the same way as above. We find that the time-like link starts increasing around $`\beta _{hs}=0.6,`$ which is in the region where the passage in model A is located in figures 1 and 2. The time-like plaquette $`P_t`$ does not seem to “feel” the passage and assumes a constant value. But the most important characteristic is that both quantities are really very small: their maximal value is $`0.03,`$ while their space-like partners $`L_s`$ and $`P_s`$ take on much larger values: $`0.13<L_s<1.0,0.90<P_s<1.0.`$ This is of course a consequence of our choice of bare parameters and confirms the fact that quantities related to communication between the planes are almost negligible in this region of the parameter space.
It is possible to show in a very explicit way the decoupling of the planes if we restrict ourselves to the Coulomb phase. The Wilson loops $`W(N,N)`$ can be calculated analytically in the 2D (pure) U(1) model and the resulting values can then be compared to the (space-like) plaquettes found by Monte Carlo simulations of the asymmetric Higgs model in the extremely layered Coulomb phase. In the table that follows we have put the results for the Wilson loops with sizes $`N\times N`$ from $`N=1`$ to $`N=6.`$ The values in parentheses are the statistical errors in the last one or two digits. The values of the parameters are reported in the table and the quantity $`\zeta `$ is set to 0.1. We also present in the table the values $`W(N,N)`$ from the analytical expressions : $`W(N,N)=(\frac{I_1(\beta _{gs})}{I_0(\beta _{gs})})^{(N\times N)}.`$ It is obvious that the system has moved to essentially two-dimensional dynamics, since most of the discrepancies lie within the error bars of the Monte Carlo data.
$$\begin{array}{ccc}& & \\ \multicolumn{3}{c}{\mathrm{Table}1}\\ & & \\ \multicolumn{3}{c}{\mathrm{Volume}=12^3,\zeta =0.1,\beta _{hs}=0.30,x=2.0}\\ & & \\ \mathrm{Loop}\mathrm{size}& \mathrm{Measured}& \mathrm{Analytic}(2\mathrm{D},\text{[15]})\\ & & \\ \multicolumn{3}{c}{\beta _{gs}=2.0}\\ & & \\ 1\times 1& 0.6981(1)& 0.697775\\ 2\times 2& 0.2380(2)& 0.237061\\ 3\times 3& 0.0395(2)& 0.0392136\\ 4\times 4& 0.0033(2)& 0.00315823\\ 5\times 5& 0.0002(2)& 0.000123845\\ 6\times 6& 0.0001(1)& 0.0000023645\\ & & \\ \multicolumn{3}{c}{\beta _{gs}=4.0}\\ & & \\ 1\times 1& 0.8635(1)& 0.863523\\ 2\times 2& 0.5562(2)& 0.556026\\ 3\times 3& 0.2669(4)& 0.266971\\ 4\times 4& 0.0955(5)& 0.0955827\\ 5\times 5& 0.0251(4)& 0.0255178\\ 6\times 6& 0.0053(2)& 0.00507988\\ & & \\ \multicolumn{3}{c}{\beta _{gs}=6.0}\\ & & \\ 1\times 1& 0.9123(1)& 0.912359\\ 2\times 2& 0.6931(2)& 0.692889\\ 3\times 3& 0.4384(4)& 0.438019\\ 4\times 4& 0.2307(9)& 0.230491\\ 5\times 5& 0.1008(13)& 0.10096\\ 6\times 6& 0.0360(11)& 0.0368106\\ & & \\ \multicolumn{3}{c}{\beta _{gs}=8.0}\\ & & \\ 1\times 1& 0.9352(1)& 0.935235\\ 2\times 2& 0.7655(1)& 0.76504\\ 3\times 3& 0.5490(3)& 0.54738\\ 4\times 4& 0.3454(6)& 0.342559\\ 5\times 5& 0.1909(9)& 0.18751\\ 6\times 6& 0.0925(10)& 0.089775\end{array}$$
### 4.2 Runs with varying $`\zeta `$
Now we come to a somewhat more detailed examination of the phase diagram where the “space-like” coupling constants are kept fixed and the ratio $`\zeta `$ changes, also forcing $`\beta _{gt}`$ and $`\beta _{ht}`$ to change. Of course, the result will strongly depend on the region in which the fixed “space–like” coupling constants lie; we will present results for the two interesting choices that we described in the discussion of figure 1: $`\beta _{gs}=6.0`$ and the values 0.5 and 0.7 for $`\beta _{hs}.`$
In figure 4 we show the space-like link $`L_s`$ for the case with $`\beta _{hs}=0.5;`$ again both compact and non-compact results are displayed. The parameter $`\zeta `$ starts from zero, where we expect a Coulomb phase with fully separated layers, that is a 2D model, to 0.8, which corresponds to a model close to the symmetric 3D one. The differences between the compact and the non–compact version are localized in the transition region. The non–compact results are somewhat bigger, the maximum difference being about 0.05, but qualitatively the results are similar for both versions. The transition from the symmetric to the broken phase takes place at about $`\zeta 0.5,`$ where 2D Coulomb becomes 3D Higgs: The behaviour of the space-like link may be interpreted with the help of figure 2. As we change $`\zeta `$, we move continuously to new curves, of the sort of the ones in figure 2, but characterizing the $`\zeta `$ under study. When $`\zeta =0,`$ the corresponding curve (close to the lowest one in figure 2) will inform us that the value $`\beta _{hs}=0.5`$ corresponds to the Coulomb phase for this value of $`\zeta .`$ The curves for larger $`\zeta `$ will differ from the initial one in an important aspect: the transition region will move to the left (to smaller $`\beta _{hs}`$) and will be steeper, until finally they will coincide with the uppermost curve in figure 2. When the transition region reaches $`\beta _{hs}=0.5,`$ the value of $`L_s`$ will increase abruptly and the system will pass over to a Higgs region. It is not clear at this point whether this is 2D or 3D Higgs, but this will be clarified by figure 5, which follows.
The time-like link $`L_t,`$ shown in figure 5 both for the compact and non-compact versions, seems a more sensitive detector of the phase transition: it changes already in the region $`0<\zeta <0.5,`$ where the space-like link is essentially constant, and its variation is bigger, since it starts from zero. Of course, it has to become equal to the space-like link (about one) in the limit $`\zeta 1,`$ since there one regains the symmetric model. The non-compact results lie systematically above the compact ones. The comparison of this figure with the previous one shows that the time–like link also gets its 3D value at the same time when the space–like link jumps. This fact indicates that we have a transition from the 2D Coulomb phase to a 3D Higgs, that is, both changes (in dimension and in the breaking of symmetry) appear to happen simultaneously.
In figure 6 we show the space-like link $`L_s`$ versus $`\zeta `$ for the second of the two chosen values, $`\beta _{hs}=0.7.`$ In this case $`\beta _{hs}`$ lies to the right of the transition already for the curve corresponding to $`\zeta =0,`$ as may be seen in figure 2. Thus we do not expect a phenomenon as in figure 4, where the transition region was passing through $`\beta _{hs}.`$ We start from a layered Higgs phase already at $`\zeta =0`$ and we may isolate the transition from the 2D to the 3D Higgs phase (the space–like link has big values, from 0.90 to 0.97, in all the range of $`\zeta ,`$ in striking contrast with figure 4). The transition happens at $`\beta _{hs}0.25`$ for the compact case. It seems smooth, however one may observe a volume dependence in the results of figure 6, where we have put mesurements from both $`8^3`$ and $`16^3`$ lattices. In addition, we remark that the points for the compact $`8^3`$ have a non–monotonous region between 0.15 and 0.25; this disappears in the $`16^3`$ lattice, so we interpret it as finite size artifact. The same effect is seen in the non–compact $`8^3`$ case, but seems unimportant, in view of the experience with the compact case. The behaviour of the non–compact $`8^3`$ is smoother than its compact partner.
Figure 7 contains the results for the time-like link $`L_t`$ versus $`\zeta `$ for the two volumes $`8^3`$ and $`16^3.`$ This quantity is a good indicator of the 2D to 3D phase transition, as we said above; its volume dependence is almost negligible, but the difference between the compact and non–compact version is very pronounced at the region of the transition. The non–compact version appears steeper, in contrast with the previous figure 6.
### 4.3 Vortex dynamics in the compact version
We will now consider the system under the influence of a homogeneous external magnetic field; thus one should construct first a lattice version of the homogeneous magnetic field. This has already been done before in in connection with the abelian Higgs model, or (2+1) QED . We follow a slightly different prescription, which we describe below.
Since we would like to impose an external homogeneous magnetic field in the $`z`$ direction, we choose the external gauge potential in such a way that the plaquettes in the $`xy`$ plane equal B, while all other plaquettes equal zero. One way in which this can be achieved is through the choice: $`A_z(x,y,z)=0`$, for all integers $`x,y,z,`$ running from 1 to N and
$$A_x(x,y,z)=\frac{B}{2}(y1),xN,A_x(N,y,z)=\frac{B}{2}(N+1)(y1),$$
(18)
$$A_y(x,y,z)=+\frac{B}{2}(x1),yN,A_y(x,N,z)=+\frac{B}{2}(N+1)(x1).$$
(19)
where $`N^3`$ is the number of points on the (cubic) lattice and the coordinates $`x,y,z`$ are integers running from 1 to N. It is trivial to check out that all plaquettes starting at $`(x,y,z),`$ with the exception of the one starting at $`(N,N,z),`$ equal $`B.`$ The latter plaquette equals $`(1N^2)B=B(N^2B).`$ One may say that the flux is homogeneous over the entire $`xy`$ cross section of the lattice and equals $`B.`$ The additional flux of $`(N^2B)`$ can be understood by the fact that the lattice is a torus, that is a closed surface, and the Maxwell equation $`𝐁=0`$ implies that the magnetic flux through the lattice should vanish. This means that, if periodic boundary conditions are used for the gauge field, the total flux of any configuration should be zero, so the (positive, say) flux $`B,`$ penetrating the majority of the plaquettes, will be accompanied by a compensating negative flux $`(N^2B)`$ in a single plaquette. This compensating flux should be “invisible”, that is it should have no observable physical effects. This is the case if the flux is an integer multiple of $`2\pi :N^2B=n2\pi B=n\frac{2\pi }{N^2},`$ where $`n`$ is an integer. Thus we may say (disregarding the “invisible” flux) that the magnetic field is homogeneous over the entire cross section of the lattice. The integer $`n`$ may be chosen to lie in the interval $`[0,\frac{N^2}{2}],`$ with the understanding that the model with integers $`m`$ between $`\frac{N^2}{2}`$ and $`N^2`$ is equivalent to the model with integers taking on the values $`N^2n,`$ which are among the ones that have already been considered. It follows that the magnetic field strength B in lattice units lies between 0 and $`\pi .`$ The physical magnetic field $`B_{phys}`$ is related to $`B`$ through $`B=g_sa_s^2B_{phys},`$ and the physical field may go to infinity letting the lattice spacing $`a_s`$ go to zero, while $`B`$ is kept constant.
One of the first questions which may be asked in connection with the existence of the external filed is of course its penetration in the bulk of the lattice. It is this question which is treated in table 2. The first column in the table is the integer determining the magnetic field; according to the previous discussion, its range is from 0 to 32 for the $`8^3`$ lattice that we have used (in table 2 we show the results up to $`n=20`$ only). Apart from $`P_s`$ and $`P_t`$ we show the expectation value $`P_s(B)`$ of the space–like plaquette which includes the external field. Its definition reads: $`P_s(B)<\frac{1}{N^3}_x\mathrm{cos}(F_{12}(x)+B)>`$. The last column contains the number $`Q`$ of vortices that penetrate into the lattice to shield the external magnetic field (Meissner effect). For this topological number $`Q`$ we use the definition of . We consider the space–like plaquette $`P_s(x)`$ and define the quantity $`\widehat{P}_s(x),`$ and the topological number $`q_s(x)`$ for the space–like plaquette at position $`x`$ through the equality:
$$P_s(x)=\widehat{P}_s(x)+2\pi q_s(x)$$
and the demand that $`\widehat{P}_s`$ lies in the interval $`(\pi ,+\pi ].`$ The topological number $`Q`$ is the sum of the quantities $`q_s(x)`$ over space–like plaquettes with fixed z-coordinates:
$$Q\underset{xy}{}q_s(x).$$
We show the results for a symmetric $`(\zeta =1.0)`$ $`8^3`$ lattice, with three different values for $`\beta _{hs}.`$
For $`\beta _{hs}=0.200`$ (which corresponds to the Coulomb phase) we observe that neither the space-like nor the time-like plaquettes depend on the magnetic field. (The two kinds of plaquettes could in principle be different even for a symmetric lattice, because of the presence of the external field). The same is true for $`<\rho ^2>,`$ which has a constant value corresponding to the Coulomb phase. The quantity $`P_s(B)`$ varies with the magnetic field, but its variation is trivial, in the sense that its value is the product of $`P_s`$ and $`\mathrm{cos}B.`$ We recall that:
$$P_s(B)=(\mathrm{cos}B)<\mathrm{cos}F_{12}>(\mathrm{sin}B)<\mathrm{sin}F_{12}>.$$
This suggests that the magnetic field penetrates completely and there is no sign of shielding. The behaviour of the topological number $`Q`$ is consistent with this picture: it is zero for all the values of the magnetic field, indicating that no vortices are formed by the fluctuating gauge field and the background field penetrates with no obstacles into the bulk of the lattice.
For $`\beta _{hs}=0.500`$ the system is in the Higgs phase, as may be seen from the value of $`<\rho ^2>.`$ For relatively small values of the external field (up to $`n=2`$) $`P_t`$ and $`P_s(B),`$ as well as $`<\rho ^2>,`$ are approximately constant, while $`P_s`$ changes even for small fields. It seems that the system tunes $`F_{12}`$ in such a way as to compensate the change of B, so that $`P_s(B)`$ is constant. Small magnetic fields appear to get shielded by the dynamics of the scalar and gauge bosons. The picture changes for bigger external fields: $`P_s`$ changes and moves towards a value characterizing the Coulomb phase and this phenomenon is also confirmed from $`<\rho ^2>.`$ The picture is that too big magnetic fields cannot be shielded any more and they penetrate into the lattice; then we pass to the Coulomb phase without Meissner effect. When we pass to the Coulomb phase $`P_s(B)`$ satisfies $`P_s(B)=(\mathrm{cos}B)P_s`$ approximately. One may check $`n=6.`$ This picture is further supported by the variation of $`Q.`$ This number is 1 for $`n=1,`$, which signals the creation of a vortex to exactly cancel this magnetic field, while for $`n=2`$ it equals 2, shielding the external field completely. For $`n=3`$ and $`n=4`$ no new vortices are created, so the screening is incomplete, and finally, for even bigger $`n,`$ even these two vortices disappear.
$$\begin{array}{ccccccc}& & & & & & \\ \multicolumn{7}{c}{\mathrm{Table}2}\\ & & & & & & \\ \multicolumn{7}{c}{\mathrm{Volume}=8^3,\zeta =1.0,\beta _{gs}=6.0,x=2.0}\\ & & & & & & \\ \mathrm{n}& \mathrm{cosB}& \mathrm{P}_\mathrm{s}& \mathrm{P}_\mathrm{s}(\mathrm{B})& \mathrm{P}_\mathrm{t}& \rho ^2& Q\\ & & & & & & \\ \multicolumn{7}{c}{\beta _{hs}=0.200}\\ & & & & & & \\ 0& 1.0& 0.9429& 0.9248& 0.9429& 1.03& 0\\ 4& 0.9238& 0.9428& 0.8711& 0.9429& 1.03& 0\\ 6& 0.8314& 0.9427& 0.7838& 0.9429& 1.04& 0\\ 8& 0.7071& 0.9429& 0.6667& 0.9429& 1.03& 0\\ 12& 0.3826& 0.9427& 0.3608& 0.9429& 1.04& 0\\ 16& 0.00& 0.9427& 0.0000& 0.9429& 1.04& 0\\ 20& 0.382& 0.9428& 0.3608& 0.9429& 1.03& 0\\ & & & & & & \\ \multicolumn{7}{c}{\beta _{hs}=0.500}\\ & & & & & & \\ 0& 1.0& 0.9526& 0.9526& 0.9527& 12.08& 0\\ 1& 0.9951& 0.9479& 0.9525& 0.9527& 12.09& 1\\ 2& 0.9807& 0.9338& 0.9520& 0.9526& 12.09& 2\\ 3& 0.9569& 0.9226& 0.9367& 0.9518& 11.40& 2\\ 4& 0.9238& 0.9250& 0.9003& 0.9506& 10.31& 2\\ 5& 0.8819& 0.9325& 0.8235& 0.9493& 9.05& 0\\ 6& 0.8314& 0.9325& 0.7765& 0.9488& 8.54& 0\\ 8& 0.7071& 0.9332& 0.6612& 0.9480& 7.72& 0\\ 12& 0.3826& 0.9346& 0.3585& 0.9464& 6.01& 0\\ 16& 0.00& 0.9376& 0.0004& 0.9455& 4.82& 0\\ 20& 0.3826& 0.9394& 0.3594& 0.9448& 3.90& 0\\ & & & & & & \\ \multicolumn{7}{c}{\beta _{hs}=0.700}\\ & & & & & & \\ 0& 1.0& 0.9570& 0.9570& 0.9571& 13.92& 0\\ 2& 0.9807& 0.9380& 0.9563& 0.9569& 13.92& 2\\ 4& 0.9238& 0.8825& 0.9550& 0.9568& 13.90& 4\\ 6& 0.8314& 0.7920& 0.9522& 0.9562& 13.89& 6\\ 7& 0.7730& 0.7349& 0.9502& 0.9559& 13.90& 7\\ 8& 0.7071& 0.6706& 0.9475& 0.9555& 13.90& 8\\ 9& 0.6343& 0.6904& 0.8905& 0.9546& 13.19& 7\\ 10& 0.5555& 0.8011& 0.7411& 0.9535& 11.93& 4\\ 12& 0.3826& 0.8149& 0.5614& 0.9525& 11.17& 3\\ 14& 0.1950& 0.8499& 0.2684& 0.9517& 10.36& 1\\ 16& 0.000& 0.8561& 0.0167& 0.9512& 9.89& 0\\ 18& 0.1950& 0.8593& 0.1531& 0.9509& 9.59& 0\\ 20& 0.3826& 0.8638& 0.3175& 0.9507& 9.41& 0\end{array}$$
For $`\beta _{hs}=0.700`$ the system is even deeper than before in the Higgs phase, as is obvious from the last column. This enables the system to screen the external field more efficiently: $`P_s(B)`$ is almost constant up to $`n=8,`$ as compared to $`n=2`$ previously. The variation of $`P_s`$ starts early again. The penetration of B and the transition to the Coulomb phase is slower than before. One may check easily that $`P_s(B)`$ is no longer the product of $`P_s`$ and $`\mathrm{cos}B`$, at least for small $`B`$; for large $`B`$, which is the region where the system approaches the Coulomb phase, the equation $`P_s(B)=(\mathrm{cos}B)P_s`$ is again approximately satisfied (one may check, for example, $`n=16`$). The behaviour of the topological number is interesting here: the system creates $`Q`$ vortices, with $`Q=n,`$ for $`n8.`$ As the magnetic field increases, the vortex number $`Q`$ decreases and the shielding is incomplete; for $`n16`$ no vortex is present any more and the system moves towards complete penetration.
Finally, we note that $`P_t`$ is not very sensitive to the presence of the magnetic field anyway.
$$\begin{array}{ccccccccccccccccc}& & & & & & & & & & & & & & & & \\ \multicolumn{17}{c}{Table3}\\ & & & & & & & & & & & & & & & & \\ \multicolumn{17}{c}{q_s(x)\mathrm{at}\mathrm{random}\mathrm{z}\mathrm{coordinates}\mathrm{for}n=2}\\ & & & & & & & & & & & & & & & & \\ \multicolumn{8}{c}{\zeta =1.0}& & \multicolumn{8}{c}{\zeta =0.1}\\ & & & & & & & & & & & & & & & & \\ 0& 0& 0& 0& 0& 0& 0& 0& & 0& 1& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& & 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0& & 0& 0& 0& 1& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& & 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& & 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0& & 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& & 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& & 0& 0& 0& 1& 0& 0& 0& 0\end{array}$$
To follow more closely the vortex formation, we have concentrated on $`n=2`$ and studied the topological numbers $`q_s(x)`$ characterizing each plaquette in the $`xy`$ planes for both the S and the A models for $`\beta _{gs}=6.0,\beta _{hs}=0.700.`$ The parameters are chosen such that the system is deeply in the Higgs phase; thus, we expect that the system will form its own vortices to shield the external magnetic field. These vortices will reflect the winding number of the background field, namely $`n=2`$. In table 3 we show these numbers for two $`xy`$ cross sections of the lattice (one for $`\zeta =1.0`$ and one for $`\zeta =0.1`$); the z–coordinates of the planes are chosen at random. It is obvious that the system adjusts itself, so that it screens the (small) magnetic field. Significantly more vortices appear in the A model; this is presumably due to energetics, since in the A model each layer acts on its own. The most important difference is that the (two) vortices appearing in the S model are at the same positions for all the $`xy`$ planes; on the contrary, the planes in the A model appear decoupled and the vortices on each one of them have no relation to the ones of its nearest neighbours. In addition, we point out that in the A model the total sum of the topological numbers over the cross section need not necessarily yield two, since there are big fluctuations; for instance, in the case we show in table 3, it is $`+1`$. We have also used another definition for the winding number, in which one first brings the gauge invariant link variables in the interval $`(\pi ,+\pi ]`$ and afterwards follows the same steps as above; this definition has the advantage of being additive, so the sum over a closed surface, such as the cross section of a lattice should be zero. The results (not shown) have been that in the S model the two vortices shown before are still present, but two more vortices with $`Q=1`$ appear, so that the sum vanishes; this is also the case for the A model.
We illustrate in figure 8 the evolution of the winding number $`Q`$ in the S model ($`\zeta =1.0`$). We see that for the symmetric model the topological number changes abruptly and quickly from 0 to $`n=2`$.
The picture is totally different in the A model. In figures 9 and 10 we show the evolution of Q for two neighbouring layers. We observe that the topological number changes very slowly and the final value $`n=2`$ is reached only after passing slowly through intermediate values with significant fluctuations. Moreover, it is by no means sure that the system has settled in its final topological number and that no new fluctuations are going to occur. This behaviour is presumably due to the decoupling of the $`xy`$ layers. We have found that the decoupling takes place at $`\zeta 0.6`$, as one decreases $`\zeta `$ from one to zero.
In figures 11 and 12 we have shown the evolution of $`<\rho ^2>`$ for the S and the A model respectively. They reflect the features of the corresponding changes in Q shown in the figures 8, 9 and 10: a sharp and quick increase for the S model and a gradual one for the A model, since there the layers decouple and do not move together to the new state characterized by a non-zero topological number. This may be an indication for the generation of the so–called pancake vortices in the layered phase.
Acknowledgements
The authors would like to thank the TMR project “Finite temperature phase transitions in Particle Physics”, EU contact number: FMRX-CT97-0122 for financial support. Stimulating discussions with F.Karsch, A. Kehagias, C.P. Korthals–Altes, T. Neuhaus, S. Nicolis and N. Tetradis are gratefully acknowledged.
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# Spin tunneling of trigonal and hexagonal ferromagnets in an arbitrarily directed magnetic field
## I. Introduction
Recently, there has been great experimental and theoretical effort to observe and interpret macroscopic quantum tunneling (MQT) and coherence (MQC) in nanometer-scale magnets at sufficiently low temperature. Theoretical investigations based on the spin-coherent-state path integral were performed for the single-domain ferromagnetic (FM) nanoparticles, which showed that MQT and MQC were possible in magnets containing as much as $`10^510^6`$ spins. Several experiments involving resonance measurements, magnetic relaxation, and hysteresis loop study for various systems showed either temperature-independent relaxation phenomena or a well-defined resonance depending exponentially on the number of total spins, which supported the idea of magnetic quantum tunneling.
More recently, the tunneling behaviors of the magnetization vector were studied extensively for the single-domain FM nanoparticles in the presence of an external magnetic field applied at an arbitrary angle. The MQT problem for FM particles with uniaxial crystal symmetry was first studied by Zaslavskii who calculated the tunneling exponent, the preexponential factors and their temperature dependences in the low barrier limit with the help of mapping the spin system onto a one-dimensional particle system. For the same crystal symmetry, Miguel and Chudnovsky calculated the tunneling rate by applying the imaginary-time path integral, and demonstrated that the angular and field dependences of the tunneling exponent obtained by Zaslavskii’s method and by the path-integral method coincide precisely. They also discussed the tunneling rate at finite temperature and suggested experimental procedures. Kim and Hwang performed a calculation based on the instanton technique for FM particles with biaxial and tetragonal crystal symmetry, and Kim extended the tunneling rate for biaxial crystal symmetry to a finite temperature. The quantum-classical transition of the escape rate for FM particles with uniaxial crystal symmetry in an arbitrarily directed field was investigated by Garanin, Hidalgo and Chudnovsky with the help of mapping onto a particle moving in a double-well potential. The switching field measurement was carried out on single-domain FM nanoparticles of Barium ferrite (BaFeCoTiO) containing about $`10^510^6`$ spins. The measured angular dependance of the crossover temperature was found to be in excellent agreement with the theoretical prediction, which strongly suggests the MQT of magnetization in the BaFeCoTiO nanoparticles. Lü et al. studied the MQT and MQC of the Néel vector in single-domain antiferromagnetic (AFM) nanoparticles with biaxial, tetragonal, and hexagonal crystal symmetry in an arbitrarily directed field.
In this paper, we extend the previous theoretical results obtained for the single-domain FM particles with biaxial and tetragonal symmetry to those for FM particles with a much more complex structure placed in an external magnetic field at an arbitrarily directed angle in the $`ZX`$ plane, based on the instanton technique in the spin-coherent-state path-integral representation. We consider the magnetocrystalline anisotropies with trigonal and hexagonal crystal symmetry, respectively. Both the Wentzel-Kramers-Brillouin (WKB) exponents and the preexponential factors are evaluated analytically in the tunneling rates for MQT and the tunnel splittings for MQC in FM particles for different angle ranges of the external magnetic field $`(\theta _H=\pi /2`$, $`\pi /2+O\left(ϵ^{3/2}\right)<\theta _H<\pi O(ϵ^{3/2})`$, and $`\theta _H=\pi ),`$ and the temperature which corresponds to the crossover from the thermal to the quantum regime is clearly shown for each case. Our results show that the distinct angular dependence, together with the dependence of the WKB tunneling rate and the crossover temperature on the strength of the external magnetic field, may provide an independent experimental test for the magnetic tunneling in single-domain FM nanoparticles. The calculations performed in this paper are semiclassical in nature, i.e., valid for large spins and in the continuum limit. We analyze the validity of the semiclassical approximation, and find that the semiclassical approximation is rather good for the typical values of parameters for single-domain FM nanoparticles.
This paper is structured in the following way. In Sec. II, we briefly review the basic ideas of the MQT and MQC in single-domain FM particles. In Secs. III and IV, we study the quantum tunneling of the magnetization vector for FM particles with trigonal and hexagonal crystal symmetry in the presence of an external magnetic field applied in the $`ZX`$ plane with a range of angles $`\pi /2\theta _H\pi `$. The conclusions are presented in Sec. V. In Appendix A, we explain briefly the computation of the preexponential factors in the WKB tunneling rate, and then apply this approach to obtain the tunnel splittings for FM particles with trigonal crystal symmetry in a magnetic field applied perpendicular to the anisotropy axis $`(\theta _H=\pi /2)`$ in detail.
## II. MQT and MQC of the magnetization vector in FM particles
In this section we briefly review some basic ideas of MQT and MQC of the magnetization vector in single-domain FM nanoparticles, based on the instanton technique in the spin-coherent-state path integral.
The system of interest is a nanometer-scale single-domain ferromagnet at a temperature well below its anisotropy gap. For such a FM particle, the tunnel splitting for MQC or the tunneling rate for MQT is determined by the imaginary-time transition amplitude from an initial state $`|i`$ to a final state $`|f`$ as
$$U_{fi}=f\left|e^{HT}\right|i=D\mathrm{\Omega }\mathrm{exp}\left(S_E\right),$$
(1)
where $`S_E`$ is the Euclidean action and $`D\mathrm{\Omega }`$ is the measurement of the path integral. In the spin-coherent-state path-integral representation, the Euclidean action can be expressed as
$$S_E(\theta ,\varphi )=\frac{V}{\mathrm{}}𝑑\tau \left[i\frac{M_0}{\gamma }\left(\frac{d\varphi }{d\tau }\right)i\frac{M_0}{\gamma }\left(\frac{d\varphi }{d\tau }\right)\mathrm{cos}\theta +E(\theta ,\varphi )\right],$$
(2)
where $`V`$ is the volume of the FM particle and $`\gamma `$ is the gyromagnetic ratio. $`M_0=\left|\stackrel{}{M}\right|=\mathrm{}\gamma S/V`$, where $`S`$ is the total spin of FM particles. It is noted that the first two terms in Eq. (2) define the topological Berry or Wess-Zumino, Chern-Simons term which arises from the nonorthogonality of spin coherent states. The Wess-Zumino term has a simple topological interpretation. For a closed path, this term equals $`iS`$ times the area swept out on the unit sphere between the path and the north pole. The first term in Eq. (2) is a total imaginary-time derivative, which has no effect on the classical equations of motion, but it is crucial for the spin-parity effects. However, for the closed instanton or bounce trajectory described in this paper (as shown in the following), this time derivative gives a zero contribution to the path integral, and therefore can be omitted.
In discussing macroscopic quantum phenomena, it is essential to distinguish between two types of processes: MQC (i.e., coherent tunneling) and MQT (i.e., incoherent tunneling). In the case of MQC, the system in question performs coherent NH<sub>3</sub>-type oscillations between two degenerate wells separated by a classically impenetrable barrier. Tunneling between neighboring degenerate vacua can be described by the instanton configuration with nonzero topological charge and leads to a level splitting of the ground states. The tunneling removes the degeneracy of the original ground states, and the true ground state is a superposition of the previous degenerate ground states. For the case of MQT, the system escapes from a metastable potential well into a continuum by quantum tunneling at sufficiently low temperatures, and the tunneling results in an imaginary part of the energy which is dominated by the so-called bounce configuration with zero topological charge. As emphasized by Leggett, the two phenomena of MQC and MQT are physically very different, particularly from the viewpoint of experimental feasibility. MQC is a far more delicate phenomenon than MQT, as it is much more easily destroyed by an environment, and by very small $`c`$-number symmetry breaking fields that spoil the degeneracy.
In the semiclassical limit, the dominant contribution to the transition amplitude comes from the finite action solution (instanton) of the classical equation of motion. The motion of the magnetization vector $`\stackrel{}{M}`$ is determined by the Landau-Lifshitz equation,
$$i\frac{d\stackrel{}{M}}{d\tau }=\gamma \stackrel{}{M}\times \frac{dE\left(\stackrel{}{M}\right)}{d\stackrel{}{M}},$$
(3)
which can also be expressed as the following equations in the spherical coordinate system,
$`i\left({\displaystyle \frac{d\overline{\theta }}{d\tau }}\right)\mathrm{sin}\overline{\theta }`$ $`=`$ $`{\displaystyle \frac{\gamma }{M_0}}{\displaystyle \frac{E}{\overline{\varphi }}},`$ ()
$`i\left({\displaystyle \frac{d\overline{\varphi }}{d\tau }}\right)\mathrm{sin}\overline{\theta }`$ $`=`$ $`{\displaystyle \frac{\gamma }{M_0}}{\displaystyle \frac{E}{\overline{\theta }}},`$ ()
where $`\overline{\theta }`$ and $`\overline{\varphi }`$ denote the classical path. Note that the Euclidean action Eq. (2) describes the $`\left(11\right)`$-dimensional dynamics in the Hamiltonian formulation with canonical variables $`\varphi `$ and $`P_\varphi =S(1`$cos$`\theta )`$. The instanton’s contribution to the tunneling rate $`\mathrm{\Gamma }`$ for MQT or the tunnel splitting $`\mathrm{\Delta }`$ for MQC (not including the topological Wess-Zumino or Berry phase) is given by
$$\mathrm{\Gamma }(\text{or }\mathrm{\Delta })=A\omega _p\left(\frac{S_{cl}}{2\pi }\right)^{1/2}e^{S_{cl}},$$
(4)
where $`\omega _p`$ is the frequency of small oscillations near the bottom of the inverted potential, and $`S_{cl}`$ is the classical action. The preexponential factor $`A`$ originates from the quantum fluctuations about the classical path, which can be evaluated by expanding the Euclidean action to second order in the small fluctuations. In Ref. 12, Garg and Kim studied the general formalism for calculating both the exponent and the preexponential factors in the WKB tunneling rates for MQT and MQC in single-domain FM nanoparticles. In Appendix A, we explain briefly the basic idea of this calculation, and then apply this approach to calculate the instanton’s contribution to the tunnel splittings for MQC of the magnetization vector in FM particles with trigonal crystal symmetry in an external magnetic field perpendicular to the anisotropy axis (considered in Sec. III) in detail.
## III. MQC and MQT for trigonal crystal symmetry
In this section, we study the tunneling behaviors of the magnetization vector in single-domain FM nanoparticle with trigonal crystal symmetry. The external magnetic field is applied in the $`ZX`$ plane, at an angle in the range of $`\pi /2\theta _H<\pi `$. Now the total energy $`E(\theta ,\varphi )`$ can be written as
$$E(\theta ,\varphi )=K_1\mathrm{sin}^2\theta K_2\mathrm{sin}^3\theta \mathrm{cos}\left(3\varphi \right)M_0H_x\mathrm{sin}\theta \mathrm{cos}\varphi M_0H_z\mathrm{cos}\theta +E_0,$$
(5)
where $`K_1`$ and $`K_2`$ are the magnetic anisotropy constants satisfying $`K_1K_2>0`$, and $`E_0`$ is a constant which makes $`E(\theta ,\varphi )`$ zero at the initial orientation. As the magnetic field is applied in the $`ZX`$ plane, $`H_x=H\mathrm{sin}\theta _H`$ and $`H_z=H\mathrm{cos}\theta _H`$, where $`H`$ is the magnitude of the field and $`\theta _H`$ is the angle between the magnetic field and the $`\widehat{z}`$ axis.
In the absence of the external magnetic field, the system reduces to one with threefold rotational symmetry around $`\widehat{z}`$ axis and reflection symmetry in the $`XY`$ plane. The unit vectors $`\widehat{z}`$ and $`\widehat{z}`$ define the two classical ground state configurations. The transition amplitude between degenerate ground states can be suppressed to zero resulting from the destructive Wess-Zumino phase if the system has time-reversal invariance at zero magnetic field. However, for the closed instanton or bounce trajectory described in this paper (as shown in the following) the phase term in Eq. (2), proportional to $`d\varphi /d\tau `$ (not $`\left(d\varphi /d\tau \right)\mathrm{cos}\theta `$ term) gives a zero contribution to the integral Eq. (2) and, therefore, can be omitted.
By introducing the dimensionless parameters as
$$\overline{K}_2=K_2/2K_1,\overline{H}_x=H_x/H_0,\overline{H}_z=H_z/H_0,$$
(6)
Eq. (6) can be rewritten as
$$\overline{E}(\theta ,\varphi )=\frac{1}{2}\mathrm{sin}^2\theta \overline{K}_2\mathrm{sin}^3\theta \mathrm{cos}\left(3\varphi \right)\overline{H}_x\mathrm{sin}\theta \mathrm{cos}\varphi \overline{H}_z\mathrm{cos}\theta +\overline{E}_0,$$
(7)
where $`E(\theta ,\varphi )=2K_1\overline{E}(\theta ,\varphi )`$, and $`H_0=2K_1/M_0`$. At finite magnetic field, the plane given by $`\varphi =0`$ is the easy plane, on which $`\overline{E}(\theta ,\varphi )`$ reduces to
$$\overline{E}\left(\theta ,\varphi =0\right)=\frac{1}{2}\mathrm{sin}^2\theta \overline{K}_2\mathrm{sin}^3\theta \overline{H}\mathrm{cos}\left(\theta \theta _H\right)+\overline{E}_0.$$
(8)
We denote $`\theta _0`$ to be the initial angle and $`\theta _c`$ the critical angle at which the energy barrier vanishes when the external magnetic field is close to the critical value $`\overline{H}_c\left(\theta _H\right)`$ (to be calculated in the following). Then, the initial angle $`\theta _0`$ satisfies $`\left[d\overline{E}\left(\theta ,\varphi =0\right)/d\theta \right]_{\theta =\theta _0}=0`$, the critical angle $`\theta _c`$ and the dimensionless critical field $`\overline{H}_c`$ satisfy both $`\left[d\overline{E}\left(\theta ,\varphi =0\right)/d\theta \right]_{\theta =\theta _c,\overline{H}=\overline{H}_c}=0`$ and $`\left[d^2\overline{E}\left(\theta ,\varphi =0\right)/d\theta ^2\right]_{\theta =\theta _c,\overline{H}=\overline{H}_c}=0`$, which leads to
$`{\displaystyle \frac{1}{2}}\mathrm{sin}\left(2\theta _0\right)3\overline{K}_2\mathrm{sin}^2\theta _0\mathrm{cos}\theta _0+\overline{H}\mathrm{sin}\left(\theta _0\theta _H\right)`$ $`=`$ $`0,`$ ()
$`{\displaystyle \frac{1}{2}}\mathrm{sin}\left(2\theta _c\right)3\overline{K}_2\mathrm{sin}^2\theta _c\mathrm{cos}\theta _c+\overline{H}_c\mathrm{sin}\left(\theta _c\theta _H\right)`$ $`=`$ $`0,`$ ()
$`\mathrm{cos}\left(2\theta _c\right)3\overline{K}_2\left(2\mathrm{sin}\theta _c\mathrm{cos}^2\theta _c\mathrm{sin}^3\theta _c\right)+\overline{H}_c\mathrm{cos}\left(\theta _c\theta _H\right)`$ $`=`$ $`0.`$ ()
After some algebra, $`\overline{H}_c\left(\theta _H\right)`$ and $`\theta _c`$ are found to be
$`\overline{H}_c`$ $`=`$ $`{\displaystyle \frac{1}{\left[\left(\mathrm{sin}\theta _H\right)^{2/3}+\left|\mathrm{cos}\theta _H\right|^{2/3}\right]^{3/2}}}[1+3\overline{K}_2{\displaystyle \frac{1}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{1/2}}}`$ ()
$`+6\overline{K}_2{\displaystyle \frac{1}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{3/2}}}],`$
$`\mathrm{sin}^2\theta _c`$ $`=`$ $`{\displaystyle \frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}\left[12\overline{K}_2{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{1/2}}}4\overline{K}_2{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{3/2}}}\right].`$ ()
Now we consider the limiting case that the external magnetic field is slightly lower than the critical field, i.e., $`ϵ=1\overline{H}/\overline{H}_c1`$. At this practically interesting situation, the barrier height is low and the width is narrow, and therefore the tunneling rate in MQT or the tunnel splitting in MQC is large. Introducing $`\eta \theta _c\theta _0`$ $`\left(\left|\eta \right|1\text{ in the limit of }ϵ1\right)`$, expanding $`\left[d\overline{E}\left(\theta ,\varphi =0\right)/d\theta \right]_{\theta =\theta _0}=0`$ about $`\theta _c`$, and using the relations $`\left[d\overline{E}\left(\theta ,\varphi =0\right)/d\theta \right]_{\theta =\theta _c,\overline{H}=\overline{H}_c}=0`$ and $`\left[d^2\overline{E}\left(\theta ,\varphi =0\right)/d\theta ^2\right]_{\theta =\theta _c,\overline{H}=\overline{H}_c}=0`$, we obtain the approximation equation for $`\eta `$ in the order of $`ϵ^{3/2}`$,
$`ϵ\overline{H}_c\mathrm{sin}\left(\theta _c\theta _H\right)\eta ^2\left({\displaystyle \frac{3}{4}}\mathrm{sin}2\theta _c+3\overline{K}_2\mathrm{cos}3\theta _c\right)`$ (10)
$`+\eta \left[ϵ\overline{H}_c\mathrm{cos}\left(\theta _c\theta _H\right)+\eta ^2\left({\displaystyle \frac{1}{2}}\mathrm{cos}2\theta _c3\overline{K}_2\mathrm{sin}3\theta _c\right)\right]=0.`$ ()
Then $`\overline{E}(\theta ,\varphi )`$ reduces to the following equation in the limit of small $`ϵ`$,
$$\overline{E}(\delta ,\varphi )=2\overline{K}_2\mathrm{sin}^2\left(3\varphi /2\right)\mathrm{sin}^3\left(\theta _0+\delta \right)+\overline{H}_x\mathrm{sin}\left(\theta _0+\delta \right)\left(1\mathrm{cos}\varphi \right)+\overline{E}_1\left(\delta \right),$$
(11)
where $`\delta \theta \theta _0`$ $`\left(\left|\delta \right|1\text{ in the limit of }ϵ1\right)`$, and $`\overline{E}_1\left(\delta \right)`$ is a function of only $`\delta `$ given by
$`\overline{E}_1\left(\delta \right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\overline{H}_c\mathrm{sin}\left(\theta _c\theta _H\right)\overline{K}_2\left(\mathrm{cos}^3\theta _c{\displaystyle \frac{3}{2}}\mathrm{sin}^2\theta _c\mathrm{cos}\theta _c\right)\right]\left(3\delta ^2\eta \delta ^3\right)`$ ()
$`{\displaystyle \frac{1}{2}}\left[\overline{H}_c\mathrm{cos}\left(\theta _c\theta _H\right)3\overline{K}_2\left(\mathrm{sin}^3\theta _c4\mathrm{sin}\theta _c\mathrm{cos}^2\theta _c\right)\right]\left[\delta ^2\left(ϵ{\displaystyle \frac{3}{2}}\eta ^2\right)+\delta ^3\eta {\displaystyle \frac{1}{4}}\delta ^4\right]`$
$`{\displaystyle \frac{3}{2}}\overline{K}_2\left(\mathrm{sin}^3\theta _c4\mathrm{sin}\theta _c\mathrm{cos}^2\theta _c\right)\delta ^2ϵ.`$
In the following, we will investigate the tunneling behaviors of the magnetization vector in FM particles with trigonal crystal symmetry at different angle ranges of the external magnetic field as $`\theta _H=\pi /2`$ and $`\pi /2<\theta _H<\pi `$, respectively.
### A. $`\theta _H=\pi /2`$
For $`\theta _H=\pi /2`$, we have $`\theta _c=\pi /2`$ from Eq. (11b) and $`\eta =\sqrt{2ϵ}\left(1+9\overline{K}_2/2\right)`$ from Eq. (12). Then $`\overline{E}_1\left(\delta \right)`$ of Eq. (14) reduces to
$$\overline{E}_1\left(\delta \right)=\frac{1}{8}\delta ^2\left(\delta 2\eta \right)^2.$$
(14)
The plot of the effective potential $`\overline{E}_1\left(\delta \right)`$ as a function of $`\delta \left(=\theta \theta _0\right)`$ for $`\theta _H=\pi /2`$ is shown in Fig. 1. Now the problem is one of MQC, where the magnetization vector resonates coherently between the energetically degenerate easy directions at $`\delta =0`$ and $`\delta =2\sqrt{2ϵ}\left(1+9\overline{K}_2/2\right)`$ separated by a classically impenetrable barrier at $`\delta =\sqrt{2ϵ}\left(1+9\overline{K}_2/2\right)`$. Substituting Eq. (15) into the classical equations of motion, we obtain the classical solution called instanton as
$`\overline{\varphi }`$ $`=`$ $`iϵ\left(1+3\overline{K}_2+{\displaystyle \frac{1}{2}}ϵ\right){\displaystyle \frac{1}{\mathrm{cosh}^2\left(\overline{\omega }_c\overline{\tau }\right)}},`$ (15)
$`\overline{\delta }`$ $`=`$ $`\sqrt{2ϵ}\left(1+{\displaystyle \frac{9}{2}}\overline{K}_2\right)\left[1+\mathrm{tanh}\left(\overline{\omega }_c\overline{\tau }\right)\right],`$ ()
where $`\overline{\omega }_c=\sqrt{ϵ/2}\left(1+21\overline{K}_2/2ϵ/2\right)`$, $`\overline{\tau }=\omega _0\tau `$, and $`\omega _0=2K_1V/\mathrm{}S`$. We can calculate the classical action by integrating the Euclidean action Eq. (2) with the above classical trajectory, and the result is found to be
$$S_{cl}=\frac{2^{5/2}}{3}ϵ^{3/2}S\left(1+\frac{15}{2}\overline{K}_2+\frac{1}{2}ϵ\right).$$
(16)
Now we consider the transition exponent which is usually addressed by experiments. Transitions between two states in a bistable system or escaping from a metastable state can occur either due to the quantum tunneling or via the classical thermal activation. In the limit of temperature $`T0`$, the transitions are purely quantum-mechanical and the rate goes as $`\mathrm{\Gamma }\mathrm{exp}\left(S_{cl}\right)`$, with $`S_{cl}`$ being the classical action or the WKB exponent which is independent of temperature. As the temperature increases from zero, thermal effects enter in the quantum tunneling process. If the temperature is sufficiently high, the decay from a metastable state is determined by processes of thermal activation, and the transition rate follows the Arrhenius law, $`\mathrm{\Gamma }\mathrm{exp}\left(U/k_BT\right)`$, with $`k_B`$ being the Boltzmann constant and $`U`$ being the height of energy barrier between the two states. Because of the exponential dependence of the thermal rate on $`T`$, the temperature $`T_c`$ characterizing the crossover from quantum to thermal regime can be estimated as $`k_BT_c=U/S_{cl}`$. For a quasiparticle with the effective mass $`M`$ moving in one-dimensional potential $`U\left(x\right)`$, a more accurate definition of the crossover temperature in the absence of any dissipation was presented by Goldanskii, $`k_BT_c^{}=\mathrm{}\omega _b/2\pi `$, where $`\omega _b=\sqrt{U^{\prime \prime }\left(x_b\right)/M}`$ is the frequency of small oscillations near the bottom of the inverted potential, $`U\left(x\right)`$, and $`x_b`$ corresponds to the bottom of inverted potential. Below $`T_c^{}`$, thermally assisted quantum tunneling occurs from the excited levels, that further reduces to the quantum tunneling from the ground-state level as the temperature decreases to zero. Above $`T_c^{}`$, quantum tunneling effects are small and the transitions occur due to the thermal activation to the top of the barrier. For the MQT problem, i.e., the problem of decay from the metastable state, both $`T_c`$ and $`T_c^{}`$ can be used as the definition of the crossover temperature corresponding to the crossover from classical to quantum behavior since the quantum escaping from a metastable state is one process of incoherent tunneling. However, for the MQC problem, i.e., the problem of resonance between degenerate states, the situation is different. As the temperature growing from zero, three kinds of transitions should be taken into account: quantum coherence between the degenerate ground-state levels (coherent tunneling), quantum tunneling from the excited levels (thermally assisted tunneling or incoherent tunneling), and classical over-barrier transition (incoherent transition). Two kinds of crossover temperatures can be defined to distinguish the three regimes. The Goldanskii definition $`T_c^{}`$ for MQC problem corresponds to the crossover from quantum coherence between the degenerate ground-state levels (coherent tunneling) to quantum tunneling from the excited levels (thermally assisted tunneling or incoherent tunneling), while $`T_c`$ corresponds to the crossover from quantum coherence between the degenerate ground-state levels (coherent tunneling) to classical over-barrier transition (incoherent transition). Experiments involving magnetic relaxation and resonance measurements for various systems have shown either temperature-independent relaxation phenomena (in MQT) or a well-defined resonance (in MQC) below some crossover temperature, which strongly support the existence of quantum tunneling processes. And more recently, the crossover from quantum to classical behavior and associated phase transition have been investigated extensively in MQT and MQC in single-domain FM particles. It is noted that the sharpness of the crossover between thermal and quantum regimes also depends on the strength of the dissipation with environment. In the case of the low dissipation which is common for the magnetic systems, its effect on the crossover is small.
For the single-domain FM nanoparticle in a magnetic field applied at $`\theta _H=\pi /2`$, the magnetization vector resonates coherently between the energetically degenerate easy directions at $`\delta =0`$ and $`\delta =2\sqrt{2ϵ}\left(1+9\overline{K}_2/2\right)`$ separated by a classically impenetrable barrier at $`\delta =\sqrt{2ϵ}\left(1+9\overline{K}_2/2\right)`$, and the height of energy barrier is found to be: $`U=K_1Vϵ^2\left(1+18\overline{K}_2\right)`$. Then, equating $`S_{cl}`$ to $`U/k_BT`$, we obtain that the crossover from quantum coherence between the degenerate ground-state levels (coherent tunneling) to classical over-barrier transition (incoherent transition) occurs at
$$k_BT_c=\frac{3}{2^{5/2}}ϵ^{1/2}\frac{K_1V}{S}\left(1+\frac{21}{2}\overline{K}_2\frac{1}{2}ϵ\right).$$
(17)
For this MQC problem, the Goldanskii definition $`T_c^{}`$ corresponding the crossover from quantum coherence between the degenerate ground-state levels (coherent tunneling) to quantum tunneling from excited levels (thermally assisted tunneling or incoherent tunneling) becomes $`k_BT_c^{}=\mathrm{}\omega _b/2\pi `$, where $`\omega _b=\overline{\omega }_b\omega _0`$, with $`\overline{\omega }_b\sqrt{\overline{E}^{\prime \prime }\left(\delta _m\right)/M}`$ is the frequency of small oscillations of the magnetization vector near the bottom of the inverted potential, $`M^1=\left(1+12\overline{K}_2ϵ\right)`$, and $`\delta _m`$ is the position of the energy barrier. For the present case, $`\delta _m=\sqrt{2ϵ}\left(1+9\overline{K}_2/2\right)`$ and $`\overline{\omega }_b=\sqrt{ϵ}\left(1+21\overline{K}_2/2ϵ/2\right)=\sqrt{2}\overline{\omega }_c`$. Then it is easy to obtain that
$$k_BT_c^{}=\frac{1}{\pi }ϵ^{1/2}\frac{K_1V}{S}\left(1+\frac{21}{2}\overline{K}_2\frac{1}{2}ϵ\right).$$
(18)
The comparison of Eqs. (18) and (19) shows that $`T_c1.67T_c^{}`$, which is consistent with the physical interpretation for quantum-classical transition in the MQC problem.
It is noted that the quantum tunneling of the magnetization vector in single-domain FM nanoparticles are studied with the help of the instanton technique in the spin-coherent-state path-integral representation, which is semiclassical in nature, i.e., valid for large spins and in the continuum limit. Therefore, one should analyze the validity of the semiclassical approximation. It is well known that for this approach to be valid, the tunneling rate must be small, which indicates that the WKB exponent or the classical action $`S_{cl}1`$. Moreover, the energy $`\mathrm{}\omega _b`$ of zero-point oscillations around the minimum of the inverted potential $`\overline{E}_1\left(\delta \right)`$ should be sufficiently small compared to the height of the barrier, $`U=2K_1V\overline{E}_1\left(\delta _m\right)`$ . For the single-domain FM nanoparticle with trigonal crystal symmetry in a magnetic field applied at $`\theta _H=\pi /2`$, it is easy to show that the WKB exponent is approximately given by
$$B\frac{U}{\mathrm{}\omega _b}=\frac{1}{2}ϵ^{3/2}S\left(1+\frac{15}{2}\overline{K}_2+\frac{1}{2}ϵ\right),$$
(19)
which agrees up to the numerical factor with the result of the classical action in Eq. (17) obtained by applying the explicit instanton solution. For the typical values of parameters for single-domain FM nanoparticles, $`K_110^8`$ erg/cm<sup>3</sup>, $`K_210^5`$ erg/cm<sup>3</sup>, and the total spin $`S=10^6`$, we obtain that $`BU/\mathrm{}\omega _b15.8`$ from Eq. (20) and $`S_{cl}59.6`$ for $`ϵ=0.001`$ from Eq. (17). In this case the semiclassical approximation should be already rather good.
By applying the instanton technique for FM particles in the spin-coherent-state path-integral representation, we obtain the instanton’s contribution to the tunnel splitting as (for detailed calculation see Appendix A),
$$\mathrm{}\mathrm{\Delta }_0=\frac{2^{13/4}}{\pi ^{1/2}}\left(K_1V\right)ϵ^{5/4}S^{1/2}\left(1+\frac{57}{4}\overline{K}_2\frac{1}{4}ϵ\right)e^{S_{cl}},$$
(20)
where the WKB exponent or the classical action $`S_{cl}`$ has been presented in Eq. (17).
Now we apply the effective Hamiltonian approach to evaluate the ground-state tunnel splitting. For the present case, the effective Hamiltonian can be written as
$$H_{eff}=\left[\begin{array}{cc}0\hfill & \mathrm{}\mathrm{\Delta }_0\hfill \\ \mathrm{}\mathrm{\Delta }_0\hfill & 0\hfill \end{array}\right].$$
(21)
A simple diagonalization of $`H_{eff}`$ shows that the eigenvalues of this system are $`\pm \mathrm{}\mathrm{\Delta }_0`$. Therefore, the splitting of ground state due to resonant coherently quantum tunneling of the magnetization vector between energetically degenerate states is $`\mathrm{}\mathrm{\Delta }=2\mathrm{}\mathrm{\Delta }_0`$, where $`\mathrm{}\mathrm{\Delta }_0`$ is shown in Eq. (21) with Eq. (17) for single-domain FM particles with trigonal crystal symmetry in a magnetic field applied perpendicular to the anisotropy axis $`\left(\theta _H=\pi /2\right)`$.
### B. $`\pi /2<\theta _H<\pi `$
For $`\pi /2<\theta _H<\pi `$, the critical angle $`\theta _c`$ is in the range of $`0<\theta _c<\pi /2`$, and $`\eta \sqrt{2ϵ/3}`$. Then $`\overline{E}_1\left(\delta \right)`$ of Eq. (14) reduces to
$$\overline{E}_1\left(\delta \right)=\frac{1}{2}\frac{\left|\mathrm{cot}\theta _H\right|^{1/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}\left[1\frac{15}{2}\overline{K}_2\frac{1}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{1/2}}\right]\left(\sqrt{6ϵ}\delta ^2\delta ^3\right).$$
(22)
The dependence of the effective potential $`\overline{E}_1\left(\delta \right)`$ on $`\delta \left(=\theta \theta _0\right)`$ for $`\theta _H=3\pi /4`$ is plotted in Fig. 2. Here, $`\overline{K}_2=0.001`$. Now the problem becomes one of MQT, where the magnetization vector escapes from the metastable state at $`\delta =0`$, $`\varphi =0`$ through the barrier by quantum tunneling. Substituting Eq. (23) into the classical equations of motion, the classical solution called bounce is found to be
$`\overline{\varphi }`$ $`=`$ $`i\left(6ϵ\right)^{3/4}|\mathrm{cot}\theta _H|^{1/6}(1+|\mathrm{cot}\theta _H|^{2/3})^{1/2}[1+{\displaystyle \frac{ϵ}{2}}{\displaystyle \frac{9}{2}}\overline{K}_2(1+|\mathrm{cot}\theta _H|^{2/3})^{1/2}`$ (24)
$`+{\displaystyle \frac{\overline{K}_2}{4}}{\displaystyle \frac{2\left|\mathrm{cot}\theta _H\right|^{2/3}9}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{1/2}}}+\overline{K}_2{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}3}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{3/2}}}]{\displaystyle \frac{\mathrm{sinh}\left(\overline{\omega }_c\overline{\tau }\right)}{\mathrm{cosh}^3\left(\overline{\omega }_c\overline{\tau }\right)}},`$
$`\overline{\delta }`$ $`=`$ $`\sqrt{6ϵ}/\mathrm{cosh}^2\left(\overline{\omega }_c\overline{\tau }\right),`$ ()
which corresponds to the variation of $`\delta `$ from $`\delta =0`$ at $`\tau =\mathrm{}`$ to the turning point $`\delta =\sqrt{6ϵ}`$ at $`\tau =0`$, and then back to $`\delta =0`$ at $`\tau =+\mathrm{}`$, where
$`\overline{\omega }_c`$ $`=`$ $`3^{1/4}\times 2^{3/4}ϵ^{1/4}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{1/6}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}[1{\displaystyle \frac{ϵ}{2}}+{\displaystyle \frac{9}{2}}\overline{K}_2(1+|\mathrm{cot}\theta _H|^{2/3})^{1/2}`$ (26)
$`+{\displaystyle \frac{\overline{K}_2}{4}}{\displaystyle \frac{2\left|\mathrm{cot}\theta _H\right|^{2/3}21}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{1/2}}}+\overline{K}_2{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}+3}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{3/2}}}].`$
The associated classical action is then given by
$`S_{cl}`$ $`=`$ $`{\displaystyle \frac{3^{1/4}\times 2^{17/4}}{5}}Sϵ^{5/4}|\mathrm{cot}\theta _H|^{1/6}[1+{\displaystyle \frac{ϵ}{2}}{\displaystyle \frac{9}{2}}\overline{K}_2(1+|\mathrm{cot}\theta _H|^{2/3})^{1/2}`$ ()
$`{\displaystyle \frac{\overline{K}_2}{2}}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}+9/2}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{1/2}}}\overline{K}_2{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}+3}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{3/2}}}].`$
For this case, the barrier height is
$`U`$ $`=`$ $`2K_1V\overline{E}_1\left(\delta _m\right)`$
$`=`$ $`{\displaystyle \frac{2^{7/2}}{3^{3/2}}}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{1/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}\left[1{\displaystyle \frac{15}{2}}\overline{K}_2{\displaystyle \frac{1}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{1/2}}}\right]ϵ^{3/2}\left(K_1V\right),`$
at $`\delta _m=2\left(6ϵ\right)^{1/2}/3`$, and the frequency of small oscillations of the magnetization vector around the bottom of the metastable well is
$`\overline{\omega }_b`$ $`=`$ $`3^{1/4}\times 2^{1/4}ϵ^{1/4}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{1/6}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}[1{\displaystyle \frac{ϵ}{2}}+{\displaystyle \frac{9}{2}}\overline{K}_2(1+|\mathrm{cot}\theta _H|^{2/3})^{1/2}`$
$`+{\displaystyle \frac{\overline{K}_2}{4}}{\displaystyle \frac{2\left|\mathrm{cot}\theta _H\right|^{2/3}21}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{1/2}}}+\overline{K}_2{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}+3}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{3/2}}}]`$
$`=`$ $`2\overline{\omega }_c.`$
Then the WKB exponent or the classical action $`B`$ is approximately given by
$`B`$ $``$ $`{\displaystyle \frac{U}{\mathrm{}\omega _b}}`$ (28)
$`=`$ $`{\displaystyle \frac{2^{9/4}}{3^{7/4}}}Sϵ^{5/4}|\mathrm{cot}\theta _H|^{1/6}[1+{\displaystyle \frac{ϵ}{2}}{\displaystyle \frac{9}{2}}\overline{K}_2(1+|\mathrm{cot}\theta _H|^{2/3})^{1/2}`$ ()
$`{\displaystyle \frac{\overline{K}_2}{2}}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}+9/2}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{1/2}}}\overline{K}_2{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}+3}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{3/2}}}],`$
which is consistent with Eq. (25) up to the numerical factor. After a simple calculation, we obtain the crossover temperature as
$`k_BT_c`$ $`=`$ $`{\displaystyle \frac{5}{2^{3/4}\times 3^{7/4}}}ϵ^{1/4}{\displaystyle \frac{K_1V}{S}}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{1/6}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}[1{\displaystyle \frac{ϵ}{2}}+{\displaystyle \frac{9}{2}}\overline{K}_2(1+|\mathrm{cot}\theta _H|^{2/3})^{1/2}`$ ()
$`+{\displaystyle \frac{\overline{K}_2}{2}}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}21/2}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{1/2}}}+\overline{K}_2{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}+3}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{3/2}}}],`$
corresponding to the transition from quantum to thermal regime. For a nanometer-scale single-domain FM particle, the typical values of parameters for the magnetic anisotropy coefficients are $`K_1=10^8`$ erg/cm<sup>3</sup>, and $`K_2=10^5`$ erg/cm<sup>3</sup>. The radius of the FM particle is about 12 nm and the sublattice spin is $`10^6`$. If $`ϵ=0.001`$, we obtain that $`T_c\left(135^{}\right)203`$mK corresponding to the crossover from quantum to classical regime, which compares well with the experimental result of $`0.31`$K on single-domain FM nanoparticles of Barium ferrite (BaFeCoTiO). Note that, even for $`ϵ`$ as small as $`10^3`$, the angle corresponding to an appreciable change of the orientation of the magnetization vector by quantum tunneling is $`\delta _2=\sqrt{6ϵ}`$ rad$`>4^{}`$.
The classical action $`S_{cl}`$ can be obtained by solving numerically the equations of motion (4a) and (4b). In Fig. 3 we present the $`\theta _H`$ dependence of $`S_{cl}`$ with $`ϵ=0.001`$ and $`\overline{K}_2=0.001`$ for $`\pi /2<\theta _H<\pi `$ by numerical and analytical calculations, respectively. As is noted in the figure, the analytical result obtained from Eq. (25) is almost valid in the whole range of angles $`\pi /2<\theta _H<\pi `$.
By applying the formulas in Ref. 12, and using Eq. (25) for the WKB exponent or the classical action, we obtain the tunneling rate $`\mathrm{\Gamma }`$ of the magnetization vector in single-domain FM nanoparticles with trigonal crystal symmetry in a magnetic field applied in the range of $`\pi /2<\theta _H<\pi `$ as
$`\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \frac{2^{31/8}\times 3^{7/8}}{\pi ^{1/2}}}{\displaystyle \frac{V}{\mathrm{}}}K_1S^{1/2}ϵ^{7/8}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{1/4}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}[1{\displaystyle \frac{ϵ}{4}}+{\displaystyle \frac{9}{4}}\overline{K}_2(1+|\mathrm{cot}\theta _H|^{2/3})^{1/2}`$ ()
$`+{\displaystyle \frac{\overline{K}_2}{4}}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}51/2}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{1/2}}}+{\displaystyle \frac{\overline{K}_2}{2}}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}+3}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{3/2}}}]e^{S_{cl}}.`$
## IV. MQC and MQT for hexagonal crystal symmetry
In this section, we study the quantum tunneling of the magnetization vector in single-domain FM particles with hexagonal crystal symmetry whose magnetocrystalline anisotropy energy $`E_a(\theta ,\varphi )`$ at zero magnetic field can be written as
$$E_a(\theta ,\varphi )=K_1\mathrm{sin}^2\theta +K_2\mathrm{sin}^4\theta +K_3\mathrm{sin}^6\theta K_3^{}\mathrm{sin}^6\theta \mathrm{cos}\left(6\varphi \right),$$
(32)
where $`K_1`$, $`K_2`$, $`K_3`$, and $`K_3^{}`$ are the magnetic anisotropic coefficients. The easy axes are $`\pm \widehat{z}`$ for $`K_1>0`$. When we apply an external magnetic field at an arbitrarily directed angle in the $`ZX`$ plane, the total energy of this system is given by
$$E(\theta ,\varphi )=E_a(\theta ,\varphi )M_0H_x\mathrm{sin}\theta \mathrm{cos}\varphi M_0H_z\mathrm{cos}\theta +E_0,$$
(33)
By choosing $`K_3^{}>0`$, we take $`\varphi =0`$ to be the easy plane, at which the potential energy can be written in terms of the dimensionless parameters as
$$\overline{E}\left(\theta ,\varphi =0\right)=\frac{1}{2}\mathrm{sin}^2\theta +\overline{K}_2\mathrm{sin}^4\theta +\left(\overline{K}_3\overline{K}_3^{}\right)\mathrm{sin}^6\theta \overline{H}\mathrm{cos}\left(\theta \theta _H\right)+\overline{E}_0,$$
(34)
where $`\overline{K}_3=K_3/2K_1`$ and $`\overline{K}_3^{}=K_3^{}/2K_1`$.
Then the initial angle $`\theta _0`$ is determined by $`\left[d\overline{E}(\theta ,0)/d\theta \right]_{\theta =\theta _0}=0`$, and the critical angle $`\theta _c`$ and the dimensionless critical field $`\overline{H}_c`$ by both $`\left[d\overline{E}(\theta ,0)/d\theta \right]_{\theta =\theta _c,\overline{H}=\overline{H}_c}=0`$ and $`\left[d^2\overline{E}(\theta ,0)/d\theta ^2\right]_{\theta =\theta _c,\overline{H}=\overline{H}_c}=0`$, which leads to
$`{\displaystyle \frac{1}{2}}\mathrm{sin}\left(2\theta _0\right)+\overline{H}\mathrm{sin}\left(\theta _0\theta _H\right)+4\overline{K}_2\mathrm{sin}^4\theta _0+6\left(\overline{K}_3\overline{K}_3^{}\right)\mathrm{sin}^5\theta _0\mathrm{cos}\theta _0=0,`$ ()
$`{\displaystyle \frac{1}{2}}\mathrm{sin}\left(2\theta _c\right)+\overline{H}_c\mathrm{sin}\left(\theta _c\theta _H\right)+4\overline{K}_2\mathrm{sin}^4\theta _c+6\left(\overline{K}_3\overline{K}_3^{}\right)\mathrm{sin}^5\theta _c\mathrm{cos}\theta _c=0,`$ ()
$`\mathrm{cos}\left(2\theta _c\right)+\overline{H}_c\mathrm{cos}\left(\theta _c\theta _H\right)+4\overline{K}_2\left(3\mathrm{sin}^2\theta _c\mathrm{cos}^2\theta _c\mathrm{sin}^4\theta _c\right)`$ (35)
$`+6\left(\overline{K}_3\overline{K}_3^{}\right)\left(5\mathrm{sin}^4\theta _c\mathrm{cos}^2\theta _c\mathrm{sin}^6\theta _c\right)=0,`$ ()
Under the assumption that $`\left|\overline{K}_2\right|`$, $`\left|\overline{K}_3\overline{K}_3^{}\right|1`$, we obtain the dimensionless critical field $`\overline{H}_c`$ as
$$\overline{H}_c=\frac{1}{\left[\left(\mathrm{sin}\theta _H\right)^{2/3}+\left|\mathrm{cos}\theta _H\right|^{2/3}\right]^{3/2}}\left[1+\frac{4\overline{K}_2}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}+\frac{6\left(\overline{K}_3\overline{K}_3^{}\right)}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}\right].$$
(36)
In the limit of small $`ϵ=1\overline{H}/\overline{H}_c`$, Eq. (32a) becomes
$`ϵ\overline{H}_c\mathrm{sin}(\theta _c\theta _H)+\eta ^2[(3/2)\overline{H}_c\mathrm{sin}(\theta _c\theta _H)+3\overline{K}_2\mathrm{sin}\left(4\theta _c\right)`$ (37)
$`+12(\overline{K}_3\overline{K}_3^{})\mathrm{sin}^3\theta _c\mathrm{cos}\theta _c(58\mathrm{sin}^2\theta _c)]+\eta \{ϵ\overline{H}_c\mathrm{cos}(\theta _c\theta _H)`$ (38)
$`\eta ^2[(1/2)\overline{H}_c\mathrm{cos}(\theta _c\theta _H)+4\overline{K}_2\mathrm{cos}\left(4\theta _c\right)`$ (39)
$`+12(\overline{K}_3\overline{K}_3^{})\mathrm{sin}^2\theta _c(520\mathrm{sin}^2\theta _c+16\mathrm{sin}^4\theta _c)]\}=0,`$ ()
where $`\eta \theta _c\theta _0`$ which is small for $`ϵ1`$. By introducing a small variable $`\delta \theta \theta _0`$ $`\left(\left|\delta \right|1\text{ in the limit of }ϵ1\right)`$, the total energy becomes
$$\overline{E}(\delta ,\varphi )=\overline{K}_3^{}\left[1\mathrm{cos}\left(6\varphi \right)\right]\mathrm{sin}^6\left(\theta _0+\delta \right)+\overline{H}_x\left(1\mathrm{cos}\varphi \right)\mathrm{sin}\left(\theta _0+\delta \right)+\overline{E}_1\left(\delta \right),$$
(40)
where $`\overline{E}_1\left(\delta \right)`$ is a function of only $`\delta `$ given by
$`\overline{E}_1\left(\delta \right)`$ $`=`$ $`\left[{\displaystyle \frac{1}{2}}\overline{H}_c\mathrm{sin}\left(\theta _c\theta _H\right)+\overline{K}_2\mathrm{sin}\left(4\theta _c\right)+4\left(\overline{K}_3\overline{K}_3^{}\right)\left(5\mathrm{sin}^3\theta _c\mathrm{cos}^3\theta _c3\mathrm{sin}^5\theta _c\mathrm{cos}\theta _c\right)\right]`$ ()
$`\times (\delta ^33\delta ^2\eta )+[{\displaystyle \frac{1}{8}}\overline{H}_c\mathrm{cos}(\theta _c\theta _H)+\overline{K}_2\mathrm{cos}\left(4\theta _c\right)+3(\overline{K}_3\overline{K}_3^{})\mathrm{sin}^2\theta _c(\mathrm{sin}^4\theta _c`$
$`10\mathrm{sin}^2\theta _c\mathrm{cos}^2\theta _c+5\mathrm{cos}^4\theta _c)\left]\right(\delta ^44\delta ^3\eta +6\delta ^2\eta ^24\delta ^2ϵ)+ϵ\delta ^2[4\overline{K}_2\mathrm{cos}\left(4\theta _c\right)`$
$`+12(\overline{K}_3\overline{K}_3^{})\mathrm{sin}^2\theta _c(\mathrm{sin}^4\theta _c10\mathrm{sin}^2\theta _c\mathrm{cos}^2\theta _c+5\mathrm{cos}^4\theta _c)].`$
In the following we investigate the MQC and MQT of the magnetization vector in FM particles with hexagonal crystal symmetry for different angle ranges of the external magnetic field: $`\theta _H=\pi /2`$, $`\pi /2+O\left(ϵ^{3/2}\right)<\theta _H<\pi O\left(ϵ^{3/2}\right)`$, and $`\theta _H=\pi `$, respectively.
### A. $`\theta _H=\pi /2`$
For $`\theta _H=\pi /2`$, i.e., the external magnetic field is applied perpendicular to the anisotropy axis, we obtain that $`\theta _c=\pi /2`$ and $`\eta =\sqrt{2ϵ}\left[14\overline{K}_212\left(\overline{K}_3\overline{K}_3^{}\right)\right]`$. Now $`\overline{E}_1\left(\delta \right)`$ becomes
$$\overline{E}_1\left(\delta \right)=\frac{1}{8}\left[1+12\overline{K}_2+30\left(\overline{K}_3\overline{K}_3^{}\right)\right]\delta ^2\left\{\delta 2\sqrt{2ϵ}\left[14\overline{K}_212\left(\overline{K}_3\overline{K}_3^{}\right)\right]\right\}^2.$$
(44)
Substituting Eq. (37) into the classical equations of motion, we obtain the following instanton solution
$`\overline{\varphi }`$ $`=`$ $`iϵ\left[1+{\displaystyle \frac{ϵ}{2}}4\overline{K}_218\overline{K}_3^{}6\left(\overline{K}_3\overline{K}_3^{}\right)\right]{\displaystyle \frac{1}{\mathrm{cosh}^2\left(\overline{\omega }_c\overline{\tau }\right)}},`$ (45)
$`\overline{\delta }`$ $`=`$ $`\sqrt{2ϵ}\left[14\overline{K}_212\left(\overline{K}_3\overline{K}_3^{}\right)\right]\left[1+\mathrm{tanh}\left(\overline{\omega }_c\overline{\tau }\right)\right],`$ ()
which corresponds to the variation of $`\delta `$ from $`\delta =0`$ at $`\tau =\mathrm{}`$ to $`\delta =2\sqrt{2ϵ}\left[14\overline{K}_212\left(\overline{K}_3\overline{K}_3^{}\right)\right]`$ at $`\tau =+\mathrm{}`$, where
$`\overline{\omega }_c=\sqrt{{\displaystyle \frac{ϵ}{2}}}\left[1{\displaystyle \frac{ϵ}{2}}+4\overline{K}_2+18\overline{K}_3^{}+6\left(\overline{K}_3\overline{K}_3^{}\right)\right].`$
We can calculate the classical action by integrating the Euclidean action of Eq. (2) with the above instanton solution, and the result is found to be
$$S_{cl}=\frac{2^{5/2}}{3}Sϵ^{3/2}\left[1+\frac{ϵ}{2}8\overline{K}_218\overline{K}_3^{}24\left(\overline{K}_3\overline{K}_3^{}\right)\right].$$
(46)
From Eq. (37) we obtain that the height of barrier is $`U=2K_1V\overline{E}_1\left(\delta _m\right)=K_1Vϵ^2\left[14\overline{K}_218\left(\overline{K}_3\overline{K}_3^{}\right)\right]`$ at $`\delta _m=\sqrt{2ϵ}\left[14\overline{K}_212\left(\overline{K}_3\overline{K}_3^{}\right)\right]`$, and the oscillation frequency around the minimum of the inverted potential $`\overline{E}_1\left(\delta \right)`$ is
$`\overline{\omega }_b=\sqrt{ϵ}\left[1{\displaystyle \frac{ϵ}{2}}+4\overline{K}_2+18\overline{K}_3^{}+6\left(\overline{K}_3\overline{K}_3^{}\right)\right]=\sqrt{2}\overline{\omega }_c.`$
Then the WKB exponent is approximately given by
$`B`$ $``$ $`{\displaystyle \frac{U}{\mathrm{}\omega _b}}`$ (47)
$`=`$ $`{\displaystyle \frac{1}{2}}Sϵ^{3/2}\left[1+{\displaystyle \frac{ϵ}{2}}8\overline{K}_218\overline{K}_3^{}24\left(\overline{K}_3\overline{K}_3^{}\right)\right],`$ ()
which agrees up to the numerical factor with Eq. (39) obtained by applying the explicit instanton solution. The temperature corresponding to the crossover from the quantum coherence between the degenerate ground-state levels (coherent tunneling) to the classical over-barrier transition (incoherent transition) is found to be
$$k_BT_c=\frac{3}{2^{5/2}}ϵ^{1/2}\frac{K_1V}{S}\left[1\frac{ϵ}{2}+4\overline{K}_2+18\overline{K}_3^{}+6\left(\overline{K}_3\overline{K}_3^{}\right)\right],$$
(48)
and the temperature corresponding to the crossover from quantum coherence between the degenerate ground-state levels (coherent tunneling) to quantum tunneling from excited levels (thermally assisted tunneling or incoherent tunneling) is found to be
$$k_BT_c^{}=\frac{1}{\pi }ϵ^{1/2}\frac{K_1V}{S}\left[1\frac{ϵ}{2}+4\overline{K}_2+18\overline{K}_3^{}+6\left(\overline{K}_3\overline{K}_3^{}\right)\right].$$
(49)
By applying the instanton technique for single-domain FM particles in the spin-coherent-state path-integral representation, we obtain the instanton’s contribution to the tunnel splitting, $`\mathrm{}\mathrm{\Delta }_0`$ as
$$\mathrm{}\mathrm{\Delta }_0=\frac{2^{13/4}}{\pi ^{1/2}}\left(VK_1\right)S^{1/2}ϵ^{5/4}\left[1\frac{ϵ}{4}+9\overline{K}_3^{}6\left(\overline{K}_3\overline{K}_3^{}\right)\right]e^{S_{cl}},$$
(50)
where the WKB exponent or the classical action $`S_{cl}`$ is clearly shown in Eq. (39). Then the splitting of ground state due to resonant coherently quantum tunneling of the magnetization vector between energetically degenerate states is found to be $`\mathrm{}\mathrm{\Delta }=2\mathrm{}\mathrm{\Delta }_0`$ for FM particles with hexagonal crystal symmetry in a magnetic field applied perpendicular to the anisotropy axis $`\left(\theta _H=\pi /2\right)`$ with the help of the effective Hamiltonian approach.
### B. $`\pi /2+O\left(ϵ^{3/2}\right)<\theta _H<\pi O\left(ϵ^{3/2}\right)`$
For this case, $`\eta \sqrt{2ϵ}/3`$ and the critical angle $`\theta _c`$ is found to be
$`\mathrm{sin}\theta _c={\displaystyle \frac{1}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{1/2}}}\left[1+{\displaystyle \frac{8}{3}}\overline{K}_2{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}+8\left(\overline{K}_3\overline{K}_3^{}\right){\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}\right].`$
Now $`\overline{E}_1\left(\delta \right)`$ becomes
$`\overline{E}_1\left(\delta \right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{1/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}[1{\displaystyle \frac{4}{3}}\overline{K}_2{\displaystyle \frac{74\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}`$ ()
$`+2(\overline{K}_3\overline{K}_3^{}){\displaystyle \frac{1116\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}\left]\right(\sqrt{6ϵ}\delta ^2\delta ^3).`$
Then the classical equations of motion have the following bounce solution
$`\overline{\varphi }`$ $`=`$ $`i\left(6ϵ\right)^{3/4}|\mathrm{cot}\theta _H|^{1/6}(1+|\mathrm{cot}\theta _H|^{2/3})^{1/2}[1+{\displaystyle \frac{ϵ}{2}}{\displaystyle \frac{4}{3}}\overline{K}_2{\displaystyle \frac{5\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}`$ (53)
$`18\overline{K}_3^{}{\displaystyle \frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}2(\overline{K}_3\overline{K}_3^{}){\displaystyle \frac{76\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}]{\displaystyle \frac{\mathrm{sinh}\left(\overline{\omega }_c\overline{\tau }\right)}{\mathrm{cosh}^3\left(\overline{\omega }_c\overline{\tau }\right)}},`$
$`\overline{\delta }`$ $`=`$ $`\sqrt{6ϵ}/\mathrm{cosh}^2\left(\overline{\omega }_c\overline{\tau }\right),`$ ()
corresponding to the variation of $`\delta `$ from $`\delta =0`$ at $`\tau =\mathrm{}`$ to the turning point $`\delta =\sqrt{6ϵ}`$ at $`\tau =0`$, and then back to $`\delta =0`$ at $`\tau =+\mathrm{}`$, where
$`\overline{\omega }_c`$ $`=`$ $`\left({\displaystyle \frac{3}{8}}\right)^{1/4}ϵ^{1/4}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{1/6}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}[1{\displaystyle \frac{ϵ}{2}}+{\displaystyle \frac{4}{3}}\overline{K}_2{\displaystyle \frac{53\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}`$
$`+18\overline{K}_3^{}{\displaystyle \frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}+2(\overline{K}_3\overline{K}_3^{}){\displaystyle \frac{710\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}].`$
The classical action associated with this bounce solution is found to be
$`S_{cl}`$ $`=`$ $`{\displaystyle \frac{2^{17/4}\times 3^{1/4}}{5}}Sϵ^{5/4}|\mathrm{cot}\theta _H|^{1/6}[1+{\displaystyle \frac{ϵ}{2}}+{\displaystyle \frac{4}{3}}\overline{K}_2{\displaystyle \frac{2\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}`$ ()
$`18\overline{K}_3^{}{\displaystyle \frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}+4(\overline{K}_3\overline{K}_3^{}){\displaystyle \frac{23\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}].`$
For this case, the barrier height $`U\left(=2K_1V\overline{E}_1\left(\delta _m=2\sqrt{6ϵ}/3\right)\right)`$ is given by
$`U`$ $`=`$ $`{\displaystyle \frac{2^{7/2}}{3^{3/2}}}\left(K_1V\right)ϵ^{3/2}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{1/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}[1+{\displaystyle \frac{4}{3}}\overline{K}_2{\displaystyle \frac{74\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}`$
$`+2(\overline{K}_3\overline{K}_3^{}){\displaystyle \frac{1116\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}],`$
and the frequency of small oscillations of the magnetization vector around the minimum of the inverted potential $`\overline{E}_1\left(\delta \right)`$ is
$`\overline{\omega }_b`$ $`=`$ $`3^{1/4}\times 2^{1/4}ϵ^{1/4}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{1/6}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}[1{\displaystyle \frac{ϵ}{2}}+{\displaystyle \frac{4}{3}}\overline{K}_2{\displaystyle \frac{53\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}`$
$`+18\overline{K}_3^{}{\displaystyle \frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}+2(\overline{K}_3\overline{K}_3^{}){\displaystyle \frac{710\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}]`$
$`=`$ $`2\overline{\omega }_c.`$
Then the WKB exponent is approximately given by
$`B`$ $``$ $`{\displaystyle \frac{U}{\mathrm{}\omega _b}}`$ (55)
$`=`$ $`{\displaystyle \frac{2^{9/4}}{3^{7/4}}}Sϵ^{5/4}|\mathrm{cot}\theta _H|^{1/6}[1+{\displaystyle \frac{ϵ}{2}}+{\displaystyle \frac{4}{3}}\overline{K}_2{\displaystyle \frac{2\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}`$ ()
$`18\overline{K}_3^{}{\displaystyle \frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}+4(\overline{K}_3\overline{K}_3^{}){\displaystyle \frac{23\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}],`$
which agrees with Eq. (46) up to the numerical factor. Equating the classical action $`S_{cl}`$ to $`U/k_BT_c`$, where $`U`$ is the barrier height, we obtain that the crossover from quantum to classical behavior occurs at
$`k_BT_c`$ $`=`$ $`{\displaystyle \frac{5}{2^{3/4}\times 3^{7/4}}}ϵ^{1/4}{\displaystyle \frac{K_1V}{S}}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{1/6}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}[1{\displaystyle \frac{ϵ}{2}}+{\displaystyle \frac{4}{3}}\overline{K}_2{\displaystyle \frac{53\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}`$ ()
$`+18\overline{K}_3^{}{\displaystyle \frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}+2(\overline{K}_3\overline{K}_3^{}){\displaystyle \frac{710\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}].`$
Based on the instanton technique, we obtain the tunneling rate corresponding to the escaping of the magnetization vector from the metastable state for single-domain FM nanoparticles with hexagonal crystal symmetry in a magnetic field applied in the range of $`\pi /2+O\left(ϵ^{3/2}\right)<\theta _H<\pi O\left(ϵ^{3/2}\right)`$ as the following equation,
$`\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \frac{2^{31/8}\times 3^{7/8}}{\pi ^{1/2}}}{\displaystyle \frac{V}{\mathrm{}}}K_1S^{1/2}ϵ^{7/8}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{1/4}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}[1{\displaystyle \frac{ϵ}{4}}+9\overline{K}_3^{}{\displaystyle \frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}`$ ()
$`{\displaystyle \frac{2}{3}}\overline{K}_2{\displaystyle \frac{127\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}`$
$`+2(\overline{K}_3\overline{K}_3^{}){\displaystyle \frac{913\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}]e^{S_{cl}},`$
where the WKB exponent or the classical action $`S_{cl}`$ has been clearly shown in Eq. (46).
### C. $`\theta _H=\pi `$
Finally, we study the MQT of the magnetization vector corresponding to the escaping from the metastable state in single-domain FM nanoparticles with hexagonal crystal symmetry in a magnetic field applied at $`\theta _H=\pi `$, i.e., antiparallel to the anisotropy axis. Now the total energy become
$`\overline{E}(\delta ,\varphi )`$ $`=`$ $`\overline{K}_3^{}\left[1\mathrm{cos}\left(6\varphi \right)\right]\delta ^6+{\displaystyle \frac{1}{2}}\delta ^2\left[ϵ{\displaystyle \frac{1}{4}}\left(18\overline{K}_2\right)\delta ^2\right]`$ ()
$`{\displaystyle \frac{1}{24}}\delta ^4\left\{ϵ{\displaystyle \frac{1}{2}}\left[132\overline{K}_2+48\left(\overline{K}_3\overline{K}_3^{}\right)\right]\delta ^2\right\}.`$
The classical equations of motion have the bounce solution
$`\overline{\varphi }`$ $`=`$ $`i\overline{\omega }_c\overline{\tau }+{\displaystyle \frac{n\pi }{3}},`$ (61)
$`\overline{\delta }`$ $`=`$ $`\sqrt{{\displaystyle \frac{4ϵ}{18\overline{K}_2ϵ\left[32\overline{K}_3^{}\left(1\mathrm{cosh}\left(6\overline{\omega }_c\overline{\tau }\right)\right)+\frac{1}{3}\left(148\overline{K}_2+96\left(\overline{K}_3\overline{K}_3^{}\right)\right)\right]}}},`$ ()
where $`n=0,1,2,3,4,5`$, and $`\overline{\omega }_c=ϵ`$. The corresponding classical action is found to be
$$S_{cl}=\frac{2}{3}Sϵ\frac{1}{\mathrm{\Delta }_1}\mathrm{ln}\left(\frac{2\mathrm{\Delta }_1}{\mathrm{\Delta }_2}\right),$$
(62)
with
$$\mathrm{\Delta }_1=18\overline{K}_232\overline{K}_3^{}ϵ\frac{1}{3}ϵ\left[148\overline{K}_2+96\left(\overline{K}_3\overline{K}_3^{}\right)\right],$$
(63)
and
$$\mathrm{\Delta }_2=32\overline{K}_3^{}ϵ.$$
(64)
According to the formulas in Ref. 12, we obtain the tunneling rate of the magnetization vector escaping from the metastable state for single-domain FM nanoparticles with hexagonal crystal symmetry in a magnetic field applied antiparallel to the anisotropy axis $`\left(\theta _H=\pi \right)`$ as
$`\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \frac{2^{13/2}\times 3^{1/2}}{\pi ^{1/2}}}{\displaystyle \frac{V}{\mathrm{}}}K_1S^{1/2}ϵ\left(1+4\overline{K}_2\right)`$ ()
$`{\displaystyle \frac{1}{116\overline{K}_264\overline{K}_3^{}ϵ\frac{2}{3}\left[148\overline{K}_2+96\left(\overline{K}_3\overline{K}_3^{}\right)\right]}}e^{S_{cl}},`$
where the WKB exponent or the classical action $`S_{cl}`$ is shown in Eq. (52). Eq. (50) shows that in this case $`\left|\varphi \right|1`$ is not valid, and therefore the problem can not be reduced to the one-dimensional motion problem. And the effective potential energy and the effective mass in one-dimensional form are not appropriate for the present case.
Now we discuss the range of angles that Eq. (46) is valid. Introducing $`\theta _1=\theta _H\pi /2`$ and $`\theta _2=\pi \theta _H`$, from Eqs. (39), (46) and (52), we find that $`\theta _1\left(5^6\times 2^{21/2}\times 3^{15/2}\right)ϵ^{3/2}`$ and $`\theta _2\left(5^6\times 2^{39/2}\times 3^{15/2}\right)ϵ^{3/2}`$. This means that Eq. (46) is almost valid in a wide range of angles $`91^{}\theta _H179^\text{ }`$ for $`ϵ=0.001`$.
For the single-domain FM nanoparticle with hexagonal crystal symmetry in the presence of an external magnetic field at arbitrarily directed angle, by using Eqs. (39) and (43) for $`\theta _H=\pi /2`$, Eqs. (46) and (49) for $`\pi /2+O\left(ϵ^{3/2}\right)<\theta _H<\pi O\left(ϵ^{3/2}\right)`$, and Eqs. (52) and (55) for $`\theta _H=\pi `$, we obtain the ground-state tunnel splitting for MQC and the tunneling rate for MQT of the magnetization vector. Our results show that the tunnel splitting and the tunneling rate depend on the orientation of the external magnetic field distinctly. When $`\theta _H=\pi /2`$, the magnetic field is applied perpendicular to the anisotropy axis, and when $`\theta _H=\pi `$, the field is antiparallel to the anisotropy axis. It is found that even a very small misalignment of the field with the above two orientations can completely change the results of tunneling rates. Another interesting observation concerns the dependence of the WKB exponent or the classical action with the strength of the external magnetic field. In a wide range of angles, the $`ϵ\left(=1\overline{H}/\overline{H}_c\right)`$ dependence of the WKB exponent $`S_{cl}`$ is given by $`ϵ^{5/4}`$, not $`ϵ^{3/2}`$ for $`\theta _H=\pi /2`$, and $`ϵ`$ for $`\theta _H=\pi `$. Therefore, both the orientation and the strength of the external magnetic field are the control parameters for the experimental test for MQT and MQC of the magnetization vector in single-domain FM nanoparticles.
## V. Conclusions
In summary we have investigated the tunneling behaviors of the magnetization vector in single-domain FM nanoparticles in the presence of an external magnetic field at arbitrarily directed angle. We consider the magnetocrystalline anisotropy with the trigonal crystal symmetry and that with the hexagonal crystal symmetry. By applying the instanton technique in the spin-coherent-state path-integral representation, we obtain both the WKB exponent and the preexponential factors in the tunnel splitting between energetically degenerate states in MQC and the tunneling rate escaping from a metastable state in MQT of the magnetization vector in the low barrier limit for the external magnetic field perpendicular to the easy axis $`\left(\theta _H=\pi /2\right)`$, for the field antiparallel to the initial easy axis $`\left(\theta _H=\pi \right)`$, and for the field at an angle between these two orientations $`\left(\pi /2+O\left(ϵ^{3/2}\right)<\theta _H<\pi O\left(ϵ^{3/2}\right)\right)`$. One important conclusion is that the tunneling rate and the tunnel splitting depend on the orientation of the external magnetic field distinctly. Another interesting conclusion concerns the field strength dependence of the WKB exponent or the classical action. We have found that in a wide range of angles, the $`ϵ\left(=1\overline{H}/\overline{H}_c\right)`$ dependence of the WKB exponent or the classical action $`S_{cl}`$ is given by $`ϵ^{5/4}`$, not $`ϵ^{3/2}`$ for $`\theta _H=\pi /2`$, and $`ϵ`$ for $`\theta _H=\pi `$. We have obtained the temperatures corresponding to the crossover from quantum to thermal regime which are found to depend on the orientation of the external magnetic field distinctly. As a result, we conclude that both the orientation and the strength of the external magnetic field are the controllable parameters for the experimental test of the phenomena of macroscopic quantum tunneling and coherence of the magnetization vector in single-domain FM nanoparticles with trigonal and hexagonal symmetries at a temperature well bellow the crossover temperature. We have analyzed the validity of the semiclassical approximation performed in the present work, and have found that the semiclassical approximation should be already rather good for the typical values of parameters for single-domain FM nanoparticles.
Recently, Wernsdorfer and co-workers performed the switching field measurements on individual ferrimagnetic and insulating BaFeCoTiO nanoparticles containing about $`10^5`$-$`10^6`$ spins at very low temperatures (0.1-6K). They found that above 0.4K, the magnetization reversal of these particles is unambiguously described by the Néel-Brown theory of thermal activated rotation of the particle’s moment over a well defined anisotropy energy barrier. Below 0.4K, strong deviations from this model are evidenced which are quantitatively in agreement with the predictions of the MQT theory without dissipation. The BaFeCoTiO nanoparticles have a strong uniaxial magnetocrystalline anisotropy. However, the theoretical results presented here may be useful for checking the general theory in a wide range of systems, with more general symmetries. The experimental procedures on single-domain FM nanoparticles of Barium ferrite with uniaxial symmetry may be applied to the systems with more general symmetries. Note that the inverse of the WKB exponent $`B^1`$ is the magnetic viscosity $`S`$ at the quantum-tunneling-dominated regime $`TT_c`$ studied by magnetic relaxation measurements. Therefore, the quantum tunneling of the magnetization should be checked at any $`\theta _H`$ by magnetic relaxation measurements. Over the past years a lot of experimental and theoretical works were performed on the spin tunneling in molecular Mn<sub>12</sub>-Ac and Fe<sub>8</sub> clusters having a collective spin state $`S=10`$ (in this paper $`S=10^6`$). Further experiments should focus on the level quantization of collective spin states of $`S=10^2`$-$`10^4`$. We hope that the theoretical results presented in this paper may stimulate more experiments whose aim is observing macroscopic quantum phenomena in nanometer-scale single-domain ferromagnets.
## Acknowledgments
R.L. would like to acknowledge Dr. Su-Peng Kou, Professor Zhan Xu, Professor Jiu-Qing Liang and Professor Fu-Cho Pu for stimulating discussions.
## Appendix A: Evaluation of the preexponential factors in WKB tunneling rate
In this appendix, we review briefly the procedure on how to calculate the preexponential factors in the WKB rate of quantum tunneling of the magnetization vector in single-domain FM particles, based on the instanton technique in the spin-coherent-state path-integral representation. The preexponential factors in tunneling rate (MQT) or the tunnel splitting (MQC) are due to the quantum fluctuations about the classical path, which can be evaluated by expanding the Euclidean action to second order in small fluctuations. Then we apply this approach to obtain the instanton’s contribution to the ground-state tunnel splitting for resonant coherently quantum tunneling of the magnetization vector in FM particles with trigonal crystal symmetry in an external magnetic field applied perpendicular to the anisotropy axis (considered in Sec. III) in detail.
In Ref. 12, Garg and Kim have studied the general formulas for evaluating both the WKB exponent and the preexponential factors in the tunneling rate or the tunnel splitting in the single-domain FM particles based on the instanton technique in the spin-coherent-state path-integral representation, without assuming a specific form of the magnetocrystalline anisotropy and the external magnetic field. Here we explain briefly the basic idea of this calculation. Such a calculation consists of two major steps. The first step is to find the classical, or least-action path (instanton) from the classical equations of motion, which gives the exponent or the classical action in the WKB tunneling rate. Instantons in one-dimensional field theory can be viewed as pseudoparticles with trajectories existing in the energy barrier, and are therefore responsible for quantum tunneling. The second step is to expand the Euclidean action to second order in the small fluctuations about the classical path, and then evaluate the Van Vleck determinant of resulting quadratic form. For single-domain FM particles, writing $`\theta \left(\tau \right)=\overline{\theta }\left(\tau \right)+\theta _1\left(\tau \right)`$ and $`\varphi \left(\tau \right)=\overline{\varphi }\left(\tau \right)+\varphi _1\left(\tau \right)`$, where $`\overline{\theta }`$ and $`\overline{\varphi }`$ denote the classical path, one obtains the Euclidean action of Eq. (2) as $`S_E[\theta \left(\tau \right),\varphi \left(\tau \right)]S_{cl}+\delta ^2S`$ with $`S_{cl}`$ being the classical action or the WKB exponent and $`\delta ^2S`$ being a functional of small fluctuations $`\theta _1`$ and $`\varphi _1`$,
$`\delta ^2S`$ $`=`$ $`iS{\displaystyle \frac{d}{d\tau }\left[\mathrm{sin}\overline{\theta }\theta _1\right]\varphi _1𝑑\tau }+{\displaystyle \frac{i}{2}}S{\displaystyle \mathrm{cos}\overline{\theta }\left(\frac{d\overline{\varphi }}{d\tau }\right)\theta _1^2𝑑\tau }`$ ()
$`+{\displaystyle \frac{V_0}{2\mathrm{}}}{\displaystyle \left(E_{\theta \theta }\theta _1^2+2E_{\theta \varphi }\theta _1\varphi _1+E_{\varphi \varphi }\varphi _1^2\right)𝑑\tau },`$
where $`E_{\theta \theta }=\left(^2E/\theta ^2\right)_{\theta =\overline{\theta },\varphi =\overline{\varphi }}`$, $`E_{\theta \varphi }=\left(^2E/\theta \varphi \right)_{\theta =\overline{\theta },\varphi =\overline{\varphi }}`$, and $`E_{\varphi \varphi }=\left(^2E/\varphi ^2\right)_{\theta =\overline{\theta },\varphi =\overline{\varphi }}`$. Under the condition that $`E_{\varphi \varphi }>0`$, the Gaussian integration can be performed over $`\varphi _1`$, and the remaining $`\theta _1`$ path integral can be casted into the standard form for a one-dimensional motion problem. As usual there exists a zero-mode, $`d\overline{\theta }/d\tau `$, corresponding to a translation of the center of the instanton, and a negative eigenvalue in the MQT problem. This leads to the imaginary part of the energy, which corresponds to the quantum escaping rate from the metastable state through the classically impenetrable barrier to a stable one. The resonant tunnel splittings of the ground state for the MQC problem can be evaluated by applying the similar technique. What is need for the calculation of the tunneling rate (in MQT) and the tunnel splitting (in MQC) is the asymptotic relation of the zero mode, $`d\overline{\theta }/d\tau `$, for large $`\tau `$,
$$d\overline{\theta }/d\tau ae^{\mu \zeta },\text{ as }\zeta \mathrm{}.$$
(67)
The new time variable $`\zeta `$ in Eq. (A2) is related to $`\tau `$ as
$$d\zeta =d\tau /2A(\overline{\theta }\left(\tau \right),\overline{\varphi }\left(\tau \right)),$$
(68)
where
$$A(\overline{\theta },\overline{\varphi })=\mathrm{}S^2\mathrm{sin}^2\overline{\theta }/2VE_{\varphi \varphi }.$$
(69)
The partial derivatives are evaluated at the classical path. Then the instanton’s contribution to the tunneling rate for MQT or the tunnel splitting for MQC of the magnetization vector in single-domain FM nanoparticles (without the contribution of the topological Wess-Zumino, or Berry phase term in the Euclidean action) is given by
$$\left|a\right|\left(\mu /\pi \right)^{1/2}e^{S_{cl}}.$$
(70)
Therefore, all that is necessary is to differentiate the classical path (instanton) to obtain $`d\overline{\theta }/d\tau `$, then convert from $`\tau `$ to the new time variable $`\zeta `$ according to Eqs. (A3) and (A4), and read off $`a`$ and $`\mu `$ by comparison with Eq. (A2). If the condition $`E_{\varphi \varphi }>0`$ is not satisfied, one can always perform the Gaussian integration over $`\theta _1`$ and end up with a one-dimensional path integral over $`\varphi _1`$.
Now we apply this approach to the problem of resonant coherently quantum tunneling of the magnetization vector between energetically degenerate easy directions in single-domain FM nanoparticle with trigonal crystal symmetry in an external magnetic field applied perpendicular to the anisotropy axis. After some algebra, we find that
$$E_{\varphi \varphi }2K_1\left(1+12\overline{K}_2ϵ\right),$$
(71)
which is positive. So we can perform the Gaussian integration over $`\varphi _1`$ directly. The relation between $`\tau `$ and the new imaginary-time variable $`\zeta `$ for this MQC problem is found to be
$$\tau =\frac{\mathrm{}S^2}{2K_1V\left(1+12\overline{K}_2ϵ\right)}\zeta .$$
(72)
It is easy to differentiate the instanton solution to obtain
$$\frac{d\overline{\delta }}{d\tau }=8\frac{K_1V}{\mathrm{}S}ϵ\left(1+15\overline{K}_2\frac{ϵ}{2}\right)\mathrm{exp}\left[\sqrt{2ϵ}S\left(1\frac{3}{2}\overline{K}_2+\frac{ϵ}{2}\right)\zeta \right],$$
(73)
as $`\zeta \mathrm{}`$. Thus,
$$\left|a\right|=8\frac{K_1V}{\mathrm{}S}ϵ\left(1+15\overline{K}_2\frac{ϵ}{2}\right),$$
(74)
and
$$\mu =\sqrt{2ϵ}S\left(1\frac{3}{2}\overline{K}_2+\frac{ϵ}{2}\right).$$
(75)
Substituting Eqs. (A9) and (A10) into the general formula (A5), and using Eq. (17) for the classical action or the WKB exponent, we obtain the instanton’s contribution to the tunnel splitting $`\mathrm{}\mathrm{\Delta }_0`$ as expressed in Eq. (21) for nanometer-scale single-domain ferromagnets with trigonal crystal symmetry in the presence of an external magnetic field applied perpendicular to the anisotropy axis.
The calculations of the tunnel splitting and the tunneling rate of the magnetization vector for other MQT and MQC problems considered in the present work can be performed by applying the similar techniques, and we will not discuss them in any further.
Figure Captions:
Fig. 1 The $`\delta \left(=\theta \theta _0\right)`$ dependence of the effective potential $`\overline{E}_1\left(\delta \right)`$ for $`\theta _H=\pi /2`$ (MQC).
Fig. 2 The $`\delta \left(=\theta \theta _0\right)`$ dependence of the effective potential $`\overline{E}_1\left(\delta \right)`$ for $`\theta _H=3\pi /4`$ (MQT). Here, $`\overline{K}_2=0.001`$.
Fig. 3 The $`\theta _H`$ dependence of the relative classical action $`S_{cl}\left(\theta _H\right)/S_{cl}\left(\theta _H=3\pi /4\right)`$ in the trigonal symmetry with $`ϵ=0.001`$ and $`\overline{K}_2=0.001`$ by numerical and analytical calculations.
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# On two simple criteria for recognizing complete intersections in codimension 2
## 1 Introduction
The problem of detecting (global) complete intersections is a key question in projective algebraic geometry and commutative algebra. Up to now, this problem is far from being solved and a complete answer is known only in trivial cases, such that of hypersurfaces in projective spaces or in Grassmanians. Moreover, the relevant conjecture of Hartshorne according to which projective subvarieties of small codimension, compared to their dimension, have to be complete intersections, is still not proved. However, in the last 25 years, there have been some partial results in this direction, particularly in the case of codimension 2. Essentially, the results obtained in the case $`X`$ is a smooth codimension 2 subvariety of $`\text{}^n`$, $`n6`$ can be grouped into two kinds of criteria.
The first one says that if $`X`$ is contained in a hypersurface $`V`$, such that $`deg(V)n2`$, then $`X`$ is a complete intersection (see or the recent improvement in , where it is shown that the bound on the degree of $`V`$ can be increased to $`n1`$, in the case of codimension 2 subvarieties of $`\text{}^6`$); using this kind of criterion one can give also a bound on the degree of $`X`$, so that to assure that $`X`$ is a complete intersection.
The second kind of criterion is based on giving a bound on the number $`p`$ of generators, not for the homogeneous ideal $`I(X)`$, but for an ideal $`I_{sch}(X)`$ which coincide with $`I(X)`$ only in high degree, that is $`[I_{sch}(X)]_d=[I(X)]_d`$, for $`d0`$. We call $`I_{sch}(X)`$ the schematic ideal of $`X`$, in that its generators define $`X`$ scheme-theoretically. Following this approach, Faltings proved in that if $`pn2`$ and $`n8`$ and $`X`$ is a (possibly singular) subcanonical local complete intersection, then it is a complete intersection (in any characteristic). Some years later, this result was improved in , proving in characteristic zero that if $`pn1`$, $`n8`$, then $`X`$ is a complete intersection (but assuming $`X`$ smooth).
The aim of the present work is twofold: on one hand we would like to give a different (in that we use Serre’s correspondence) and simpler proof of the result announced in , hoping to give a crystal clear version of some obscure (to our opinion) arguments. Moreover, assuming only that $`X`$ is a (possibly singular) subcanonical l.c.i., we prove in any characteristic that if $`n3`$ and $`pn1`$, then $`X`$ is a complete intersection and we give some more results working in characteristic zero, assuming that the normal bundle of $`X`$ extends to a numerically split bundle $`E`$ on $`\text{}^n`$ (i.e. the Chern classes of $`E`$ are those of a split bundle), $`n3`$, $`pn`$. On the other hand, as an application of our result we answer to a question posed recently by Franco, Kleiman and Lascu in , (neglecting the case of space curves). Unfortunately, our result shows that the characterization given by Faltings is not peculiar of two codimension embeddings in high dimensional projective spaces.
Our proof is based in constructing and exploiting an exact sequence of locally free sheaves, (sequence (3)), which relates the rank 2 vector bundle $`E`$ appearing in Serre’s correspondence with the generators of the scheme-theoretic ideal of $`X`$.
## 2 Main result: a scheme-theoretic criterion
From now on, $`X`$ will denote a codimension 2 subcanonical l.c.i. (possibly singular) closed subscheme of a projective space $`\mathrm{}^n`$ over an algebraically closed field $`k`$ of any characteristic $`p0`$, where, as usual $`\mathrm{}^n=Proj(k[x_0,\mathrm{},x_n])`$. Then we prove the following result:
Theorem: A) If $`X\mathrm{}^n`$, ($`n3`$) is a scheme-theoretic intersection of $`pn1`$ hypersurfaces, then $`X`$ is a complete intersection.
B) Assume moreover that char($`k`$)=$`0`$, $`n3`$, and $`X`$ scheme-theoretically defined by $`pn`$ hypersurfaces $`V_1,\mathrm{},V_n`$ of degrees $`d_1,\mathrm{},d_p`$, respectively. If the normal bundle of $`X`$ can be extended to a rank 2 vector bundle $`E`$ on $`\text{}^n`$ which is numerically split (i.e. $`c_1(E)=a+b`$ and $`c_2(E)=ab`$, $`a,b\text{}`$) and $`a`$ or $`b`$ is in $`(d_1,\mathrm{},d_p)`$, then $`X`$ is a complete intersection.
Proof of part A: Since $`X`$ is assumed to be subcanonical (i.e. its dualizing sheaf $`\omega _X`$, which is locally free, is of the form $`𝒪_X(e)`$) , by Serre’s correspondence there exists an algebraic vector bundle $`E`$ of rank 2 over $`\mathrm{}^n`$ and a section $`sH^0(\text{},E)`$ such that $`X`$ is identified with the scheme of zeroes of $`s`$, $`Z(s)`$. The Koszul complex for this section gives a projective resolution of the ideal sheaf of $`Z(s)`$, hence of the ideal sheaf of $`X`$:
$$0\stackrel{2}{}E^{}E^{}_X0.$$
(1)
Since $`_X`$ is not itself projective, by (1) it turns out that the projective dimension of $`_X`$ is 1. On the other hand, if $`X`$ is schematically cut out by $`pn1`$ hypersurfaces of degrees $`d_1,\mathrm{},d_p`$, we have an exact sequence:
$$0Ker(f)𝒪(d_i)\stackrel{f}{}_X0.$$
(2)
Since pd($`_X`$)=1, then the first syzygy $`Ker(f)`$ is also projective (see for example ), hence it corresponds to a locally free sheaf. Certainly, we have a morphism $`hHom(𝒪(d_i)𝒪(c_1),_X)`$ which is given first by projecting to $`𝒪(d_i)`$ and then composing with $`f`$, ($`c_1`$ is the first Chern class of $`E`$). Moreover, since $`Ext^1(𝒪(d_i)𝒪(c_1),𝒪(c_1))=0`$, then $`h`$ comes from an element $`g`$ in $`Hom(𝒪(d_i)𝒪(c_1),E^{})`$; indeed, due to the fact that $`^2E^{}𝒪(c_1)`$, we have the following commutative diagram:
$$\begin{array}{ccccccccc}0& & 𝒪(c_1)& & 𝒪(d_i)𝒪(c_1)& & 𝒪(d_i)& & 0\\ & & & & g& & f& & \\ 0& & ^2E^{}& & E^{}& & _X& & 0\end{array}$$
Applying the snake lemma to the previous commutative diagram, we see that $`g`$ is surjectve and that $`Ker(g)Ker(f)`$, so that dualizing the sequence $`0Ker(g)𝒪(d_i)𝒪(c_1)\stackrel{g}{}E^{}0`$, we get a short exact sequence of locally free sheaves:
$$0E\stackrel{\alpha _1}{}𝒪(d_i)𝒪(c_1)\stackrel{\alpha _2}{}C0,$$
(3)
where $`C`$ is just the cokernel sheaf. Since $`C`$ is locally free, it can be identified with a vector bundle of rank equal to $`p1`$ ($`pn1`$). Let $`in_j`$ denote the canonical injection of $`𝒪(d_j)`$ into $`𝒪(d_i)𝒪(c_1)`$, and $`pr_j`$ the corresponding projection from $`𝒪(d_i)𝒪(c_1)`$ to $`𝒪(d_j)`$. Considering the maps $`f_j:=pr_j\alpha _1`$ and $`g_j:=\alpha _2in_j`$ we have a diagram like the following:
$$\begin{array}{ccccccccc}0& & E& \stackrel{\alpha _1}{}& 𝒪(d_i)𝒪(c_1)& \stackrel{\alpha _2}{}& C& & 0\\ & & & \stackrel{f_j}{}& pr_jin_j& \stackrel{g_j}{}& & & \\ & & & & 𝒪(d_j)& & & & \end{array}$$
Now consider the morphisms $`E\stackrel{f_1}{}𝒪(d_1)`$ and $`𝒪(d_1)\stackrel{g_1}{}C`$ and denote by $`Z(f_1)`$ and $`Z(g_1)`$, their respective degeneracy loci. Since in general we have that $`codim(Z(f_j))2`$ and $`codim(Z(g_j))p1`$, it turns out that if $`pn1`$, then $`Z(f_1)Z(g_1)\mathrm{}`$. On the other hand, by exactness of (3), it is clear that $`Z(f_1)Z(g_1)=\mathrm{}`$. Indeed, if it exists $`xZ(f_1)Z(g_1)`$, then $`Ker(f_1)_x=E_x`$, but since the morphism $`\alpha _1`$ can not degenerate at any point (due to the fact that the cokernel $`C`$ is locally free), we have that $`Im(f_1)_x_{i2}𝒪_x(d_i)𝒪_x(C_1)`$. On the other hand, $`Ker(g_1)_x=𝒪_x(d_1)`$, but since $`in_1`$ is always injective, we have that $`Ker(\alpha _2)_x=in_1(𝒪(d_1))_x`$. Hence, if exists $`xZ(f_1)Z(g_1)`$, then $`Ker(\alpha _2)_xIm(\alpha _1)_x=\mathrm{}`$ and so the sequence (3) can not be exact in the middle, at $`x`$. Absurd.
Hence $`Z(f_1)=\mathrm{}`$ or $`Z(g_1)=\mathrm{}`$.
If $`Z(f_1)=\mathrm{}`$, then $`f_1`$ is never degenerate, so dualizing $`E\stackrel{f_1}{}𝒪(d_1)0`$ we get $`0𝒪(d_1)\stackrel{(f_1)^T}{}E^{}`$, where the map $`(f_1)^T`$ is never degenerate; hence $`E^{}`$ splits, so $`E`$ splits and $`X`$ is a complete intersection.
If instead $`Z(g_1)=\mathrm{}`$, we build up the following commutative diagram:
$$\begin{array}{ccccccccc}& & 0& & 𝒪(d_1)& \stackrel{}{}& 𝒪(d_1)& & \mathrm{}\\ & & & & in_1& & g_1& & \\ 0& & E& & 𝒪(d_i)𝒪(c_1)& & C& & 0\\ & & & & & & & & \\ 0& & Ker(\psi )& & _{i2}𝒪(d_i)𝒪(c_1)& \stackrel{\psi }{}& C^{^{}}& & 0\\ & & & & & & & & \\ \mathrm{}& & 0& & 0& & 0& & \end{array}$$
By applying the snake lemma to the two central rows, we see that $`Ker(\psi )E`$, $`C^{^{}}`$ is locally free since $`Z(g_1)=\mathrm{}`$, and we obtain a short exact sequence of locally free sheaves:
$$0E\underset{i2}{}𝒪(d_i)𝒪(c_1)C^{^{}}0.$$
(4)
Repeating the previous reasoning, we can consider the morphisms $`E\stackrel{f_2}{}𝒪(d_2)`$ and $`𝒪(d_2)\stackrel{g_2}{}C^{^{}}`$ and as before $`Z(f_2)=\mathrm{}`$ or $`Z(g_2)=\mathrm{}`$. If $`Z(f_2)=\mathrm{}`$, then $`E`$ splits and $`X`$ is a complete intersection. On the other hand, if $`Z(g_2)=\mathrm{}`$, arguing as before, we obtain a short exact sequence of locally free sheaves:
$$0E\underset{i3}{}𝒪(d_i)𝒪(c_1)C^{^{\prime \prime }}0.$$
In this way we obtain a sequence $`Z(f_1)`$, …, $`Z(f_{p1})`$. If one of these is empty, we are done; otherwise, if all are not empty, then, necessarily, $`Z(g_{p1})=\mathrm{}`$ and as before we obtain:
$$0E𝒪(d_p)𝒪(c_1)0,$$
and we are done.
Proof of part B: To deal with the more restrictive case of characteristic zero, we can assume $`k=\mathrm{}`$ in view of Lefschetz’s principle. In our assumption $`X`$ is cut out schematically by $`n`$ hypersurfaces $`V_1,\mathrm{},V_n`$ of degrees $`d_1,\mathrm{},d_n`$ and we have that $`c_1(E):=a+b=d_k+b`$ and $`c_2(E):=ab=d_kb`$ for some $`k(1,\mathrm{},n)`$. (It is not restrictive to assume that $`a(d_1,\mathrm{},d_n)`$). Reordering the hypersurfaces, we can assume that $`c_1(E)=d_1+b`$ and $`c_2(E)=d_1b`$. From the exact sequence (3), it is clear that the rank of $`C`$ is $`n1`$, so that the morphism $`𝒪(d_1)\stackrel{g_1}{}C`$ degenerates at most in codimension $`n1`$. On the other hand, the morphism $`E\stackrel{f_1}{}𝒪(d_1)`$ can not degenerate in codimension $`2`$, otherwise, if $`Z(f_1)`$ is its degeneracy locus, the Poincaré dual $`[Z(f_1)]H^4(\text{},\mathrm{})`$ would represent $`c_2(E^{}𝒪(d_1))=c_1(E^{})d_1+c_2(E^{})+d_1^2=0`$, so that either the morphism $`f_1`$ does not degenerate at all, and in this case we are done as before, or it degenerates in codimension $`1`$.
So if the morphism $`f_1`$ degenerates in codimension one, we have that $`Z(f_1)Z(g_1)\mathrm{}`$, provided that $`Z(g_1)\mathrm{}`$. On the other hand, by exactness of (3), we must have $`Z(f_1)Z(g_1)=\mathrm{}`$, so that we conclude that $`Z(g_1)=\mathrm{}`$. Thus, arguing as in part A, we obtain the following short exact sequence:
$$0E\underset{i2}{\overset{n}{}}𝒪(d_i)𝒪(c_1)C^{^{}}0,$$
and, from this, we conclude as in part A.
The result of part B of the previous theorem can be interpreted as a relation between degree $`deg(X)`$ and subcanonicity $`e`$, recalling the well-known fact that if $`E`$ is the vector bundle associated to $`X`$ via Serre’s correspondence, then $`deg(X)=c_2(E)`$, while $`e+n+1=c_1(E)`$.
Corollary A: Let $`X\text{}^n`$ ($`X`$ as in the hypotheses of Theorem B), $`n3`$ be scheme-theoretically defined by $`n`$ hypersurfaces of degrees $`d_1,\mathrm{},d_n`$, and let $`l`$ be an integer in the set $`(d_1,\mathrm{},d_n)`$. If the following relation is satisfied:
$$deg(X)+l^2(e+n+1)l=0,$$
(5)
then $`X`$ is a complete intersection.
Proof: Arguing as in part B, it is clear that to show that $`E`$ splits is sufficient to show that $`c_2(E^{}𝒪(l))=0`$ for some $`l`$ as above. But the vanishing of the second Chern class of $`E^{}𝒪(l)`$ is exactly the relation (5), as an easy computation can show.
Remark 1: The approach of giving bounds on the degree of a subvariety to detect a complete intersection is particularly ”effective”, but it is obviously hopeless if one pretend to solve Hartshorne’s conjecture. On the other hand, since any closed subscheme (irreducible or not) of $`\text{}^n`$ which is a local complete intersection is always scheme-theoretically defined by $`n+1`$ hypersurfaces, as proved in , the approach of recognizing a complete intersection via the number of generators of its scheme-theoretic ideal, could be in principle useful to solve the conjecture. Unfortunately, the cases $`p=n`$ and in particular $`p=n+1`$ (the generic case) appear completely intractable, at least up to now, since it is very difficult to relate the algebro-geometric properties of a small codimension embedding, with those of its scheme-theoretic ideal.
Remark 2: In the light of the previous remark and of Theorem B, it would be nice to know when a subvariety can be scheme-theoretically defined by $`n`$ equations. To get a sufficient condition, we can use the theory of excess and residual intersections as developed in . For example, let us consider $`4`$ hypersurfaces $`\{V_1,\mathrm{},V_4\}`$ in $`\text{}^4`$, such that $`V_i=X\{p_1,\mathrm{},p_k\}`$, where $`X`$ is a smooth subcanonical surface and $`\{p_1,\mathrm{},p_k\}`$ are (possibly non reduced) points, i.e. the four hypersurfaces define scheme-theoretically the union of $`X`$ and a bunch of points outside $`X`$. The theory of residual intersections enables us to predict the (weighted) number of residual points as a function of the degrees of the hypersurfaces, of the degree of $`X`$ and of the degrees of the Chern classes of $`T_X`$, the tangent sheaf of $`X`$. Imposing that the number of residual points is zero, gives us a sufficient condition for a surface (a subvariety in general) to be scheme-theoretically defined by $`n`$ equations. Thus combining Proposition 9.12 (page 154) of with Example 9.1.5. we get (for a surface in $`\text{}^4`$):
$$deg(c_2(T_X))+(\sigma _1(g_i)5)deg(c_1(T_X))+$$
$$+(\sigma _2(g_i)5\sigma _1(g_i)+15)deg(X)+W(p_1,\mathrm{},p_k)=\sigma _4(g_i),$$
(6)
where $`W(p_1,\mathrm{},p_k)`$ is the weighted number of residual points and $`\sigma _j(g_i)`$ is the j-th elementary symmetric polynomial in the degrees $`g_i`$ of the hypersurfaces $`V_i`$. Assuming $`X`$ subcanonical, from $`c_1(T_X)=K_X`$, we get $`c_1(T_X)=eH`$, where $`H`$ is the class of a hyperplane section; moreover, from the exact sequence:
$$0T_XT_\text{}^4𝒪_XN_{X/\text{}^4}0,$$
we get $`c_2(T_X)=10H^2+5HK+K^2c_2(N_{X/\text{}^4})`$, and by $`c_1(T_X)=eH`$, we have $`c_2(T_X)=(105e+e^2)H^2c_2(N_{X/\text{}^4})`$. Since $`deg(c_2(N_{X/\text{}^4}))=d^2`$, $`deg(H^2)=d`$ and $`deg(H)=d`$, (where $`d`$ is $`deg(X)`$), substituting in (6), taking degrees and imposing $`W(p_1,\mathrm{},p_k)=0`$ reads:
$$[25+\sigma _2(g_i)(5+e)\sigma _1(g_i)+e^2d]d=\sigma _4(g_i).$$
(7)
Thus, the relation (7) gives a sufficient condition for a subcanonical surface in $`\text{}^4`$ to be scheme-theoretically defined by $`4`$ equations.
### 2.1 An application: the linkage criterion
As usual, if $`X`$ and $`Y`$ are l.c.i. of codimension $`2`$ in $`\text{}^n`$, we say that $`X`$ is (directly) linked to $`Y`$ if there exists a complete intersection $`(F_1,F_2)`$, such that $`Y`$ is the residual scheme of $`X`$ in the intersection $`F_1F_2`$, and viceversa. In , working in characteristic zero and assuming $`X`$ smooth and subcanonical $`dimX1`$, Beorchia and Ellia proved that $`X`$ is a complete intersection iff it is self-linked, i.e. iff there exists complete intersection $`(F_1,F_2)`$ such that $`F_1F_2=2X`$ ($`F_1`$ and $`F_2`$ define on $`X`$ a double structure which is a complete intersection). They also asked if the same criterion holds also for possibly singular l.c.i.. Recently, in , Franco, Kleiman and Lascu have given a positive answer to this question proving that the same criterion holds avoiding smoothness: $`X`$ can be reducible and non reduced. Their proof works only in characteristic zero (unless $`dimX4`$, where it holds over any algebraically closed field, due to a previous result of Faltings), so they ask if the same holds in positive characteristic, for lower dimensional $`X`$. Using our Theorem A we prove the following:
Proposition A: Let $`X`$ be a subcanonical (possibly singular) l.c.i. subscheme of codimension 2 in $`\text{}^n`$, $`n4`$, defined over an algebraically closed field of any characteristic. Then $`X`$ is a complete intersection iff it is self-linked.
Proof: According to the Gherardelli linkage theorem , which holds over any algebraically closed field (see for its proof) we know that $`XF_1F_2`$ is subcanonical iff its residual scheme $`Y`$ (in the complete intersection $`F_1F_2`$) is scheme-theoretically defined by the intersection of $`F_1`$ and $`F_2`$ with a third hypersurface $`F_3`$. On the other hand, if $`X`$ is self-linked by $`F_1`$ and $`F_2`$, then, by definition $`X`$ is equal to its own residual scheme in the complete intersection of $`F_1`$ and $`F_2`$, and since $`X`$ is assumed subcanonical, by the Gherardelli theorem it is scheme theoretically defined by $`F_1,F_2`$ and $`F_3`$; hence, by Theorem A, it is a complete intersection as soon as $`dimX2`$. Viceversa, if $`X`$ is a complete intersection, it is immediate to see that it is self-linked (just consider the intersection of $`F_1`$ and $`2F_2`$ if $`X=F_1F_2`$).
There is an immediate generalization of the previous proposition, which is the following:
Proposition B: Let $`X`$ as in Proposition A. Then $`X`$ is a complete intersection iff it can be (directly) linked to $`Y`$, where $`Y`$ is any subcanonical (possibly singular) l.c.i. subscheme.
Proof: It is sufficient to use again the Gherardelli linkage and Theorem A.
Remark 3: The $`X`$’s as in Proposition B are self-linked iff they are scheme-theoretically defined by three hypersurfaces. Indeed, if $`X`$ is self-linked, then by Gherardelli it is schematically defined by $`3`$ equations; viceversa, if $`X`$ is defined by $`3`$ equations it is a complete intersection by Theorem A and then it is self-linked by Proposition B.
Remark 4: The most difficult case, in order to characterize complete intersection via self-linking is that of curves in $`\text{}^3`$. Beorchia and Ellia proved their criterion (in characteristic zero) also for curves, assuming that they are smooth, while Franco, Kleiman and Lascu extended this result to l.c.i. curves (always working in characteristic zero). In positive characteristic (characteristic 2), however, there is certainly a counterexample for this criterion to hold in the case of curves, due to Migliore (see the discussion at the end of ). So our extension of this criterion over a field of any characteristic is the best possible for low dimensional subvarieties: that is surfaces are the lowest dimensional subvarieties where this criterion holds without exceptions. Unfortunately, up to now, there is no positive result for the case of space curves in characteristic grater than zero.
Acknowledgements: It is a pleasure to thank Philippe Ellia and Alexandru Lascu for useful discussions, and Steve Kleiman for valuable remarks and corrections; I would like also to thank the members of the Department of Mathematics of Ferrara University for the kind hospitality and the warm and stimulating atmosphere.
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# 1 Rule table for the DK PCA. The first line gives the initial neighbourhood, the other two lines give the probability at which the state listed at left is reached by the central bit.
The Domany-Kinzel (DK) probabilistic cellular automaton (PCA) is one of the most studied PCA in the physics literature, because it is the most general left-right symmetric one-dimensional PCA, and has the interesting property of having the mixed site-bond directed percolation (DP) process on the square lattice as one of its instances. Its defining rules are given in table 1. The mixed site-bond DP process is given by assigning a probability $`s[\mathrm{\hspace{0.17em}0},1]`$ for a site to be present in the lattice, and a probability $`b[\mathrm{\hspace{0.17em}0},1]`$ for a bond to exist between any two sites on the lattice. The mixed DP problem consists in finding the values of $`s`$ and $`b`$ for which an infinite cluster of sites connected by the bonds occurs such that one can walk unidirectionally in it, say to the south and to the east indefinitely. In the DK PCA this problem is obtained by choosing the transition probabilities $`x=0`$, $`y=sb`$, and $`z=sb(2b)`$. For $`s=1`$ one obtains the pure bond DP problem, whilst for $`b=1`$ one obtains the pure site DP problem.
In this paper we investigate numerically the critical behaviour of a continuous-time one-dimensional non-attractive lattice gas for which some lower bounds on the critical points of the PCA version was given recently . The model is related with the DK PCA in one of its parameter subspaces, and although the model is not on the mixed site-bond DP parameter subspace of the DK PCA, it presents DP exponents, as expected on the basis of the DP conjecture .
Let $`n_{\mathrm{}}(t)\{0,1\}`$ denote the occupation number of the site $`\mathrm{}\mathrm{\Lambda }`$ at the integer instant $`t`$, with $`\mathrm{\Lambda }`$ a finite lattice with $`|\mathrm{\Lambda }|=L`$ sites and periodic boundary conditions $`\mathrm{}+L\mathrm{}`$. The model we are interested in is the continuous-time version of the PCA defined by the rules
$$n_{\mathrm{}}(t+1)=\{\begin{array}{cc}(n_\mathrm{}1(t)+n_{\mathrm{}+1}(t))mod2& \text{with probability }p,\hfill \\ 0& \text{with probability }1p.\hfill \end{array}$$
(1)
The rule table for this PCA is given in table 2. From tables 1 and 2 we see that our PCA is equivalent to the DK PCA with rates $`x=0`$, $`y=p`$, and $`z=0`$. We thus see that unless we take the unphysical value $`b=2`$ in the site-bond DP subspace of the DK PCA, this model does not belong to that subspace.
Our approach in constructing the continuous-time version for the above PCA is to take its non-diagonal transitions, i.e., those transitions for which the final state differs from the initial state, and associate with them a stochastic lattice gas with transition rates given by the original PCA rules. This approach has been used before in the PCA literature , and is equivalent to the so-called ‘Hamiltonian’ or ‘strong anisotropic’ limit for the transfer matrixes of equilibrium lattice models .
As is well known , we may write the master equation for interacting lattice gases as a Schrödinger-like equation in Euclidean time,
$$\frac{\mathrm{d}}{\mathrm{d}t}|P(t)=H|P(t),$$
(2)
with $`|P(t)`$ the generating vector of the probabilities $`P(𝐧,t)=𝐧|P(t)`$ of observing the configuration $`𝐧=(n_1,n_2,\mathrm{},n_L)\{0,1\}^\mathrm{\Lambda }`$ at instant $`t`$, and with the infinitesimal generator $`H`$ of the Markov semigroup playing the role of the Hamiltonian. For the non-diagonal transitions of the PCA defined by Eq. (1), the operator $`H`$ can be written as
$$H=\underset{\mathrm{}=1}{\overset{L}{}}H_{\mathrm{}1,\mathrm{},\mathrm{}+1},$$
(3)
with the three-body stochastic transition matrix $`H_{\mathrm{}1,\mathrm{},\mathrm{}+1}`$ given by
$$H_{\mathrm{}1,\mathrm{},\mathrm{}+1}=\left(\begin{array}{cccccccc}& & 1& & & & & \\ & p& & 1p& & & & \\ & & 1& & & & & \\ & p& & 1+p& & & & \\ & & & & p& & 1p& \\ & & & & & & & 1\\ & & & & p& & 1+p& \\ & & & & & & & 1\end{array}\right),$$
(4)
where the three-site basis vectors are ordered as usual, $`(0,0,0)(0,0,1)\mathrm{}(1,1,0)(1,1,1)`$, and the dots indicate null entries. Proper tensorization of the above three-body matrix with unit matrixes in order to obtain the full matrix $`H`$ is understood. Notice that the elements in the columns of $`H_{\mathrm{}1,\mathrm{},\mathrm{}+1}`$ (and consequently of $`H`$) add to zero due to the conservation of probabilities, and that its non-diagonal elements are positive, since $`0p1`$. Identifying a particle with the up spin state and a hole with the down spin state in the $`\sigma ^z`$ basis, the transition matrix $`H_{\mathrm{}1,\mathrm{},\mathrm{}+1}`$ above is seen to be equivalent to the non-Hermitian quantum spin operator
$$H_{\mathrm{}1,\mathrm{},\mathrm{}+1}=\frac{1}{2}(\sigma _{\mathrm{}}^x1)\left[1+(1p)\sigma _{\mathrm{}}^z+p\sigma _\mathrm{}1^z\sigma _{\mathrm{}}^z\sigma _{\mathrm{}+1}^z\right],$$
(5)
where $`\sigma ^x`$ and $`\sigma ^z`$ are the usual Pauli spin-$`\frac{1}{2}`$ matrices. The transition matrix (4) resembles the analogous matrix for the basic contact process, but with non-standard rates and with the elementary process $`101111`$ lacking. This lack is the root of the non-attractiveness of the process. (Loosely speaking, attractive interacting particle systems present a tendency for clustering, as it occurs in ferromagnetic models or in the basic contact process. The precise mathematical statement of attractiveness can be found in .)
The lowest gap in the spectrum of $`H`$ may be used to perform a finite-size scaling analysis in the same way as one does in equilibrium problems . Around the critical point $`pp^{}`$, the correlation lengths of the infinite system behave like
$$\xi _{}\xi _{}^z(pp^{})^\nu _{}(pp^{})^{\nu _{}z},$$
(6)
where $`\xi _{}`$ and $`\xi _{}`$ are the correlation lengths respectively in the time and space directions, $`\nu _{}`$ and $`\nu _{}`$ are the corresponding critical exponents, and $`z=\nu _{}/\nu _{}`$ is the dynamical critical exponent. For finite systems of size $`L`$ we expect that
$$\xi _{,L}^1=L^{z_L}\mathrm{\Phi }\left(|pp_L^{}|L^{1/\nu _{,L}}\right),$$
(7)
where $`z_L`$ and $`\nu _{,L}`$ are the finite versions of $`z`$ and $`\nu _{}`$, and $`\mathrm{\Phi }(u)`$ is a scaling function with $`\mathrm{\Phi }(u1)u^\nu _{}`$. On general grounds one expects $`lim_L\mathrm{}p_L^{},z_L,\nu _{,L}=p^{},z,\nu _{}`$. From Eqs. (6) and (7) we obtain
$$\frac{\mathrm{ln}\left[\xi _{,L}(p_L^{})/\xi _{,L^{}}(p_L^{})\right]}{\mathrm{ln}(L/L^{})}=\frac{\mathrm{ln}\left[\xi _{,L^{\prime \prime }}(p_L^{})/\xi _{,L}(p_L^{})\right]}{\mathrm{ln}(L^{\prime \prime }/L)}=z_L,$$
(8)
which through the comparison of three different system sizes $`L^{}<L<L^{\prime \prime }`$ furnishes simultaneously $`p_L^{}`$ and $`z_L`$. Of course, $`\xi _{,L}`$ and the gap $`E_L^{(1)}E_L^{(0)}=E_L^{(1)}`$ of $`H`$ are related by $`\xi _{,L}^1=\mathrm{Re}\{E_L^{(1)}\}`$.
We calculated the gaps of $`H`$ with the power method, which requires only matrix-by-vector multiplications that can be carried out efficiently and does not require a diagonalization in the usual, ‘QR’ sense, a step that may lessen the quality of the data. The version of the power method we use takes advantage of the presence of absorbing states, and is also suitable for the investigation of time dependent properties of Markov chains .
Our results for $`p^{}`$ and $`z`$ are summarized in figure 1. Curiously enough, despite the translational invariance of the lattice gas rules the finite-size estimates behaved better for triplets of lengths of the form $`L^{},L,L^{\prime \prime }=2l1,2l,2l+1`$, $`l`$. Both sets of data behaved irregularly with the system sizes, preventing us from applying the usual extrapolation algorithms to them. We are presumably in the presence of strong finite-size effects and corrections to scaling. Notice that the last three points of the data seem to be converging monotonically, but unfortunately we were not able to go beyond $`L=22`$ in our diagonalizations. The $`L=\mathrm{}`$ values for $`p_L^{}`$ were obtained through a least-squares fit to the curve $`x_L=x_{\mathrm{}}+aL^1`$, since the data scale well with $`L^1`$ (although it does so with $`L^2`$ also, but with a smaller correlation coefficient), whilst for $`z_L`$ we only estimated the mean value of our data. As expected on the basis of the DP conjecture, namely, that the phase transition about a single absorbing state in single-component systems with a scalar order parameter and in the absence of internal symmetries should be in the DP universality class of critical behaviour , our lattice gas shows a DP-compatible exponent $`z=1.58\pm 0.04`$. The most precise value of $`z`$ for the DP universality class to date is given by $`z=\mathrm{1.580\hspace{0.17em}745}\pm \mathrm{0.000\hspace{0.17em}010}`$ . The critical point of the model is estimated as $`p^{}=0.926\pm 0.004`$ (LS correlation coefficient $`\gamma =0.944`$). This value is slightly higher than the critical value $`p_{\mathrm{DK}}^{}=0.82\pm 0.01`$ of the corresponding point ($`x=0`$, $`y=p`$, $`z=0`$) in the DK PCA . This shift in the critical point for the lattice gas version of the PCA was observed before in a study similar to the present one, where the properties of the lattice gas on the line $`x=0`$, $`y=z`$ (corresponding to the pure site DP problem) was investigated , and is probably a general feature, since asynchronous dynamics tend to be more noisy than synchronous dynamics.
In summary, we conducted numerical diagonalizations of the infinitesimal generator of the continuous-time version of a non-attractive probabilistic cellular automaton (PCA) that is an instance of the Domany-Kinzel (DK) PCA. Although our PCA is not in the site-bond directed percolation (DP) subspace of the DK PCA, its continuous-time version shows DP critical behaviour. Our finite-size data showed an irregular approach to the infinite system limit, and this prevented us from obtaining good estimates of the critical values. It would be interesting to study this lattice gas by time-dependent Monte Carlo methods in order to obtain more accurate critical values.
This work was supported by the Fundação de Amparo à Pesquisa do Estado de São Paulo (FAPESP), Brazil.
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# Pion photo- and electroproduction and the partially-conserved axial current
## Abstract
The relevance of the axial current for pion production processes off the nucleon with real or virtual photons is revisited. Employing the hypothesis of a partially-conserved axial current (PCAC), it is shown that, when all of the relevant contributions are taken into account, PCAC does not provide any additional constraint for threshold production processes that goes beyond the Goldberger–Treiman relation. In particular, it is shown that pion electroproduction processes at threshold cannot be used to extract any information regarding the weak axial form factor. The relationships found in previous investigations are seen to be an accident of the approximations usually made in this context.
The hypothesis of a partially-conserved axial current (PCAC) has been employed in many investigations for constraining scattering processes involving pions at threshold. One of its early successes was the relation by Goldberger and Treiman between the strength of the weak decay of the nucleon $`g_\mathrm{a}`$ and the strong-interaction $`\pi NN`$ coupling constant $`g_{\pi NN}`$, i.e.,
$$\frac{g_\mathrm{a}}{f_\pi }\frac{g_{\pi NN}}{m},$$
(1)
where $`f_\pi `$ is the weak decay constant of the pion and $`m`$ is the nucleon mass. Experimentally this relation is found to be satisfied to better than 10%.
Recently, PCAC relations were employed to extract the properties of the nucleon’s weak decay axial form factor $`G_\mathrm{a}`$ from threshold pion electroproduction data (see , and references therein). These extractions are based on the assumption that $`G_\mathrm{a}`$ is related to the electromagnetic structure of the Kroll–Ruderman contact term .
I will show here that the previous derivations of this relationship are based on an incomplete evaluation of the relevant PCAC expressions and that, if all mechanisms are taken into account, the dependence on $`G_\mathrm{a}`$ vanishes. Thus, it will become obvious that the identification of $`G_\mathrm{a}`$ as the form factor entering the Kroll–Ruderman term is an accident of the usually employed approximations.
To set the stage, the basic PCAC relations will be recapitulated first. Excluding ‘second class’ (i.e., tensor) currents, the general form of the weak axial current is given by
$`j_\mathrm{a}^\mu `$ $`=`$ $`\overline{u}_f\gamma _5\left[\gamma ^\mu G_\mathrm{a}+(pp^{})^\mu G_\mathrm{p}\right]{\displaystyle \frac{\tau }{2}}u_i.`$ (2)
The PCAC hypothesis constrains this current by
$`(p^{}p)_\mu j_\mathrm{a}^\mu `$ $`=`$ $`{\displaystyle \frac{f_\pi \mu ^2}{t\mu ^2}}\overline{u}_f\gamma _5G_tu_i\tau ,`$ (3)
which provides a conserved current for vanishing pion mass $`\mu `$. $`G_t`$ is the $`\pi NN`$ vertex function; other than the $`\gamma _5`$ which has been pulled out explicitly, I make no assumptions about the internal structure of $`G_t`$. Of course, within the present context, i.e., between on-shell spinors, $`G_t`$ is a function of $`t=(pp^{})^2`$ only, where $`p`$ and $`p^{}`$ are the initial and final nucleon momenta, respectively. $`\tau `$ is the vertex isospin operator. Here and throughout the present work, a Cartesian isospin basis is being employed, and the corresponding indices are suppressed; summation over these indices is implied when quantities carrying isospin indices are multiplied with each other.
The two form factors of the axial current are related via Eq. (3), i.e.,
$$2mG_\mathrm{a}+tG_\mathrm{p}=2f_\pi \frac{\mu ^2}{t\mu ^2}G_t.$$
(4)
Evaluated at $`t=0`$, this provides the Goldberger–Treiman relation ,
$$\frac{G_\mathrm{a}(0)}{f_\pi }=\frac{G_t(0)}{m}.$$
(5)
Equation (1) assumes that $`G_\mathrm{a}(0)G_\mathrm{a}(\mu ^2)g_\mathrm{a}`$ and $`G_t(0)G_t(\mu ^2)g_{\pi NN}`$ since the pion mass is small.
While strictly speaking, Eq. (3) is presumed to be valid only for $`t`$ values up to order $`\mu ^2`$ , I will in the following take all of the preceding relations at face value, assuming them to be valid at the operator level, and will consider limits of small $`t`$, etc., only at the end.
Introducing an operator $`\widehat{ȷ}_\mathrm{a}^\mu `$ for the axial current, i.e.,
$$j_\mathrm{a}^\mu =\overline{u}_f\widehat{ȷ}_\mathrm{a}^\mu u_i,$$
(6)
it can be split into weak and hadronic parts according to
$$\widehat{ȷ}_\mathrm{a}^\mu =\widehat{ȷ}_{\mathrm{a},\mathrm{w}}^\mu +\widehat{ȷ}_{\mathrm{a},\mathrm{h}}^\mu ,$$
(7)
where, having eliminated $`G_\mathrm{p}`$ with the help of Eq. (4),
$`\widehat{ȷ}_{\mathrm{a},\mathrm{w}}^\mu `$ $`=`$ $`\gamma _5\left[\gamma ^\mu +(p^{}p)^\mu {\displaystyle \frac{2m}{t}}\right]G_\mathrm{a}{\displaystyle \frac{\tau }{2}},`$ (9)
$`\widehat{ȷ}_{\mathrm{a},\mathrm{h}}^\mu `$ $`=`$ $`f_\pi (p^{}p)^\mu {\displaystyle \frac{\mu ^2}{t}}{\displaystyle \frac{1}{t\mu ^2}}\gamma _5G_t\tau `$ (10)
separate the dependence on the weak and hadronic form factors, respectively. The divergence of the weak part,
$`(p^{}p)_\mu \widehat{ȷ}_{\mathrm{a},\mathrm{w}}^\mu `$ $`=`$ $`[\gamma _5(p/m)+(p/^{}m)\gamma _5]G_\mathrm{a}{\displaystyle \frac{\tau }{2}},`$ (11)
which vanishes between nucleon spinors, signifies the conserved part of the current and
$`(p^{}p)_\mu \widehat{ȷ}_{\mathrm{a},\mathrm{h}}^\mu `$ $`=`$ $`f_\pi {\displaystyle \frac{\mu ^2}{t\mu ^2}}\gamma _5G_t\tau `$ (12)
provides the PCAC divergence of Eq. (3).
Note that the two contributions $`\widehat{ȷ}_{\mathrm{a},\mathrm{w}}^\mu `$ and $`\widehat{ȷ}_{\mathrm{a},\mathrm{h}}^\mu `$ may be interpreted as resulting from the two diagrams of Fig. 1. The hadronic current $`\widehat{ȷ}_{\mathrm{a},\mathrm{h}}^\mu `$, in particular, provides the straightforward interpretation of the pion-pole-dominated diagram of Fig. 1: It describes the creation of the pion of mass $`\mu `$, with coupling operator $`f_\pi (p^{}p)^\mu `$ and associated normalized ‘form factor’ $`\mu ^2/t`$, and the subsequent propagation of the pion and its final absorption in the nucleon. In other words,
$`\widehat{ȷ}_\pi ^\mu `$ $`=`$ $`f_\pi \widehat{q}^\mu {\displaystyle \frac{\mu ^2}{\widehat{q}^2}}`$ (13)
corresponds to the circle labeled H in Fig. 1, with $`\widehat{q}=p^{}p`$ being the pion’s four-momentum flowing out of H.
The preceding operator-level description of the axial current, and its diagrammatic interpretation, will provide precise meaning for the following of how the photon couples to the axial current.
Turning now to the main issue of the present work, i.e., the production of pions off the nucleon with real or virtual photons, the corresponding amplitude $``$ is determined by the four diagrams in Fig. 2 , i.e.,
$$=\overline{u}_f\left(M_s^\nu +M_u^\nu +M_t^\nu +M_{\mathrm{int}}^\nu \right)u_i\epsilon _\nu .$$
(14)
Adapting the PCAC hypothesis to pion photoproduction by employing minimal substitution, Adler finds that $``$ satisfies
$`{\displaystyle \frac{f_\pi \mu ^2}{q^2\mu ^2}}`$ $`=`$ $`q_\mu J_{\mathrm{a},\gamma }^{\mu \nu }\epsilon _\nu Q_\pi j_\mathrm{a}^\nu \epsilon _\nu ,`$ (15)
where $`J_{\mathrm{a},\gamma }^{\mu \nu }`$ describes the coupling of the photon to the axial current and $`j_\mathrm{a}^\nu `$ is the nucleon matrix element (6) of the axial current. $`(Q_\pi )_{kl}=ei\epsilon _{k3l}`$ is the pion charge operator. Note that only the nucleons are on-shell here, but the pion is off-shell.
This relation between the pion photoproduction amplitude $``$ and the axial current is presumed to be valid only in the limit of vanishing pion momentum $`q`$. In the soft-pion limit $`q0`$, following Ref. , the first term on the right-hand side here is often taken as zero (see also ). If true, this immediately provides
$`|_{q=0}=e_\pi \overline{u}_f{\displaystyle \frac{\gamma _5\gamma ^\nu }{2m}}\stackrel{~}{G}_\mathrm{a}(k^2)u_i\epsilon _\nu ,`$where $`e_\pi =Q_\pi \tau `$ effectively describes the charge of the outgoing pion. The operator structure of this expression, $`\gamma _5\gamma ^\nu `$, is identical to the Kroll–Ruderman contact current , with a form factor $`\stackrel{~}{G}_\mathrm{a}(k^2)=G_t(0)G_\mathrm{a}(k^2)/G_\mathrm{a}(0)`$ that derives its normalization from the $`\pi NN`$ form factor $`G_t`$ but its functional behavior from the axial form factor $`G_\text{a}`$; $`k`$ here is the incoming photon momentum. This is taken as evidence that the electromagnetic structure of the Kroll–Ruderman term must be described by the axial form factor $`\stackrel{~}{G}_\mathrm{a}`$ and that $`Q_\pi j_\mathrm{a}^\nu \epsilon _\nu `$ may be used as the starting point for extracting the threshold behavior of pion production processes by considering expansions around $`q=0`$ .
In the following, I will show that these results follow from an incomplete treatment of the prevailing dynamical situation and that none of these conclusions is warranted. In doing so, it will become clear that Eq. (15)—which is based on the applicability of the minimal substitution procedure to the present case—may need to be modified to correctly describe the fact that the hadrons have internal structure.
To this end, I will consider the divergence of the current $`J_{\mathrm{a},\gamma }^{\mu \nu }\epsilon _\nu `$ of Eq. (15). Instead of evaluating this in the usual manner by the LSZ reduction scheme , it is much more convenient to do this in terms of Feynman diagrams, consistent with the operator approach adopted here for the axial currents.
The current $`J_{\mathrm{a},\gamma }^{\mu \nu }`$ corresponds to inserting photon lines in all possible places in the axial-current diagrams of Fig. 1. The result is shown in Fig. 3; it can be verified either by direct inspection of all relevant graphs or, in a more formal way, by using the gauge-derivative method of Ref. (which is completely equivalent to the usual expansion of the relevant Green’s functions in terms of the electromagnetic field $`A^\nu `$ and summing up all first-order contributions in $`A^\nu `$). In terms of operators, the resulting expression is
$`\widehat{J}_{\mathrm{a},\gamma }^{\mu \nu }`$ $`=`$ $`\widehat{ȷ}_{\mathrm{a},f}^\mu {\displaystyle \frac{1}{p/+k/m}}\mathrm{\Gamma }_i^\nu +\mathrm{\Gamma }_f^\nu {\displaystyle \frac{1}{p/^{}k/m}}\widehat{ȷ}_{\mathrm{a},i}^\mu `$ (18)
$`+\widehat{ȷ}_\pi ^\mu {\displaystyle \frac{1}{q^2\mu ^2}}\mathrm{\Gamma }_\pi ^\nu {\displaystyle \frac{1}{t\mu ^2}}\gamma _5G_t\tau +W^{\mu \nu }`$
$`+H^{\mu \nu }{\displaystyle \frac{1}{t\mu ^2}}\gamma _5G_t\tau +\widehat{ȷ}_\pi ^\mu {\displaystyle \frac{1}{q^2\mu ^2}}M_{\mathrm{int}}^\nu .`$
The operators $`W^{\mu \nu }`$ and $`H^{\mu \nu }`$ describe the respective contact terms from the second line of Fig. 3. $`H^{\mu \nu }`$ is given by
$`H^{\mu \nu }`$ $``$ $`\left\{\widehat{ȷ}_\pi ^\mu (p^{}p)\right\}^\nu `$ (19)
$`=`$ $`f_\pi \mu ^2\left[g^{\mu \nu }{\displaystyle \frac{q^\mu (2qk)^\nu }{q^2}}\right]{\displaystyle \frac{Q_\pi }{t}},`$ (20)
where $`\{j_\pi ^\mu \}^\nu `$ is the gauge-derivative notation of Ref. which describes the coupling of the photon to $`\widehat{ȷ}_\pi ^\mu `$ of Eq. (13). Note that this is identical to what one obtains from minimal substitution, which is appropriate here since $`\widehat{ȷ}_\pi ^\mu `$ does not contain any unknown functions. By contrast, for $`W^{\mu \nu }`$, defined analogously as
$$W^{\mu \nu }\left\{\widehat{ȷ}_{\mathrm{a},\mathrm{w}}^\mu (p^{}p)\right\}^\nu ,$$
(21)
one cannot give a result in closed form in general since the internal structure of $`G_\mathrm{a}`$ is unknown. However, for the present discussion it suffices to know that $`W^{\mu \nu }`$ only depends on $`G_\mathrm{a}`$ since this is the only form factor contained in $`\widehat{ȷ}_{\mathrm{a},\mathrm{w}}^\mu `$. \[For a discussion of the structureless limit $`G_\mathrm{a}g_\mathrm{a}`$, see remarks pertaining to Eq. (28).\]
In Eq. (18), $`\mathrm{\Gamma }_\pi ^\nu `$ is the electromagnetic current for the pion and the (gauge-invariant) current for the nucleon is
$`\mathrm{\Gamma }_\mathrm{n}^\nu `$ $`=`$ $`\gamma ^\nu Q_\mathrm{n}+T_\mathrm{n}^\nu ,`$ (23)
$`T_\mathrm{n}^\nu `$ $`=`$ $`(\gamma ^\nu k^2k^\nu k/){\displaystyle \frac{F_11}{k^2}}Q_\mathrm{n}+i{\displaystyle \frac{\sigma ^{\nu \lambda }k_\lambda }{2m}}\widehat{\kappa }_\mathrm{n}F_2.`$ (24)
$`\mathrm{N}=i,f`$ denotes the initial or the final nucleon; $`F_1`$ and $`F_2`$ are the usual Dirac and Pauli form factors; $`Q_\mathrm{n}`$ and $`\widehat{\kappa }_\mathrm{n}`$ are the nucleon charge and anomalous magnetic moment operators. Note that $`\tau Q_i=e_i`$ and $`Q_f\tau =e_f`$ provide effective (Cartesian-basis) charge operators for the nucleons in the present context and that one has $`e_i=e_f+e_\pi `$, describing charge conservation across the $`\pi NN`$ vertex.
Of particular importance in Eq. (18) is the interaction current $`M_{\mathrm{int}}^\mu `$ which originates from the photon attaching itself within the $`t`$-channel $`\pi NN`$ vertex of the pion-pole-dominated diagram of Fig. 1. In lowest order (bare vertices), this corresponds to the usual gauge-invariance-preserving Kroll–Ruderman term as obtained by minimal substitution. In higher orders, with fully dressed vertices, this term contains a dressed Kroll–Ruderman term, exchange currents, and all contributions from final-state interactions .
In evaluating the divergence
$$(p^{}pk)_\mu \overline{u}_f\widehat{J}_{\mathrm{a},\gamma }^{\mu \nu }u_i\epsilon _\nu =q_\mu J_{\mathrm{a},\gamma }^{\mu \nu }\epsilon _\nu ,$$
(25)
it is crucial to note that this involves divergences of the axial current contributions $`\widehat{ȷ}_{\mathrm{a},f}^\mu `$, $`\widehat{ȷ}_{\mathrm{a},i}^\mu `$, and $`\widehat{ȷ}_\pi ^\mu `$ according to Eqs. (11)-(13) which do not vanish even when $`q0`$. The corresponding divergences of the first three and the last terms in Eq. (18), in fact, produce the complete photoproduction amplitude $``$, plus electromagnetic contact terms arising from employing Eq. (11). Indeed, one now easily finds that
$$q_\mu J_{\mathrm{a},\gamma }^{\mu \nu }\epsilon _\nu Q_\pi j_\mathrm{a}^\nu \epsilon _\nu =\frac{f_\pi \mu ^2}{q^2\mu ^2}+\overline{u}_f𝒲^\nu u_i\epsilon _\nu ,$$
(26)
where
$$𝒲^\nu =q_\mu W^{\mu \nu }Q_\pi \widehat{ȷ}_{\mathrm{a},\mathrm{w}}^\nu \frac{\gamma _5\tau \mathrm{\Gamma }_i^\nu +\mathrm{\Gamma }_f^\nu \tau \gamma _5}{2}G_\mathrm{a}(q^2),$$
(27)
with the last term containing the electromagnetic contact contributions. Equation (26) is the desired final result and several remarks are in order now.
To conform to Eq. (15), $`𝒲^\nu `$ should vanish. However, since Eq. (15) was derived with the help of minimal substitution, this is only required in the structureless limit $`G_\mathrm{a}g_\mathrm{a}`$. Evaluating $`W^{\mu \nu }`$ of Eq. (21) in this limit, one easily finds
$$W^{\mu \nu }\gamma _5\left[g^{\mu \nu }\frac{q^\mu (2qk)^\nu }{q^2}\right]\frac{m}{t}g_\mathrm{a}e_\pi ,$$
(28)
which, employing also $`F_1=1`$ and $`F_2=0`$, indeed leads to $`𝒲^\nu =0`$ and thus verifies the validity of Adler’s relation (15) in this limit. In general, however, for nucleons with electroweak structure, it is not obvious that this remains true. One would need a microscopic description of the weak form factor $`G_\mathrm{a}`$ to determine whether $`𝒲^\nu `$ still vanishes.
The derivation of Eq. (26) clearly shows that the entire $`G_\mathrm{a}`$ dependence of its left-hand side is contained solely in $`𝒲^\nu `$ on the right-hand side. In other words, the pion-production amplitude $``$ itself does not depend on $`G_\mathrm{a}`$. Moreover, in view of the explicit expressions available for the axial photoproduction current $`\widehat{J}_{\mathrm{a},\gamma }^{\mu \nu }`$, one finds that evaluating the limit $`q0`$ in Eq. (26) produces simply an identity, but does not provide a constraint that would permit one to extract the threshold behavior of the pion-production amplitude independent of performing that limit in $``$ itself.
For $`q_\mu \widehat{J}_{\mathrm{a},\gamma }^{\mu \nu }`$, in particular, one finds
$`q_\mu \widehat{J}_{\mathrm{a},\gamma }^{\mu \nu }|_{q=0}`$ $`=`$ $`f_\pi \left[M_{\text{int}}^\nu +{\displaystyle \frac{\gamma _5\gamma ^\nu }{2m}}{\displaystyle \frac{m}{f_\pi }}G_\mathrm{a}(k^2)e_\pi \right]+𝒲^\nu `$ (30)
$`+f_\pi {\displaystyle \frac{k^\nu }{k^2}}\gamma _5\left[G_t(k^2){\displaystyle \frac{m}{f_\pi }}G_\mathrm{a}(k^2)\right]e_\pi .`$
The right-hand side here vanishes only in the extreme structureless limit, where all electroweak and hadronic vertices are bare and the interaction current $`M_{\text{int}}^\nu `$ reduces to the Kroll–Ruderman contact term. In general, however, it will be non-zero and, therefore, the often used approximation of assuming that $`q_\mu J_{\mathrm{a},\gamma }^{\mu \nu }\epsilon _\nu `$ vanishes for $`q0`$ is unjustified for physical hadrons with structure. \[Technically, the incorrect limit is obtained if one consistently reduces the entire axial current to its $`\gamma _5\gamma ^\mu `$ part when evaluating Eq. (25).\]
It should be noted that Nambu et al. do not employ this incorrect limit. Nevertheless, their results suffer from an indiscriminate interchange of the limits $`q0`$ and $`\mu 0`$. Clearly, to obtain meaningful threshold results in the chiral limit $`\mu 0`$, one must perform the limit $`q0`$ first. In terms of Eq. (26), the results of correspond to performing on the right-hand side first the $`q`$ limit and then the $`\mu `$ limit, and reversing this order on the left-hand side. On the left-hand side, putting $`\mu =0`$ at the outset makes the hadronic part $`\widehat{ȷ}_{\mathrm{a},\mathrm{h}}^\mu `$ of the axial current vanish from the very beginning and destroys PCAC \[cf. Eq. (12)\]. In other words, $`q_\mu J_{\mathrm{a},\gamma }^{\mu \nu }`$ is evaluated by omitting the third and sixth diagrams from Fig. 3. These are the terms that would normally produce the $`t`$-channel pion-pole contribution and the interaction current. The term $`Q_\pi j_\mathrm{a}^\nu `$ on the left hand-side is then erroneously interpreted as supplying these two terms since it accidentally happens to have the same structure for $`q=\mu =0`$. In view of this incorrect treatment the resulting production current is not gauge-invariant and, therefore, an additional ad hoc current was added in Ref. to repair this deficit. The present derivation shows that when performing the limits correctly, this should not be necessary since the photoproduction current that enters Eq. (26) is gauge-invariant to start with.
It should be emphasized that the derivation of Eq. (26) given here is completely model-independent. It hinges only on describing the weak axial current in terms of the operators defined in Eqs. (7). This corresponds to an effective Lagrangian description completely consistent with the general form (2) of the axial current and with the PCAC hypothesis. This consistency is necessary to avoid an incomplete or partial evaluation of all contributing mechanisms which—since many of the terms have a deceptively similar structure—would lead to erroneous conclusions almost as a matter of course .
The present considerations show that Eq. (26) is devoid of any additional dynamical content that is not already part of the original pion-production amplitude. In fact, it simply provides an alternative definition of the corresponding on-shell amplitude by the reduction formula
$`=\underset{q^2\mu ^2}{lim}{\displaystyle \frac{q^2\mu ^2}{f_\pi \mu ^2}}\overline{u}_f\left[q_\mu \widehat{J}_{\mathrm{a},\gamma }^{\mu \nu }Q_\pi \widehat{ȷ}_\mathrm{a}^\nu 𝒲^\nu \right]u_i\epsilon _\nu ,`$ (31)
where on the right-hand side the dependence on $`G_\mathrm{a}`$ cancels completely even before the limit is taken. This, therefore, does not provide any constraint that goes beyond the original PCAC equation (4) which led to the Goldberger–Treiman relation. In particular, there is no justification in modifying the Kroll–Ruderman term by multiplying it with the axial form factor when considering virtual photons with $`k^20`$.
Within PCAC, therefore, pion electroproduction data at threshold clearly cannot be interpreted in terms of $`G_\mathrm{a}`$, in contrast to what is commonly believed. How this can be reconciled with the findings of chiral perturbation theory remains an open question at present.
This work was supported in part by Grant No. DE-FG02-95ER-40907 of the U.S. Department of Energy.
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# X-ray observations of supernova remnant G54.1+0.3: X-ray spectrum and the discovery of an X-ray jet
## 1 Introduction
Radio source G54.1+0.3 was first suggested to be a Crab-like SNR by Reich et al. (1985) for its flat spectral index of $`\alpha 0.1\pm 0.1`$, filled-center morphology and significant polarization. This identification to G54.1+0.3 was confirmed by Velusamy $`\&`$ Becker (1988) with high resolution multifrequency observations with the VLA and OSRT. In the high resolution VLA maps, G54.1+0.3 has a filled-center brightness distribution peaks around R.A.(2000) =19:30:30, DEC(2000)=18:52:11 and extends to the northeast and north (Velusamy & Becker 1988). They pointed out that these extensions are reminiscent of the radio jets seen in the Crab (Velusamy 1984), CTB80 (Angerhofer et al. 1981) and G332.4+0.1 (Roger et al. 1985).
X-rays from G54.1+0.3 was detected by $`EINSTEIN`$ IPC (resolution $``$ 1$`\mathrm{}`$) with a source strength of 0.016$`\pm `$0.004 counts s<sup>-1</sup> in the energy band 0.5-4.0 keV (Seward 1989). No extent to the X-ray emission was found, due to both its small angular size (2.0$`{}_{}{}^{}\times 1.2^{}`$) (Velusamy & Becker 1988) and its low flux. A power law spectral fitting with energy index of 1.0 gives column density N<sub>H</sub> between 5$`\times `$10<sup>21</sup> and 1$`\times `$10<sup>23</sup> cm<sup>-2</sup>, with the best fit value of 3$`\times `$10<sup>22</sup> cm<sup>-2</sup>, indicating a large distance of this source.
In the paper we present the analyses of $`ROSAT`$ PSPC and $`ASCA`$ GIS and SIS observations of G54.1+0.3. We obtain its spectral information, and, with the aid of an image restoration method, we obtain a high spatial resolution X-ray map of the remnant which clearly shows an X-ray jet pointing to the northeast.
## 2 Observations and analysis method
The $`ROSAT`$ PSPC pointing observation of SNR G54.1+0.3 was carried out from April 11th to 18th, 1991 with a total acceptable observational time of 20271 seconds. We use $`EXSAS`$ (Zimmermann et al. 1998) to analyze its spectrum and produce a 0.1-2.5 keV X-ray image (figure 1) whose spatial resolution is the intrinsic resolution of PSPC (40$`\mathrm{}`$).
G54.1+0.3 was also observed with $`ASCA`$ observatory (Tanaka et al. 1994) continuously from April 27th to 28th, 1997, using the two Gas Imaging Spectrometers (GIS-2 and GIS-3) and the two Solid State Imaging Spectrometers (SIS-0 and SIS-1). Data were collected by the two GIS detectors with a photon time-of-arrival resolution of 4.88$`\times 10^4`$ s in the high bit-rate modes. An effective exposure time 16.5 ks was achieved for each detector. The SIS detectors were operated in the 1-CCD faint mode in which read-out is every 4 s. All SIS data were filtered using the standard screening criteria, which resulted in effective exposures of 19 ks and 20.7 ks for SIS-0 and SIS-1 respectively. Since the two GIS detectors were operated in the high time-of-arrival resolution model, we used the GIS data for temporal analysis. The SIS detectors which are sensitive to photons in 0.5-10.0 keV have superior energy resolution compared to the GIS, and so the SIS data are used for spectral analysis.
Due to the small angular size (120$`\mathrm{}\times 75\mathrm{}`$) of G54.1+0.3 and the limited spatial resolution ($`40\mathrm{}`$) of PSPC, the PSPC observation can not directly give even a coarsely resolved image of G54.1+0.3. In order to obtain an image with higher spatial resolution, we use the widely used Lucy-Richardson formula (Richardson 1972, Lucy 1974) to eliminate the point spread function effect in figure 1. In the iteration process we have used the mean background as the lower limit constraints, in order to improve the quality of the restorated image, as done by Li & Wu (1994), Lu et al. (1996) and Zhang et al. (1998).
## 3 Results
### 3.1 Spectrum from $`ROSAT`$ PSPC observation
The $`ROSAT`$ PSPC spectrum of G54.1+0.3 shows a lack of low energy photons and peaks at energy channel 150 (about 1.5 keV). The spectrum can be fitted with power law model and Raymond-Smith (1977) thermal plasma model. The power law model yields a photon index of -0.8 with 1 $`\sigma `$ error range of -2.8 to 0.0 and an absorption column density of 12.3$`\times `$10<sup>21</sup> cm<sup>-2</sup> with 1 $`\sigma `$ error range of 8 to 20$`\times `$10<sup>21</sup> cm<sup>-2</sup> (see figure 2). The thermal plasma model derives a plasma temperature of 1.8 keV ($`>`$ 1.2keV) and absorption column density of 21.1$`\times `$10<sup>21</sup> cm<sup>-2</sup> with 1 $`\sigma `$ error range of 15-26$`\times `$10<sup>21</sup> cm<sup>-2</sup>. The reduced $`\chi ^2`$ values are almost the same, 0.831 for power law model and 0.834 for thermal plasma model. We adopt the power law model in this paper for it gives the best and the most reasonable fit to the $`ASCA`$ SIS spectrum, as shown in the next section. It gives the absorbed and unabsorbed 0.1-2.4 keV X-ray energy fluxes of 1.0$`\times 10^{12}`$ and 3.4$`\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, respectively. Figure 3 shows the power law model fitting results.
### 3.2 Spectrum from $`ASCA`$ SIS observation
The SIS spectra of the source were extracted within a 4.5 arcminutes radius region. After subtracted the source region, another region of the CCD was used for background subtraction. The source and background spectra obtained from both SISs were added to obtain improved statistics. The spectral anlyses software is XSPEC. Energies above 8 keV were not used because of the poor signal to noise ratio. We have used power law, blackbody, single temperature bremsstrahlung and Raymond-Smith thermal plasma models to fit the spectrum, and found that only the power law model and the thermal bremsstrahlung model give acceptable and reasonable fits. The obtained parameters of the power law model are: photon index $`\alpha `$ -1.9$`{}_{0.2}{}^{}{}_{}{}^{+0.2}`$, column density $`N_H`$ 17.9$`{}_{2.5}{}^{+2.8}\times 10^{21}`$ cm<sup>-2</sup>, 0.7-2.1 keV energy flux 6.5$`\times `$10<sup>-13</sup> erg cm<sup>-2</sup> s<sup>-1</sup>, reduced $`\chi ^2`$ 0.7. Parameters of a thermal bremstrahlung model are: temperature $`T_e`$ 7.9$`{}_{3.1}{}^{}{}_{}{}^{+3.9}`$ keV, column density $`N_H`$ 15.4$`{}_{1.9}{}^{+2.0}\times 10^{21}`$ cm<sup>-2</sup>, 0.7-2.1 keV energy flux 8.6$`\times `$10<sup>-13</sup> erg cm<sup>-2</sup> s<sup>-1</sup>, reduced $`\chi ^2`$ 0.8. We choose the power law model in this paper for that it has the smallest $`\chi ^2`$ and a power law X-ray spectrum is the typical property of the X-ray emission of a Crab-like SNR. For illustration, we show in figure 4 the best fit power law model and its residuals.
### 3.3 Temporal analysis
We examined the $`ASCA`$ GIS data for temporal variability by extracting photons from a 6 arcminutes radius circle centered on the source. A search for coherent pulsations from the source was made by combing the two GIS high-time-resolution data sets (time resolution 4.88$`\times `$10<sup>-4</sup> s) and the arrival times of the used 1805 photons were barycentered. We performed a restricted search for periodic signals between 0.01 s and 2 s using a folding technique (20 phase bins per fold), and detected no pulsation with a significance of more than 3$`\sigma `$ in this period range.
### 3.4 Image restoration results
The distribution of photons detected by $`ROSAT`$ PSPC peak at 1.5 keV and is quite symmetric. We thus use the PSPC point spread function in 1.5 keV and the method described in section 2 to restore the original image (figure 1). The iteration stops after 50 iterations (indeed the restored image is insensitive to the iteration number after 20 iterations). The restored image is shown in figure 5, in which a jet-like feature (hereafter JLF) pointing to the northeast appears, in addition to the $``$30$`\mathrm{}`$ diameter bright nebula coinciding with the brightest radio region. The angular distance from the head of the JLF to the center of the central bright nebula is about 40$`\mathrm{}`$. The total photon flux of the bright nebula is 2.07$`\times `$10<sup>-2</sup> counts s<sup>-1</sup>, that of the JLF is about 3.9$`\times `$10<sup>-3</sup> counts s<sup>-1</sup>, about 430 and 80 photons have been detected from the bright nebula and JLF respectively.
In order to exam the reliability of the restored image, we have performed a Monte-Carlo simulation. Figure 6(a) displays an object similar to G54.1+0.3 in figure 4 in shape and flux. Figure 6(b) is the simulated $`ROSAT`$ PSPC observational result with the same observing time and background level as the real observation to G54.1+0.3 , and figure 6(c) is the smoothed image from 6(b). Figure 6(d) is the restorated image of 6(b). The simulation shows that the high resolution X-ray image of G54.1+0.3 we obtained is reliable.
## 4 Discussions
### 4.1 Distance and X-ray luminosity of G54.1+0.3
Velusamy & Becker (1988) suggested that G54.1+0.3 may have a distance of about 3.2 kpc, if its progenitor is in the star-forming region G53.9+0.3. The galactic HI column density in this direction is about 14.5$`\times `$10<sup>21</sup> cm<sup>-2</sup> (Dickey $`\&`$ Lockman 1990). The best fit column density we get from $`ROSAT`$ PSPC observation is a little lower and the best fit column density of the $`ASCA`$ SIS observation is a little higher than that value. These column densities imply a distance comparable with the radius of the galaxy, similar to the result of $`EINSTEIN`$ IPC (Seward 1989), 3.2 kpc might be then too close and 10 kpc should be a reasonable estimation. The X-ray luminosity in 0.1-2.4 keV is $`L_X`$ = 3.2$`\times `$10$`{}_{}{}^{33}d_{10}^{2}`$ erg s<sup>-1</sup>, where $`d_{10}`$ is the distance to G54.1+0.3 in unit of 10 kpc. If the distance does not deviate from 10 kpc very much, its X-ray luminosity is three or four orders’ lower than Crab Nebula (Helfand & Becker 1987), lies in the lower end of Crab-like SNRs, similar to that of SNR 3C58 (Helfand & Becker 1987; Helfand et al. 1995).
The radio luminosity of G54.1+0.3 is about 5$`\times `$10<sup>33</sup> $`d_{10}^2`$ erg s<sup>-1</sup>, derived from the radio observations of Velusamy & Becker (1988). The ratio $`L_x/L_r`$ = 0.6, also similar to that of 3C58 (Helfand & Becker 1987).
Seward & Wang (1988) found that a relation between the X-ray luminosity ($`L_X`$) of a plerionic SNR and the spin-down luminosity ($`\dot{E}`$) of the central pulsar. Using that relation we can derive $`\dot{E}`$8$`\times `$10<sup>35</sup> erg s<sup>-1</sup> for the central pulsar in G54.1+0.3.
### 4.2 Electron energy distribution
The $`ROSAT`$ PSPC observation of G54.1+0.3 shows that the X-ray flux at 1 keV is about 6.4$`\times `$10<sup>-4</sup> mJy. As the radio flux at 1.4 GHz is 478 mJy, the flux index between radio and X-ray is about -0.7, a little flatter than the $`ASCA`$ SIS obtained X-ray energy index (-0.9) and much steeper than that of radio spectral index (-0.13) (Velusamy $`\&`$ Becker 1998), indicating that the spectrum contains a break between radio and X-ray. Comparing the radio to X-ray flux index, radio flux index and the X-ray flux index, we find that the break is around 10<sup>11</sup> Hz.
If the relativistic electrons have a power law energy distribution $`n_e`$ = $`E^\gamma `$, the spectral index $`\alpha `$ is $`\frac{\gamma +1}{2}`$. The spectral break in the spectral means a similar break in the electron energy distribution. The critical radiation frequency of a relativestic electron with energy $`E`$ in a magnetic field with strength $`B`$ is $`\nu _c=16.1BE^2sin\psi `$ MHz, where $`B`$ is in $`\mu `$G, $`E`$ is in GeV and $`\psi `$ the incident angle of electron (Lang 1998). The electrons whose maximum radiations are at 1 GHz have typical energies of 9.4$`B^{0.5}`$ GeV with $`B`$ the magnetic field strength in $`\mu `$G, assuming that the incident angle is 45$`\mathrm{°}`$. Similarly the X-ray (around 1 keV) emitting relativestic electrons will have typical energies of 94$`B^{0.5}`$ TeV. If the magnetic field is about 10$`\mu `$G, the above estimations show that the electron energy distribution is $``$$`E^{1.3}`$ around 30 GeV and $``$$`E^{2.8}`$ around 300 TeV. An index break exists between 30 GeV and 300 TeV, probably around 300 GeV.
The life time of a relativestic electron can be represented by $`t_{1/2}`$, the time of the electron loses half of its initial energy $`E_0`$, $`t_{1/2}=\frac{8.35\times 10^9}{(Bsin\psi )^2E_0}`$ years, where $`B`$ is in $`\mu `$G, $`E_0`$ is in GeV and $`\psi `$ the incident angle of the electron (Lang 1998). The lifetimes for the 30 GeV, 300 GeV and 300 TeV photons in 10 $`\mu `$G magnetic field are about 5.6$`\times `$10<sup>6</sup>, 5.6$`\times `$10<sup>5</sup>, 5.6$`\times `$10<sup>2</sup> years, respectively. These three typical lifetimes will be used in the discussions of the origin of the electron energy distribution index break in the next paragraph.
If the electrons from the center pulsar have a continuous power law energy distribution initially, the observed break should be due to the short lifetime of the high energy electrons. Because the low energy electrons which radiate radio emission have a long lifetime, their energy distribution represents the initial electron energy distribution well. The initial energy flux ratio $`\frac{f_{1.4GHz}}{f_{1keV}}`$ is then 11.75. The currently observed energy flux ratio $`\frac{f_{1.4GHz}}{f_{1keV}}`$ is 7.8$`\times 10^5`$, indicates that the age of G54.1+0.3 would be at least $`\frac{7.8\times 10^5}{11.75}`$$`\times `$5.6$`\times `$10<sup>2</sup>=3.7$`\times `$10<sup>7</sup> years. This large age value shows that the observed electron energy distribution break is quite probably an intrinsic property of the electrons from the central pulsar.
### 4.3 X-ray jet
For the first time an well resolved X-ray image of G54.1+0.3 has been obtained. It shows a JLF pointing to the northeast. The simulation shows that such a structure can be clearly resolved by $`ROSAT`$ PSPC with the aid of an image restoration technique. The simulation also shows that this feature can not be attributed as the fluctuations of the bright source, it is an intrinsic structure of the object.
We have studied the possibility that the JLF is indeed a separate object lies in a similar direction with G54.1+0.3. We find that there is no identified object in the 30$`\mathrm{}`$ vicinity of the JLF except G54.1+0.3. The optical plate obtained by Palomar Observatory Sky Survey and electronically reproduced by Skyview of NASA/GSFC shows no source in the JLF region too. The JLF shown in figure 5 shows some enhancements in the head. But it might be a false phenomenon caused by the low quality of the original data and the restoration process, as can be found in the simulation, although some similar structures exist in the 4.8 GHz radio map. More simulations show that the length of the JLF is quite reliable, the width of the JLF might have an uncertainty up to $``$ 50$`\%`$.
We have compared figure 4 with the 4.8 GHz VLA map obtained by Velusamy & Becker (1988) in details. The brightest point of the extended X-ray source locates at R.A.(2000)=19:30:30.0, DEC(2000)=18:52:07, which coincides with the brightest region of the radio source. The head of the JLF has a coordinate of R.A.(2000)=19:30:32.2, DEC(2000)=18:52:31, which also coincides with the northeastern enhancement in the radio map. The nice position coincidence of the X-ray and radio sources strongly favor their same origin. However, the X-ray source has a smaller extent than the radio source and no significant X-ray emission has been detected along the northward feature, which was suggested to be the most probable radio JLF by Velusamy & Becker (1988). It might be due to the intrinsic deficiency or the limited sensitivity of the present observation.
There are two possible ways to explain the origin of the X-ray JLF. One is that it is a fragment produced in the supernova explosion, like the fragments detected around the Vela SNR, especially its ‘bullet’-like fragment A. (Aschenbach et al. 1995; Strom et al. 1995). However, significant radio emission has only been detected around the head of the fragments, implies that most of the relativistic electrons are in the leading edge of the fragments, close to the shock front (Strom et al. 1995). But in the case of G54.1+0.3 the radio emission has a similar distribution with the X-ray JLF, indicating a similar distribution of relativistic electrons with the X-ray brightness. It makes the fragment origin of the X-ray JLF implausible. The second is that the X-ray JLF is due to the relativistic electrons produced by the central pulsar, like X-ray jets detected in PSR 1929+10 (Wang et al. 1993), Crab SNR (Hester et al. 1995), Vela pulsar (Markwardt & Ögelman 1995), SNR MSH 15-52 (Tamura et al. 1996), SNR CTB80 (Wang & Seward 1984; Safi-Harb et al. 1995) in the galaxy and SNR N157B in the Large Magellanic Cloud (Wang & Gotthelf 1998). The coincidence of radio and X-ray emission in the case of Vela pulsar jet (Frail et al. 1997) and that of SNR N157B (Wang & Gotthelf 1998) strongly support this scenario. We conclude that the JLF we discoveried is quite probably an X-ray jet connected with the pulsar in G54.1+0.3.
The X-ray emission of Vela pulsar jet can be fitted with both power law and thermal plasma model (Markwardt & Ögelman 1995), and the X-ray pulsar jet in MSH15-52 appears to be nonthermal. It is difficult to get the spectral properties of the X-ray jet in G54.1+0.3 with the present data. We assume that it share the same power law model with the whole remnant, and is due to the synchrotron radiation of relativistic electrons from the pulsar. The X-ray luminosity of the jet in 0.1-2.4 keV is then about 5.1$`\times `$10<sup>32</sup>$`d_{10}^2`$ erg s<sup>-1</sup>.
From the radio map of Velusamy $`\&`$ Becker (1988) we estimate that the flux of the jet at 4.8 GHz is about 40 mJy. Its X-ray flux at 1 keV is about 9.4$`\times 10^5`$ mJy. The two fluxes give a spectral index from radio to X-ray of about -0.73, quite similar to that of the whole remnant. As no significant radio spectral variation across the source has been detected (Velusamy $`\&`$ Becker 1988), the jet electrons have a break with the energy distribution too, similar to the whole remnant.
The distance of the jet head to the nebula center is about 40$`\mathrm{}`$. It corresponds to 2 pc if the SNR is 10 kpc away. Reccent distance measurements to Vela SNR obtained a distance of 250$`\pm 30`$ pc (Cha et al. 1999). If so the Vela pulsar jet is about 3 pc long (Cha et al. 1999, Markwardt & Ögelman 1995). The lengths of the two jets are quite similar.
## 5 Summary
$`ROSAT`$ PSPC and $`ASCA`$ observations of G54.1+0.3 imply a large distance comparable with the galactic radius. Its X-ray spectrum is of nonthermal origin. The comparison of the radio and X-ray emissions shows that the energy distribution of the relativistic electrons has a break around 300 GeV. This break is quite probably an intrinsic property of the relativistic electrons from the central pulsar instead of due to the energy lose in the synchrotron radiation process, if G54.1+0.3 is not as old as 3.7$`\times `$10<sup>7</sup> years.
A high spatially resolved image shows an X-ray jet pointing to the northeast, similar to the radio structures. Its nonthermal spectrum and the existence of X-ray jet confirm the formal identification of G54.1+0.3 as a Crab-like SNR, though no pulsation has been found in the X-ray observation. Future deep X-ray observations with high spatial resolution and spectral resolution telescopes such as Chandra and XMM are invaluable to find out the spectral and spatial structure of the remnant as well as the X-ray jet.
###### Acknowledgements.
F.J. Lu is supported by the exchange program between Max-Planck Society and Chinese Academy of Sciences. He thanks Professor J. Trümper for hospitality. The authors thank Drs S.D. Mao and Q.D. Wang for helpful discussions. This research is partially supported by the National Natural Science Foundation of China and the Special Funds for Major State Basic Research Projects. It has made use of the SIMBAD database, operated at CDS, Strasbourg, France and the Digitized Sky Survey operated by Skyview of NASA/GSFC.
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# Self-similar collapse of collisional gas in an expanding Universe
## 1 Introduction
On scales larger than a few megaparsecs, pressure forces in the baryonic matter in the universe are negligible, so the evolution of dark and baryonic matter is mainly determined by gravity. On small scales pressure becomes important and may segregate between the evolution of baryonic and dark matter. Pressure forces, cooling of gas, and star formation feedback, are key ingredients in galaxy formation. These ingredients combine to cause differences between the distributions of galaxies and dark matter (biasing), even on large scales where these effects are not directly important (e.g., Kaiser 1984, Dekel & Rees 1987, Kauffmann, Nusser & Steinmetz 1997, Benson et. al. 2000). On scales smaller than the Jeans length of the photo-heated intergalactic medium (IGM), pressure forces dominate gravity and can prevent the collapse of gas into dark haloes below a certain mass threshold. For haloes massive enough the temperature of the IGM can be neglected and the gas falls into the halo. The mean free path for collisions between gas particles inside a halo is $`\left(200/\delta _c\right)\left(10^{15}/\sigma \right)(1+z)^3\left(0.1/\mathrm{\Omega }_bh^2\right)1.6\mathrm{p}c,`$ where $`\delta _c`$ is the overdensity inside the virial radius and $`\sigma `$ is a typical cross section for collisions in units of $`\mathrm{c}m^2`$. This is smaller than the virial radius of a typical halo by a few orders of magnitude. Therefore, on its infall into the halo, the gas is likely to form shocks and transform its kinetic energy into heat. The hot dense gas can then cool to form stars which explode and inject energy into the halo gas. Detailed study of these processes under general conditions is not feasible. One can aim at a global parameterization based on general physical requirements which match observational data (Kauffman, White, Guiderdoni 1993, Somerville & Primack 1999, Cole et. al. 1994). Another route would be to study special aspects which can be treated by either numerical or analytical methods. Here we focus on the collapse of the baryonic gas in an Einstein-De Sitter universe, ignoring the gas initial temperature, cooling and heating processes. We assume that the collapse initiates from a symmetric scale free density peak, and that the velocity of each shell in the peak is taken to match the general expansion of the universe. The energy of each shell is negative and it will expand up to a maximum distance before it starts falling towards the center of the perturbation. The maximum distance is termed the turnaround radius. Shell crossing is not allowed and the collapse can proceed in two distinct ways, either a shock wave forms, or shells accumulate at the center. Which of these possibilities actually occurs, depends on the physical conditions at the center. If the velocity vanishes at the center than a shock wave forms. If on the other hand physical conditions allow a non vanishing velocity at the center then the shells accumulate at the center (Bertschinger 1985). Inner boundary conditions can be arranged so that a shock is accompanied by the accumulation of central mass, but a proper stability analysis is needed to determine whether or not this is possible (Bertschinger 1985). In this paper we focus on shocked collapses without the formation of a central mass. Since the initial gas pressure is negligible, the collapse eventually develops in a self-similar way where the only relevant scale at any time is the radius of the shell at maximum expansion. Bertschinger (1985), and Forcada-Miro & White (1997) have studied similarity solutions in spherically symmetric perturbations with initial relative mass excess $`\delta M/Mr^3`$, and $`r^2`$, respectively. Here we derive similarity solutions in planar, cylindrical, and spherical geometries, for the collapse of a perturbation with $`\delta M/Mr^s`$ for any $`s>0`$ and adiabatic index $`\gamma 4/3`$.
In section 2 we write the equations of motion for symmetric perturbations in planar, cylindrical, and spherical geometry. In section 3 we discuss the asymptotic behaviour of the fluid variables near the center, in the case of shocked collapse. In section 4 we present results of numerical integrations of the equations. In section 5 we derive asymptotic behaviour of the fluid variables in a universe dominated by collisionless dark matter. In section 6 we conclude with a discussion of the results and their potential astrophysical consequences.
## 2 The equations
We write the Newtonian equations of motion governing the adiabatic collapse of symmetric perturbations in a collisional fluid (gas) of adiabatic index $`\gamma `$ and zero initial pressure. Except section 5, we restrict the analysis here and throughout to the collapse in a flat universe containing collisional gas only. The initial gas pressure is zero, so the expansion scale factor of the universe is $`a(t)t^{2/3}`$, the Hubble function is $`H(t)=2/(3t)`$, and the background density is $`\rho _c=3H^2/(8\pi G)=1/(6\pi Gt^2)`$.
Denote by $`r`$ and $`\upsilon \mathrm{d}r/\mathrm{d}t`$ the physical position and velocity of a gas shell, where $`r=0`$ is the symmetry center of the perturbation. Further, let $`\rho (r,t)`$ and $`p(r,t)`$ be the gas density and pressure at $`r`$. As in Fillmore & Goldreich (1984) define the mass within a distance $`r`$ from the symmetry center by $`m(r,t)=_0^rx^{n1}\rho (x,t)dx`$, where $`n=1,2`$, and 3 refer, respectively, to planar, cylindrical, and spherical perturbations. The mass within a fixed shell varies with time like $`mt^{2(3n)/3}`$, because of the Hubble expansion along $`3n`$ of the axes. In this notation, the equations of motion are, the continuity equation,
$$\frac{\mathrm{d}(\rho t^{\frac{2(3n)}{3}})}{\mathrm{d}t}=t^{\frac{2(3n)}{3}}\rho r^{1n}_r(r^{n1}\upsilon ),$$
(1)
Euler,
$$\frac{\mathrm{d}\upsilon }{\mathrm{d}t}\frac{2}{9}\frac{3n}{n}\frac{r}{t^2}=\frac{_rp}{\rho }\frac{4\pi Gm}{r^{n1}},$$
(2)
adiabatic condition,
$$\frac{\mathrm{d}}{\mathrm{d}t}(p\rho ^\gamma )=0,$$
(3)
and the relation,
$$_rm=r^{n1}\rho .$$
(4)
These equations are satisfied everywhere except at the shock where the fluid variables are described by jump conditions obtained from mass, momentum, and energy conservation. The initial conditions leading to self-similar collapse are specified at an early time close to zero, $`t_i`$, as
$`{\displaystyle \frac{\delta M}{M}}`$ $`=`$ $`\left({\displaystyle \frac{r}{r_0}}\right)^s,`$ (5)
$`\upsilon (r,t_i)`$ $`=`$ $`{\displaystyle \frac{2}{3t_i}}r,`$ (6)
$`p(r,t_i)`$ $`=`$ $`0,`$ (7)
where $`\delta M/M`$ is the mean density contrast interior to $`r`$, and $`s>0`$. For cosmological initial conditions the initial density contrast must be tiny, so we will be interested in the solution in the region $`rr_0`$. A perturbation with $`s>n`$ can be realized by placing a high narrow positive density peak at the center ($`r_0`$) of a symmetric void with local density contrast $`(r^s)`$. The condition (6) means that a gas shell at $`r`$ moves initially with the general universal expansion. This condition can be relaxed to allow for a non vanishing initial zero peculiar velocity according to late time linear theory (e.g., Peebles 1980). However, this does not affect the details of the collapse (Peebles 1980, Bertschinger 1985), so we use (6) which is commonly adopted in the literature. Bertschinger (1985) and White & Forcada (1997), respectively, considered the collapse of spherical perturbations with $`s=3`$, and $`s=2`$.
The equations of motion (14) together with the initial conditions (57) are insufficient to completely determine the evolution of the perturbation. Still missing is an inner boundary condition specifying the velocity and mass at $`r=0`$, for $`tt_i`$. For a shock to develop without the accumulation of a central mass (a black hole for $`n=3`$) we must have $`v(r=0,tt_i)=0`$ and $`m(r=0,tt_i)=0`$. Relaxing the condition $`V(0)=0`$ leads to a non-vanishing central mass with or without the presence of a shock.
In a critical density universe ($`\mathrm{\Omega }=1`$) the only length scale relevant to the collapse is the scale of non-linearity. At any time, $`t`$, this scale can be defined as the distance of the shell at the maximum expansion, i.e., the shell with $`\upsilon =0`$ (Gunn 1977, Fillmore & Goldreich 1984, Bertschinger 1985). This radius is termed the current turnaround radius, $`r_{ta}(t)`$. Starting from tiny initial density contrast, the mean overdensity (density in units of $`\rho _c`$) interior to $`r_{ta}(t)`$ is a fixed number independent of time. For time $`tt_i`$, when shells with $`rr_0`$ reach their turnaround, the collapse develops a self-similar behaviour that depends on $`r`$ and $`t`$ through the combination $`\lambda =r/r_{ta}`$. The turnaround radius $`r_{ta}(t)`$ is given by (e.g., Fillmore & Goldreich 1984),
$$r_{ta}=r_0\frac{C_x}{C_t^{3\alpha /2}}\left(\frac{t}{t_i}\right)^\alpha ;\alpha =\frac{2}{3}\frac{s+1}{s}$$
(8)
where,
$$C_x=\frac{5}{12},0.741,1;C_t=\frac{5}{6},1.386,\left(\frac{3\pi }{4}\right)^{2/3},$$
(9)
for $`n=1,2,`$ and 3, respectively. The turnaround radius grows faster than the scale factor $`at^{2/3}`$. This is because the mass, $`\rho _cr_{ta}(t)^3`$, interior to $`r_{ta}`$ must grow with time while the mass, $`\rho _ca^3(t)`$, inside a fixed shell in a homogeneous universe is constant. For $`s<2`$ the turnaround radius grows faster than $`t`$ reaching the horizon scale in finite time. When this happens relativistic description must be used and $`r_{ta}`$ ceases to be the only scale in the problem (Fillmore & Goldreich 1984).
The equations can be cast into a non-dimensional form using the scaled variables $`V(\lambda )`$, $`D(\lambda )`$, $`P(\lambda )`$, and $`M(\lambda )`$ defined by (Bertschinger 1985),
$`\upsilon (r,t)`$ $`=`$ $`{\displaystyle \frac{r_{ta}}{t}}V(\lambda )`$ (10)
$`\rho (r,t)`$ $`=`$ $`\rho _cD(\lambda )`$ (11)
$`p(r,t)`$ $`=`$ $`\rho _c\left({\displaystyle \frac{r_{ta}}{t}}\right)^2P(\lambda )`$ (12)
$`m(r,t)`$ $`=`$ $`{\displaystyle \frac{1}{3}}\rho _cr_{ta}^nM(\lambda ).`$ (13)
Expressed in terms of these variables, the equations (1-4) become, respectively,
$$\left(V\alpha \lambda \right)D^{}+\left(\frac{n1}{\lambda }V+V^{}\frac{2n}{3}\right)D=0,$$
(14)
$$\left(\alpha 1\right)V+\left(V\alpha \lambda \right)V^{}\frac{2}{9}\frac{3n}{n}\lambda =\frac{P^{}}{D}\frac{2}{9}\frac{M}{\lambda ^{n1}},$$
(15)
$$\left(\gamma \frac{D^{}}{D}\frac{P^{}}{P}\right)\left(V\alpha \lambda \right)=2\left(\alpha 2+\gamma \right),$$
(16)
$$M^{}=3\lambda ^{n1}D,$$
(17)
where the prime symbol denotes derivatives with respect to $`\lambda `$.
We will mainly be concerned with solutions for shocked collapse with vanishing mass at the center. The inner boundary condition appropriate for this collapse are vanishing mass and velocity at $`\lambda =0`$, i.e.,
$$V(0)=0\mathrm{and}M(0)=0,$$
(18)
Self-similarity implies that the shock appears at fixed $`\lambda =\lambda _s=r_s/r_{ta}`$, so the physical radius of the shock $`r_st^\alpha `$ and its non-dimensional speed is $`(r_{ta}/t)^1(\mathrm{d}r_s/\mathrm{d}t)=\alpha \lambda _s`$. At the surface of the shock the fluid variables satisfy the jump conditions obtained from mass, momentum, and energy conservation. In terms of the non-dimensional fluid variables, the jump conditions appropriate for an adiabatic shock are,
$`V^+`$ $`=`$ $`\alpha \lambda _s+{\displaystyle \frac{\gamma 1}{\gamma +1}}(V^{}\alpha \lambda _s),`$ (19)
$`D^+`$ $`=`$ $`{\displaystyle \frac{\gamma +1}{\gamma 1}}D^{},`$ (20)
$`P^+`$ $`=`$ $`{\displaystyle \frac{2}{\gamma +1}}D^{}(V^{}\alpha \lambda _s)^2,`$ (21)
$`M^+`$ $`=`$ $`M^{},`$ (22)
where the superscripts of the minus and plus signs refer to pre- and post-shock quantities. In employing energy conservation we have taken $`\frac{P}{D(\gamma 1)}`$ as the non-dimensional internal energy per unit mass.
In section 4 we will find numerical solutions satisfying the requirements for shocked collapse without a central mass. Except spherical perturbations with $`\gamma =4/3`$ only one value $`\lambda _s`$ can yield solutions satisfying these requirements. Spherical perturbations with $`\gamma =4/3`$ allow a range of values for $`\lambda _s`$. Before presenting the numerical solutions we derive in the next section the asymptotic behaviour of the fluid variables near the center, and two integrals of motion which will be used as a check on the numerical solutions.
## 3 Integrals of motion and asymptotic behaviour near the center
Solutions to (1417) with the appropriate jump and boundary conditions for all $`\lambda `$ will be found by numerical integration. We present here an analytic treatment of the equations to derive the asymptotic behaviour of the fluid variables near $`r=0`$, and two integrals of motion (e.g., Bertschinger 1983, 1985). We restrict the analysis shocked collapses satisfying the inner boundary condition (18). All fluid variables can be expressed in terms of an auxiliary function $`K(\lambda )\mathrm{exp}_0^\lambda \frac{\mathrm{d}x}{V(x)\alpha x}`$ as follows.
$`V(\lambda )`$ $`=`$ $`{\displaystyle \frac{K}{K^{}}}+\alpha \lambda ,`$ (23)
$`{\displaystyle \frac{D(\lambda )}{D_0}}`$ $`=`$ $`\lambda ^{1n}K^{}K^{1+n(2/3\alpha )},`$ (24)
$`{\displaystyle \frac{P(\lambda )}{P_0}}`$ $`=`$ $`\lambda ^{(1n)\gamma }K^{4+(2n/33)\gamma \alpha (2+n\gamma )}K^\gamma ,`$ (25)
$`{\displaystyle \frac{M(\lambda )}{D_0}}`$ $`=`$ $`{\displaystyle \frac{9K^{n(2/3\alpha )}}{n(3\alpha 2)}},`$ (26)
where $`D_0`$ and $`P_0`$ are constants. The fluid variables in (2326) satisfy the non-dimensional equations (14), (16), and (17) for any functional form of $`K`$. The function $`K`$ is then specified by only one equation, the non-dimensional Euler equation (15).
Two integrals of motion can immediately be found from (2326). These are the mass and entropy integrals of motion (Bertschinger 1983, 1985),
$`M`$ $`=`$ $`{\displaystyle \frac{9}{n(23\alpha )}}D(V\alpha \lambda )\lambda ^{n1},`$ (27)
$`PD^\gamma M^\zeta `$ $`=`$ $`\mathrm{const};\zeta ={\displaystyle \frac{6}{n}}{\displaystyle \frac{\alpha +\gamma 2}{23\alpha }},`$ (28)
where all fluid variable are evaluated at any $`\lambda `$ inside the shock. Since $`\zeta <0`$, the entropy integral of motion means that the entropy $`\mathrm{ln}(PD^\gamma )`$ is an increasing function of the mass and hence $`\lambda `$.
The auxiliary function greatly simplifies the derivation of the asymptotic behaviour of the non-dimensional fluid variables near $`\lambda =0`$. Since $`V(0)=0`$ and $`M(0)=0`$, equations (23) and (26) imply that, to first order, $`K(\lambda )`$ must approach
$$K(\lambda )=\lambda ^{1/(V_0\alpha )}$$
(29)
as $`\lambda 0`$, where $`V_0`$ is an arbitrary constant. Substituting this expression for $`K`$ in (2326), yields
$`V`$ $`=`$ $`V_0\lambda ,`$ (30)
$`M`$ $`=`$ $`{\displaystyle \frac{9D_0}{n(3\alpha 2)}}\lambda ^{\delta +n},`$ (31)
$`D`$ $`=`$ $`{\displaystyle \frac{D_0}{\alpha V_0}}\lambda ^\delta ,`$ (32)
$`P`$ $`=`$ $`{\displaystyle \frac{P_0}{(\alpha V_0)^\gamma }}\lambda ^\eta ,`$ (33)
and the asymptotic exponents are expressed in terms of the coefficient $`V_0`$ as
$`\delta `$ $`=`$ $`{\displaystyle \frac{n(23V_0)}{3(V_0\alpha )}},`$ (34)
$`\eta `$ $`=`$ $`{\displaystyle \frac{42\alpha 2\gamma +(\frac{2}{3}V_0)n\gamma }{V_0\alpha }}.`$ (35)
These relations have been obtained without using the Euler equation (15). In order to determine the exponents uniquely we use the Euler equation which adds the following constraints,
$`\eta =2(\delta +1),\mathrm{for}\eta <0,`$
$`2(\delta +1)0,\mathrm{for}\eta =0.`$ (36)
A value $`\eta >0`$ gives $`P_0<0`$ and so, by (28), a negative entropy integral of motion. But the jump condition (21) gives positive pressure, $`P^+`$, just behind the shock and since the density and mass are also positive, the entropy integral (28) must be positive. So we rule out $`\eta >0`$.
If the solution to (34), (35), and (36) is $`\eta <0`$ then the Euler equation also provides the following constraint on the coefficients $`P_0`$ and $`D_0`$,
$$P_0=\frac{2D_0^2\alpha (V_0\alpha )^{\gamma 1}}{n\eta (3\alpha 2)}.$$
(37)
Table 1 lists the values of $`V_0`$, $`\eta `$, and $`\delta `$ in all cases. The dashed curves in figure 1 are a graphical representation of the density exponent $`\delta `$ versus $`s`$ for $`\gamma =5/3`$. In planar geometry, $`n=1`$, the only possible solution to equations is $`\eta =0`$ and the pressure is finite everywhere. In cylindrical geometry, $`n=2`$, we have $`\eta <0`$ only for $`\alpha <\frac{5\gamma 6}{3(\gamma 1)}`$, so it must be zero for other values of $`\alpha `$. In spherical geometry, $`n=3`$, if $`\alpha <2`$ and $`\gamma >4/3`$ then $`\eta =2(1\frac{2}{\alpha })`$, and $`V_0=0`$ meaning that $`V(\lambda )\lambda ^\nu `$, where $`\nu >1`$. A second order expansion gives $`\nu =1+2\alpha ^1\left[(45\alpha )\pm \gamma ^{1/2}\sqrt{8\alpha (\gamma 8)+16\gamma +\alpha ^2(32+\gamma )}\right]`$, where only one of the roots is $`\nu >1`$. In the limit of either $`\alpha 2`$ or $`\gamma 4/3`$, we have $`\nu 1`$.
For $`n=3`$ and $`\gamma =4/3`$ the relations (34), (35), and (36) allow multiple solutions for $`V_0`$ and consequently for $`\lambda _s`$. A second order calculation gives the upper limit $`V_0<45\alpha `$. Solutions exist for any positive $`\lambda _s`$ smaller than a maximal value which corresponds to the upper limit on $`V_0`$. This means that for $`\gamma =4/3`$ a shocked collapse cannot be accompanied by the presence of a non-vanishing mass at the center. In the table we list the asymptotic constants corresponding to $`V_0=45\alpha `$, i.e., the maximal value of $`\lambda _s`$ for which a shocked collapse occurs. Bertschinger (1984) does not mention that there are solutions for shocked without a central mass for a range of $`\lambda _s`$. His numerical solution with $`\gamma =4/3`$ seems to correspond to the maximal $`\lambda _s`$ and yields $`\eta =3.2`$ and $`\delta =2.6`$, instead of $`3`$ and $`2.5`$ as listed in table 1.
We now examine how the dimensional density, $`\rho (r,t)=\rho _cD`$, varies with time near the center. Using $`\rho _ct^2`$, the first order expression $`D(r/r_{ta})^\delta `$, and $`r_{ta}t^\alpha `$, we find $`\rho (r,t)r^\delta /t^{2+\alpha \delta }`$. When $`2+\delta \alpha =0`$ the density is constant with time. In spherical geometry a time independent density is equivalent to $`V_0=0`$. Because of the expansion in the $`n1`$ directions, a vanishing $`V_0`$ in planar and cylindrical geometries leads to $`\rho t^{2\alpha \delta }t^{2(n1)/3}`$. In all geometries a vanishing $`V_0`$ indicates that, to first order in the asymptotic expansion, the gas is pressure supported in hydrostatic equilibrium. According to the table, spherical perturbations have $`V_0=0`$ only for $`s>1/2`$. Note that for $`2>s>1/2`$ the asymptotic density is constant even though the temperature $`P/Dr^{\eta \delta }`$ decreases inward. For $`s<1/2`$, the density increases with time. According to table 1, planar and cylindrical perturbations have $`V_0=0`$ only at $`s=1/(2\gamma )`$ and $`s=2/(4\gamma )`$, respectively.
## 4 Numerical integration
We present results of the numerical integration of the non-dimensional equations of motion. The numerical solutions shown here describe shocked collapses without a central mass. Outside the shock the fluid variables are given from the solution for collapse with zero pressure (Zel’dovich 1970, Peebles 1980, Fillmore & Goldreich 1984).
The shock position $`\lambda _s`$ is unknown a priori. We have to find its value such that the fluid variables satisfy the equations of motion, (1417), the jump conditions, (1922), and the inner boundary condition, (18). Assuming that the pre-shock variables are given from the zero pressure solution, the value of $`\lambda _s`$ can be found as follows (Bertschinger 1985, Forcada–Miro & White 1997). For an assumed value for $`\lambda _s`$, we obtain the post-shock variables using the jump conditions. We then integrate the equations of motion from $`\lambda _s`$ inward to $`\lambda =0`$, and tune $`\lambda _s`$ so that the solution gives $`V=0`$ and $`M=0`$ at $`\lambda =0`$. In all numerical solutions we find that if $`M(0)=0`$ for a given $`\lambda _s`$ then $`V(0)=0`$, and vice verse.
The zero pressure solutions in planar and spherical geometries are known analytically ( Zel’dovich 1970, Fillmore & Goldreich 1984) so the pre-shock fluid variable can be found directly, for an assumed $`\lambda _s`$. However in cylindrical geometry an analytic solution is not available and we numerically integrate the equations with zero pressure to obtain the pre-shock quantities. In practice we use numerical integration also in planar and spherical geometries. It is convenient to integrate the zero pressure equations from $`\lambda =1`$, i.e., the turnaround radius, to $`\lambda _s`$. At $`\lambda =1`$ the fluid variables are
$$V(1)=0,P(1)=0,$$
(38)
$$D(1)=\frac{1}{s+1}\left(\frac{C_t}{C_x}\right)^n,M(1)=\frac{3}{n}\left(\frac{C_t}{C_x}\right)^n,$$
(39)
where the expressions for $`D(1)`$ and $`M(1)`$ were obtained from the zero pressure solution (e.g., Fillmore & Goldreich 1984).
The sequence of figures (25) shows the numerical solutions for the over density $`D`$, the pressure $`P`$, the velocity $`V`$, and the thermal energy $`U=\frac{P}{D(\gamma 1)}`$ in spherical, cylindrical and planar geometries for several values of $`s`$, all as a function of $`\lambda `$. All curves are obtained from the solutions with $`\gamma =5/3`$. The solid curve in each plot corresponds to $`s=n`$. The sudden change in the fluid variables indicate the location of the shock. All numerical solutions satisfy the integrals of motion (27) and (28) up to the numerical accuracy and agree with asymptotic behaviour of the previous section. The logarithmic slopes of all curves match the corresponding values listed in table 1. Spherical perturbations with $`s=3`$ and $`s=2`$ were, respectively, analyzed by Bertschinger (1984) and Forcada–Miro & White (1997). The agreement between their solutions with $`\gamma =5/3`$ and ours is excellent. In all geometries the density and pressure in the shock are higher for larger $`s`$. Figures 4 of the velocity and figure 5 of the internal energy $`U=\frac{P}{D(\gamma 1)}`$ demonstrate that particles are decelerated at the shock converting most of their kinetic energy into heat. The velocity near the center of cylindrical perturbations can be positive inside the shock, hence the linear vertical scale in velocity plot in this case.
In figure 6 we plot the location of the shock $`\lambda _s`$ as a function of $`s`$ for several values of $`\gamma `$. The shock location always increases with $`s`$ and agrees with the values obtained by Bertschinger (1985) and Forcada-Miro & White (1997) for spherical perturbations with $`s=3`$ and $`s=2`$, respectively. A special case is spherical collapse with $`\gamma =4/3`$. Here a solution for shocked collapse without a central point mass is possible for any $`\lambda _s`$ less than a maximal value $`\lambda _c`$. In this case we plot the maximal value $`\lambda _c`$. For $`s1`$, $`\lambda _s`$ can exceed unity increasing to a finite value as $`s\mathrm{}`$.
The variation of the fluid variables attached to a given fluid elements are also of interest. Using the velocity obtained from the self-similar solutions we can calculate the trajectory $`r(t)`$ of a fluid element (particle) as a function of time. The fluid variable associated with the particle at any time can then be obtained by interpolating the solutions at the particle’s position at that time. The three panels from top to bottom in figure 7 show, respectively, the trajectory, density, and pressure of a particle as obtained from the solutions with $`\gamma =5/3`$ for three values of $`s`$. The time axis in all panels is $`t/t_{ta}`$ where $`t_{ta}`$ is the time at which the particle reached its maximum expansion (the turnaround time). The particle position, $`r`$, density $`\rho `$, have been scaled by their values at $`t_{ta}`$, $`r_{ta}^{}=r(t_{ta})`$, and $`\rho =\rho (t_{ta})`$, respectively. The pressure, $`p`$, has been scaled by its value at the shock $`p_s`$. With this scaling the curves in the figures become valid for all the particles. The particle trajectory for $`s=3`$ (solid line, top panel) agrees with the corresponding curve in figure 4 of Bertschinger (1985). As expected from the asymptotic solution the particles for $`s=3`$ and $`s=10`$ tend to settle at a physical distance which is a fixed fraction of their turnaround radii. The decay of the trajectory for $`s=1/4`$ at late time can be evaluated using the asymptotic expansion. Near the origin $`V(\lambda )=V_0\lambda `$, so $`\mathrm{d}r/\mathrm{d}t=(r_{ta}/t)V_0\lambda =r/t`$ and $`rt^{V_0}`$. According to table 1 $`V_0=8/15`$ for $`s=1/4`$ and $`\gamma =5/3`$ so $`r`$ decays like $`t^{8/15}`$. The density (middle panel) and pressure (bottom) curves for $`s=3`$ and $`s=10`$ flatten at late times, consistent with the settling of the particles to a constant $`r`$.
## 5 Asymptotic behaviour in a Universe dominated by collisionless matter
So far we have considered similarity solutions for collapse involving gas only. The collapse of scale free symmetric perturbations in a an Einstein-De Sitter universe containing a mixture of gas and collisionless matter is also self-similar. Similarity solutions for mixed collapse are beyond the scope of the present paper. Here we only obtain the asymptotic behaviour of the gas variables in a universe dominated by collisionless matter. For a self-similar collapse to develop, the two matter components must start with the same initial conditions with zero initial gas pressure. Shells of gas and collisional matter then move together until they reach either a shock in the gas or the region of shell crossing in the collisionless component. The evolved density profiles of both components depend on $`r`$ and $`t`$ through $`\lambda `$. Although the global gas mass fraction is negligible, the gas can be gravitationally dominant at the center of the collapse if it has a steeper density profile than the collisionless matter. For the purpose of deriving the asymptotic exponents we proceed assuming that the gravity of the gas is negligible everywhere and check the consistency of this assumption according to the results. So we replace the gas mass $`M(\lambda )`$ in the non-dimensional Euler equation (15) by the collisionless matter mass $`\stackrel{~}{M}(\lambda )`$. Writing $`\stackrel{~}{M}\lambda ^{\stackrel{~}{\delta }+n}`$ near $`\lambda =0`$, Fillmore & Goldreich (1984) find that $`\stackrel{~}{\delta }`$ is
$$n=1:\stackrel{~}{\delta }=\frac{s}{s+1}$$
(40)
$$n=2:\stackrel{~}{\delta }=1$$
(41)
$`n=3:\stackrel{~}{\delta }`$ $`=`$ $`2,\mathrm{for}s2`$ (42)
$`=`$ $`{\displaystyle \frac{3s}{s+1}},\mathrm{for}s>2`$
Table 2 summarizes the values of the asymptotic constants obtained by substituting the asymptotic behaviour of the collisionless mass in the non-dimensional Euler equation. The dotted and dashed curves in figure 1 represent the density exponent vs $`s`$ computed from table 1 and table 2, respectively, for $`\gamma =5/3`$. The solid curve is the collisionless matter exponent $`\stackrel{~}{\delta }`$ versus $`s`$.
In spherical perturbations we see from the tables and figure 1 that the asymptotic gas density profile in a universe with collisionless matter is steeper than that containing gas only. For $`s>2`$ the density asymptotic exponents of the gas and collisionless matter are identical. For $`s<2`$ the collisionless matter density profile is steeper. This is consistent with neglecting the gravity of the gas near the center. As $`\gamma 4/3`$ gas density exponent approaches $`2`$ for $`s<2`$ and $`3s/(1+s)`$ otherwise, so the gas has the same asymptotic profile as the collisionless matter for all $`s`$.
In cylindrical perturbations with $`s`$ large enough, the gas has a steeper density profile than the collisionless matter. This means that virial motions in the collisionless matter are more effective at balancing gravity than the gas pressure force. So for large $`s`$ the derivation of the asymptotic exponent in section 3 is more suitable. A steeper gas profile also occurs in planar perturbations with $`s<1/(3\gamma )`$. However the asymptotic constants obtained here and in section 3 are identical in planar geometry.
In spherical perturbations with $`s>2`$, the density $`\rho (r,t)`$ is constant with time and the temperature $`p/\rho `$ diverges like $`r^{\frac{2s}{1+s}}`$ as $`r0`$. To first order in the asymptotic expansion where the gas can be in hydrostatic equilibrium in the dominant potential well of the collisionless matter. For $`s<2`$ the density increases with time and the temperature is constant near the center. For comparison, in the absence of collisionless matter, for $`2>s>1/2`$ the temperature decreases inward but the gas is in hydrostatic equilibrium to first order in the asymptotic expansion.
## 6 discussion
The similarity solutions are found for collapse in a flat universe with matter density parameter $`\mathrm{\Omega }=1`$. Because of Birchhoff’s theorem, the solutions for spherical collapse are valid in an open universe if the current turnaround radius is well inside the spherical region interior to which the perturbation is bound. The statement is incorrect for planar and cylindrical perturbations because of the explicit appearance of cosmology dependent terms in the equations of motion (14), like $`t^{2(3n)/3}`$ in the continuity equation (1).
The solutions are appropriate for the adiabatic collapse of perturbations with deep gravitational potential so that the initial thermal energy of the gas can be ignored. Such perturbations are probably the seeds for massive galaxies, galaxy groups and clusters. In the intergalactic medium (IGM) most of the gas is continuously photo-heated and is of moderate density. There is considerable interest in analytic modelling of the IGM in current methods for extracting cosmological information from the Lyman forest (Croft et. al. 1998, Nusser & Haehnelt 1999, 2000). So far these methods have heavily relied on linear analysis (e.g., Bi, Börner, & Chu 1992, Gnedin & Hui 1998, Nusser 2000) and hydrodynamical simulations (e.g., Petitjean et. al. 1995, Theuns et. al. 1999). Analytic treatment of the IGM beyond the linear regime is exceedingly complicated. Consider a situation in which photo-heating establishes the relation $`p=k\rho ^\gamma `$ in the IGM, where $`k`$ and $`\gamma `$ depend non-trivially on time (e.g., Theuns et. al. 1999). The pressure in this case introduces a length scale $`k^{1/2}G^{(1\gamma )/2}t^{2\gamma }`$ (e.g., Sedov 1959). If we take constant $`k`$ and $`\gamma `$, this length scale varies with time like $`r_{ta}`$ only in the special case of $`\gamma =4/3`$ and infinite $`s`$. So physically interesting situations including initial pressure in which the collapse is self-similar do not exist.
Spherical Perturbations with $`s>2`$ when $`\gamma >4/3`$, and with any $`s>0`$ when $`\gamma =4/3`$ deserve special attention. In the corresponding solutions for shocked collapse without a black hole at the center, the quantity $`V_c^2=2Gm/r`$ diverges towards the center, and so there is a point $`r=r_g`$ at which $`V_c^2=c^2`$, where $`c`$ is the speed of light. According to general relativity this implies the presence of a black hole at the center, invalidating the assumption of no central mass, made in deriving the solutions. In particular, the condition $`\upsilon (r=0)=0`$, necessary for shocked collapse, is incompatible with the presence of a central black hole. Near $`r_g`$, however, radiation pressure and angular momentum can prevent the formation of a black hole. Should this occur, we expect our solution for shocked accretion to be valid away from the central region.
The evolved gas variables in the symmetric self-similar collapse contain full information on the initial perturbation. So the system retains memory of the initial conditions, even in the highly nonlinear regime. On the other hand, a collapsing system of collisionless matter can develop density profiles which do not depend on the initial shape of the perturbation. For example, according to the solutions of Fillmore & Goldreich (1984), a spherical density perturbation develops into $`r^2`$ for $`s<2`$, and a cylindrical perturbation into $`r^1`$ for all $`s`$. Haloes identified in cosmological simulations of collisionless particles with generic initial conditions, also tend to have density profiles independent of the spectrum of the initial fluctuations (Navarro, Frenk & White 1997).
Our results are relevant for describing the gas distribution in various physical systems such as the cores of clusters or pancake-like superclusters. Over a limited range of scales, the index $`s`$ can be related to the index, $`l`$, of the three dimensional power spectrum, $`p(k)k^l`$, of the linear density fluctuations. If the initial density field is gaussian with a scale free power spectrum then the properties of the nonlinear field depend only on one scale. This is the nonlinear scale, $`R_{nl}`$, defined as the scale on which the rms value of density fluctuations is unity. This scale grows with time like<sup>1</sup><sup>1</sup>1$`R_{nl}`$ does not involve the dimension $`n`$ because $`l`$ refers to the three dimensional $`p(k)`$ so the rms value on a scale $`R`$ is $`R^{(l+3)/2}`$ independent of $`n`$. $`R_{nl}t^{\frac{2(l+5)}{3(l+3)}}`$. By matching the time dependence of $`R_{nl}`$ and $`r_{ta}t^{\frac{2(s+1)}{3s}}`$ we identify $`s=(l+3)/2`$. So the collapse of gas into galaxies and clusters can be, respectively, modeled by our solutions for $`s0.4`$ and $`0.65`$, where we have taken $`l2.2`$ and $`1.7`$ assuming a Cold Dark Matter power spectrum. Taking $`l1`$ for collapse on a pancake-like large scale superclusters gives $`s1`$. Another way to relate $`s`$ and $`l`$ is to identify symmetric perturbations with local maxima in the linear density field (Hoffman & Shaham 1985). The shape of high density peaks in a gaussian field varies with $`r`$ like the two-point correlation function, $`r^{(l+3)}`$ (e.g., Bardeen et. al. 1986). So, at least in the limit of high peaks, $`s=l+3`$. On galaxy and cluster scales the relation $`s=l+3`$, respectively, gives $`s=0.8`$ and $`1.3`$ in contrast to $`s=0.4`$ and $`0.65`$ obtained from $`s=(l+3)/2`$. Note however that gas in galaxies tends to settle into disks and so our solutions are less relevant than they are for cluster size objects.
In spherical geometry the asymptotic behaviour shows that the gas cannot be pressure supported if $`s<1/2`$, and $`s<2`$ for collapse with, and without collisionless matter, respectively. Estimates of the masses of rich galaxy clusters from X-ray observations of the intracluster gas rely on hydrostatic equilibrium (e.g., Fabian 1994). If on cluster scales $`s0.7`$$`1.3`$, then the asymptotic behaviour implies that the cluster gas may not be in hydrostatic equilibrium. How large is the error introduced in the mass estimates by assuming hydrostatic equilibrium? The following argument shows that this error is negligible. Hydrostatic equilibrium calculations neglect the term $`G^1r^2\mathrm{d}\upsilon /\mathrm{d}t`$ in the mass estimate. Using the asymptotic expansion one finds that neglecting this term amounts to a relative mass error of $`(2/\pi ^2)(t/t_{ta})^{3V_02}`$ where $`t_{ta}`$ is the turnaround time of the shell present at $`r`$ at the current time $`t`$. Shells in the inner regions have passed their maximum expansion a few dynamical times ago. Therefore $`tt_{ta}`$ and since $`V_0<0`$ we conclude that the error is negligible.
The solutions are related to modelling the structure of haloes made of self interacting dark matter (SIDM) (Spergel & Steinhardt 1999) with large interaction cross section. On scales of massive galaxies and clusters, our results predict final density profile $`\rho r^{1.2}`$ to $`r^{1.7}`$. These profiles are consistent with the results obtained by Moore et. al. in their simulations of SIDM with large cross section.
## 7 acknowledgement
We thank the referee for useful comments. This research was supported by The Israel Science Foundation founded by The Israel Academy of Sciences and Humanities, and by the Fund for the Promotion of Research at The Technion.
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# On the center of affine Hecke algebras of type A
### 0.1 Introduction.
Let $`G`$ be a simple complex algebraic group. Let $`W`$ be its Weyl group and $`\widehat{W}`$ the associated extended affine Weyl group. Let $`\widehat{𝐇}`$ be the Iwahori-Hecke algebra of $`\widehat{W}`$. It is well-known that $`\widehat{𝐇}`$ admits two presentations : the Coxeter presentation which arises naturally when $`\widehat{𝐇}`$ is realized as the convolution algebra $`L(\widehat{G},I)`$ of compactly supported functions on a p-adic group $`\widehat{G}=G(\overline{_p})`$ which are bi-invariant under action of the Iwahori subgroup $`I`$ (see \[IM\]), and the Bernstein presentation, which arises when $`\widehat{𝐇}`$ is realized in the $`G^{}\times ^{}`$-equivariant K-theory of the Steinberg variety associated to $`G^{}`$ where $`G^{}`$ is the Langlands dual group (see \[Gi\]). The interplay between these two presentations is central in the Deligne-Langlands correspondence for finite-dimensional irreducible representations of $`\widehat{𝐇}`$.
The center $`Z(\widehat{𝐇})`$ is easily described in the K-theoretic picture : it is spanned by the classes of the trivial (equivariant) bundles on $`Z`$. A geometric construction of this center in the convolution algebra presentation is given by Gaitsgory, \[Ga\]. This is in turn inspired by work of Beilinson, and Haines, Kottwitz and Rapoport in the framework of Shimura varieties, see \[H1\],\[H2\].
###
In this paper we give an explicit expression for the central elements of $`\widehat{𝐇}`$ in the Coxeter presentation when $`G=GL(r)`$ (Theorem 2.5). This expression generalizes those obtained by Haines in the minuscule case, \[H2\] and is in some sense more explicit than \[Ga\]. More generally, we obtain expressions for central elements in the “parabolic spherical” Hecke algebras $`L(\widehat{G},P)`$ where $`PI`$ is a parahoric subgroup. In particular, taking $`P=K`$ to be a maximal compact open subgroup recovers Lusztig’s description \[L1\] of the Satake isomorphism between $`Z(\widehat{𝐇})`$ and the spherical algebra $`\widehat{𝐇}_{sph}`$ (in the case $`G=GL(r)`$).
###
Our method is based on the Hall algebra of a cyclic quiver, on Uglov’s higher-level Fock spaces and on the theory of canonical bases of Kashiwara and Lusztig. Namely, we use Ginzburg and Vasserot’s geometric description of quantum affine Schur-Weyl duality to construct an embedding of (half of) the center $`Z(\widehat{𝐇})`$ in the center of the Hall algebra $`𝐔_n^{}`$ of the quiver $`\stackrel{~}{A}_{n1}`$ for $`nr`$ (see \[S\]). This embedding is compatible with the canonical bases of $`\widehat{𝐇}`$ and $`𝐔_n^{}`$. To describe the center of $`𝐔_n^{}`$ we then consider the action on the Fock spaces $`\mathrm{\Lambda }_{𝐬_l}^{\mathrm{}}`$ recently introduced by Uglov \[U\], and use the fact that this action is again compatible with the canonical bases.
###
Finally, we give a simple alternate description of the center of $`𝐔_n^{}`$ in terms of a certain desingularization of orbit closures of representations of the quiver $`\stackrel{~}{A}_{n1}`$, introduced by Varagnolo and Vasserot \[VV\]. This can be seen as a cyclic analogue of the desingularization of orbit closures recently obtained by Reineke \[Re\] for finite-type simply laced Dynkin quivers.
###
We note that the Fock spaces and their canonical bases appear to be a very fundamental object in type A representation theory : they describe Grothendieck groups and decomposition numbers of Hecke algebras of type A or B (or more generally cyclotomic Hecke algebras) at roots of unity (see \[LLT\],\[A\], \[AM\], \[Gro\]), and modular representations of symmetric groups (see \[Di\], \[J\], \[Gro\]).
### 0.2 Notations.
Set $`𝕊=[v]`$, $`𝔸=[v,v^1]`$ and $`𝕂=(v)`$. We define a $``$-linear ring involution $`u\overline{u}`$ on $`𝔸`$ by setting $`\overline{v}=v^1`$. Let $`𝔽`$ be a finite field with $`q^2`$ elements. Let $`𝔖_r`$ denote the symmetric group on $`r`$ elements and let $`\{s_i\}_{i=1,\mathrm{}r1}`$ be the set of simple reflections. Let $`\widehat{𝔖}_r=𝔖_r^r`$ be the extended affine symmetric group and let $`s_0`$ be the affine simple reflection. Let $`\mathrm{\Pi }`$ stand for the set of partitions and let $`\mathrm{\Pi }_r`$ be the set of partitions of length at most $`r`$. Elements of $`\mathrm{\Pi }^l`$ for some $`l`$ will be called $`l`$-multipartitions. Finally, we will denote by $`\overline{Y}`$ the Zariski closure of any subset $`Y`$ of an algebraic variety $`X`$.
## 1 Affine Hecke algebras and canonical bases
### 1.1
Consider the Iwahori-Hecke algebra $`\widehat{𝐇}_r`$ associated to $`\widehat{𝔖}_r`$, i.e the $`𝔸`$-algebra generated by elements $`T_\sigma `$, $`\sigma \widehat{𝔖}_r`$ with relations
$$(T_{s_i}+1)(T_{s_i}v^2)=0fori=0,\mathrm{},r1,$$
$$T_\sigma T_\gamma =T_{\sigma \gamma }ifl(\sigma \gamma )=l(\sigma )l(\gamma ).$$
We set $`\stackrel{~}{T}_\sigma =v^{l(\sigma )}T_\sigma `$ for every $`\sigma \widehat{𝔖}_r`$.
###
It is well-known that $`\widehat{𝐇}_r`$ admits another presentation (the Bernstein presentation) as the unital $`𝔸`$-algebra generated by elements $`T_i^{\pm 1},X_j^{\pm 1}`$ where $`i[1,r1]`$, $`j[1,r]`$ with the following relations
$`T_iT_i^1=1=T_i^1T_i,`$ $`(T_i+1)(T_iv^2)=0,`$
$`T_iT_{i+1}T_i=T_{i+1}T_iT_{i+1},`$ $`|ij|>1T_iT_j=T_jT_i,`$
$`X_iX_i^1=1=X_i^1X_i,`$ $`X_iX_j=X_jX_i,`$
$`T_iX_iT_i=v^2X_{i+1},`$ $`ji,i+1X_jT_i=T_iX_j.`$
The isomorphism between the two presentations is such that $`T_{s_i}T_i`$ and $`\stackrel{~}{T}_\lambda ^1X_1^{\lambda _1}\mathrm{}X_r^{\lambda _r}`$ if $`\lambda =(\lambda _1,\mathrm{},\lambda _r)`$ is dominant. The center of $`\widehat{𝐇}_r`$ is $`Z(\widehat{𝐇}_r)=𝔸[X_1^{\pm 1},\mathrm{},X_r^{\pm 1}]^{𝔖_r}`$. Set $`Z_r^{}=𝔸[X_1^1,\mathrm{},X_r^1]^{𝔖_r}.`$
### 1.2
For every $`t,s`$ define the left (resp. right) representation of $`\widehat{𝔖}_t`$ on $`^t`$ of level $`s`$ by
$`s_j(i_1,\mathrm{},i_t)`$ $`=(i_1,\mathrm{},i_{j+1},i_j,\mathrm{}i_t),1j<r,`$
$`\lambda (i_1,\mathrm{},i_t)`$ $`=(i_1+s\lambda _1,\mathrm{},i_t+s\lambda _t),\lambda ^t`$
and
$`(i_1,\mathrm{},i_t)s_j`$ $`=(i_1,\mathrm{},i_{j+1},i_j,\mathrm{}i_t),1j<r,`$
$`(i_1,\mathrm{},i_t)\lambda `$ $`=(i_1+s\lambda _1,\mathrm{},i_t+s\lambda _t),\lambda ^t`$
respectively. The set $`𝒜_t^s=\{1i_1\mathrm{}i_ts\}`$ is a fundamental domain for both actions. For each $`𝐢𝒜_t^s`$ we set $`𝔖_𝐢=Stab𝐢𝔖_t`$ and denote by $`\omega _𝐢𝔖_𝐢`$ the longest element. We also let $`𝔖^𝐢`$ be the set of all minimal length elements of the cosets $`𝔖_𝐢\backslash \widehat{𝔖}_t`$.
### 1.3
Fix some $`n^{}`$. For any $`𝐢,𝐣𝒜_r^n`$ and any $`\sigma 𝔖_𝐢\backslash \widehat{𝔖}_r/𝔖_𝐣`$ we set $`T_\sigma =_{\delta \sigma }T_\delta `$ and we let $`\widehat{𝐇}_{\mathrm{𝐢𝐣}}\widehat{𝐇}_r`$ be the $`𝔸`$-linear span of the elements $`T_\sigma `$ for $`\sigma 𝔖_𝐢\backslash \widehat{𝔖}_r/𝔖_𝐣`$. Set $`e_𝐢=_{\delta 𝔖_𝐢}T_\delta `$. Then $`\widehat{𝐇}_{\mathrm{𝐢𝐣}}=e_𝐢\widehat{𝐇}_re_𝐣`$. Put
$$\widehat{𝐒}_{n,r}=\underset{𝐢,𝐣𝒜_r^n}{}\widehat{𝐇}_{\mathrm{𝐢𝐣}}.$$
This space, equipped with the multiplication
$$e_𝐢he_𝐣e_𝐤h^{}e_𝐥=\delta _{jk}e_𝐢he_𝐣h^{}e_𝐥\widehat{𝐇}_{\mathrm{𝐢𝐥}}forallh,h^{}\widehat{𝐇}_r$$
is called the affine q-Schur algebra. It is proved in \[GV\], \[L3\] that $`\widehat{𝐒}_{n,r}`$ is a quotient of the modified quantum affine algebra $`\dot{𝐔}_v^{}(\widehat{𝔤𝔩}_n)`$.
### 1.3
Set $`𝐓_{n,r}=_{𝐢𝒜_r^n}e_𝐢\widehat{𝐇}_r`$. For $`\sigma 𝔖_𝐢\backslash \widehat{𝔖}_r`$ we put $`T_\sigma =_{\delta \sigma }T_\delta `$. Then $`\{𝐓_\sigma \},\sigma 𝔖_𝐢\backslash \widehat{𝔖}_r`$ is an $`𝔸`$-basis of $`e_𝐢\widehat{𝐇}_r`$. It will be convenient to identify the element $`\sigma `$ with $`𝐢\sigma ^r`$, so that $`\{𝐓_p\},p^r`$ is an $`𝔸`$-basis of $`𝐓_{n,r}`$.
###
The algebra $`\widehat{𝐇}_r`$ acts on $`𝐓_{n,r}`$ by multiplication on the right, and $`\widehat{𝐒}_{n,r}`$ acts on $`𝐓_{n,r}`$ on the left by
$$e_𝐢he_𝐣e_𝐤h^{}=\delta _{\mathrm{𝐣𝐤}}e_𝐢he_𝐣h^{}e_𝐢\widehat{𝐇}_rforeveryh,h^{}\widehat{𝐇}_r.$$
Let us denote these actions by $`\rho _r:\widehat{𝐒}_{n,r}\mathrm{End}(𝐓_{n,r})`$ and $`\sigma _r:\widehat{𝐇}_r\mathrm{End}(𝐓_{n,r})`$. It is obvious that these two actions commute. The following result is a quantum and affine analogue of Schur-Weyl duality.
###### Theorem (\[VV\]).
We have $`\widehat{𝐒}_{n,r}=\mathrm{End}_{\widehat{𝐇}_r}(𝐓_{n,r})`$. Moreover, we have $`\widehat{𝐇}_r=\mathrm{End}_{\widehat{𝐒}_{n,r}}(𝐓_{n,r})`$ if $`nr`$.
### 1.4
Let us now, following \[GV\] and \[IM\], give the geometric realization of the above Schur-Weyl duality. Let $`𝕃=𝔽((z))`$ and set $`𝔾=GL_r(𝕃)`$. By definition, a lattice in $`𝕃^r`$ is a free $`𝔽[[z]]`$-submodule of rank $`r`$. Consider the variety $`X`$ of sequences of lattices $`(L_i)_i`$ such that
$$L_iL_{i+1},\mathrm{dim}_𝔽(L_i/L_{i1})=1,L_{i+r}=z^1L_i$$
(the affine flag variety of type $`GL_r`$). Consider also the variety $`Y`$ of all $`n`$-step periodic flags in $`𝕃^r`$, i.e the set of all sequences of lattices $`(L_i)_i`$ such that
$$L_iL_{i+1},L_{i+n}=z^1L_i$$
(the affine partial flags variety). The group $`𝔾`$ acts (transitively) on $`X`$ and acts on $`Y`$ in obvious ways. Consider the diagonal action of $`𝔾`$ on $`X\times X`$ and $`Y\times Y`$ respectively.
###
It is well-known that the set of $`𝔾`$-orbits on $`X\times X`$ is canonically identified with $`\widehat{𝔖}_r`$. In order to describe these $`𝔾`$-orbits we let $`(e_1,\mathrm{},e_r)`$ be a fixed $`𝕃`$-basis of $`𝕃^r`$ and set $`e_{i+kr}=z^ke_i`$. Consider the right action of $`\widehat{𝔖}_r`$ on $`^r`$ of level $`r`$. To any element $`𝐱`$ in the orbit of $`\rho _r=(1,2,\mathrm{},r)`$ we associate the flag $`(L(𝐱)_i)_i`$ defined by
$$L(𝐱)_i=\underset{p(j)i}{}𝔽e_j,$$
where $`p:`$ is the bijection uniquely defined by $`p(j)=𝐱_j`$ if $`1jr`$ and $`p(j+r)=p(j)+r`$. The $`𝔾`$-orbit decomposition of $`X\times X`$ reads
$$X\times X=\underset{\sigma \widehat{𝔖}_r}{}X_\sigma $$
where $`X_\sigma =𝔾(L(\rho _r\sigma ),L(\rho _r))`$. Similarly, to each $`𝐢^r`$ we associate the map $`p:`$ uniquely defined by $`p(j)=𝐢_j`$ if $`1jr`$ and $`p(j+r)=p(j)+n`$. Consider the flag
$$L(𝐢)_i=\underset{𝐢(j)i}{}𝔽e_j.$$
Then $`Y=_{𝐢𝒜_r^n}Y_𝐢`$ where $`Y_𝐢=𝔾(L(𝐢))`$ and
$$Y_𝐢\times Y_𝐣=\underset{\sigma 𝔖_𝐢\backslash \widehat{𝔖}_r/𝔖_𝐣}{}Y_\sigma $$
where $`Y_\sigma =𝔾(L(𝐢\sigma ),L(𝐣))`$ and where the right action of $`\widehat{𝔖}_r`$ on $`^r`$ is now of level $`n`$.
###
Let $`_𝔾(X\times X)`$ (resp. $`_𝔾(Y\times Y)`$) be the space of complex-valued $`𝔾`$-invariant functions on $`X\times X`$ (resp. on $`Y\times Y`$) which are supported on finitely many orbits. The convolution product endows these spaces with an associative algebra structure. We let $`\mathrm{𝟏}_𝒪_𝔾(X\times X)`$ (resp. $`\mathrm{𝟏}_𝒪_𝔾(Y\times Y)`$) be the characteristic function of a $`𝔾`$-orbit $`𝒪X\times X`$(resp. $`𝒪Y\times Y`$).
###### Theorem (\[IM\],\[VV\]).
1. The linear map $`(\widehat{𝐇}_r)_{|v=q^1}_𝔾(X\times X)`$ defined by $`T_\sigma \mathrm{𝟏}_{X_\sigma }`$ is an algebra isomorphism.
2. The linear map $`(\widehat{S}_{n,r})_{|v=q^1}_𝔾(Y\times Y)`$ such that $`T_\sigma \mathrm{𝟏}_{Y_\sigma }`$ is an algebra isomorphism.
###
Now consider the diagonal action of $`𝔾`$ on $`Y\times X`$. The collection of orbits are parametrized by $`^r`$: to $`𝐢^r`$ corresponds the orbit $`𝒪_𝐢`$ of the pair $`(L(𝐢),L(\rho _r))`$. The algebras $`_𝔾(X\times X)`$ and $`_𝔾(Y\times Y)`$ act by convolution on $`_𝔾(Y\times X)`$ on the right and on the left respectively.
###### Theorem (\[VV\]).
The map $`(𝐓_{n,r})_{|v=q^1}_𝔾(Y\times X)`$ such that $`e_𝐢\mathrm{𝟏}_{𝒪_𝐢}`$ for $`𝐢𝒜_r^n`$ extends uniquely to an isomorphism of $`(\widehat{𝐒}_{n,r})_{|v=q^1}\times (\widehat{𝐇}_r)_{|v=q^1}`$-modules.
### 1.5
Let $`u\overline{u}`$ be the semilinear involution of $`\widehat{𝐇}_r`$ defined by $`\overline{T}_\sigma =T_{\sigma ^1}^1`$ for all $`\sigma `$. For each $`\sigma \widehat{𝔖}_r`$ there exists a unique element $`𝐜_\sigma \widehat{𝐇}_r`$ such that
$$\mathrm{i})\overline{𝐜_\sigma }=𝐜_\sigma ,\mathrm{ii})𝐜_\sigma =\stackrel{~}{T}_\sigma +_{\delta <\sigma }c_{\delta ,\sigma }(v)\stackrel{~}{T}_\delta ,c_{\delta ,\sigma }(v)v𝕊.$$
The polynomial $`c_{\sigma ,\delta }(v)`$ is the affine Kazhdan-Lusztig polynomial of type $`\stackrel{~}{A}_{r1}`$ associated to $`\sigma `$ and $`\delta `$ (this polynomial is denoted by $`h_{\sigma ,\delta }`$ in Soergel’s notation \[Soe\]).
###
For $`\sigma \widehat{𝔖}_r`$ and $`LX`$ let $`X_{\sigma ,L}`$ be the fiber of the first projection $`X_\sigma X`$. Then $`X_{\sigma ,L}`$ is the set of $`𝔽`$-points of an algebraic variety of dimension $`l(\sigma )`$ whose isomorphism class is independent of $`L`$. Then
$$𝐜_\sigma =\underset{i,\delta }{}v^{i+l(\sigma )l(\delta )}\mathrm{dim}_{X_{\delta ,L}}^i(IC_{X_{\sigma ,L}})\stackrel{~}{T}_\delta $$
where $`IC_{X_{\sigma ,L}}`$ denotes the intersection cohomology complex associated to $`X_{\sigma ,L}`$ and where $`^i`$ stands for local cohomology.
###
Similarly, let $`𝐢,𝐣𝒜_r^n`$ and let $`\sigma 𝔖_𝐢\backslash \widehat{𝔖}_r/𝔖_𝐣`$. Denote by $`Y_{\sigma ,𝐢}`$ the fiber above $`(L(𝐢))`$ of the projection of $`Y_\sigma Y`$ on the first component. This is the set of $`𝔽`$-points of an algebraic variety of dimension, say $`y(\sigma )`$ (an explicit formula for $`y(\sigma )`$ can be found in \[L3\]). Put $`\stackrel{~}{T}_\sigma =v^{y(\sigma )}T_\sigma `$. For every $`\sigma 𝔖_𝐢\backslash \widehat{𝔖}_r/𝔖_𝐣`$ set
$$𝐜_\sigma =\underset{i,\delta }{}v^{i+y(\sigma )y(\delta )}\mathrm{dim}_{Y_{\delta ,𝐢}}^i(IC_{Y_{\sigma ,𝐢}})\stackrel{~}{T}_\delta .$$
It is clear that $`\overline{\widehat{𝐇}_{\mathrm{𝐢𝐣}}}=\widehat{𝐇}_{\mathrm{𝐢𝐣}}`$. Define a semilinear involution $`\tau :\widehat{𝐇}_{\mathrm{𝐢𝐣}}\widehat{𝐇}_{\mathrm{𝐢𝐣}}`$ by $`\tau (u)=v^{2l(\omega _𝐣)}\overline{u}`$. The elements $`\{𝐜_\sigma \}`$ for all $`𝐢,𝐣𝒜_r^n`$ form the canonical basis of $`\widehat{𝐒}_{n,r}`$ and are characterized by the following two properties :
$$\mathrm{i})\tau (𝐜_\sigma )=𝐜_\sigma ,\mathrm{ii})𝐜_\sigma =\stackrel{~}{T}_\sigma +_{\delta <\sigma }c_{\delta ,\sigma }(v)\stackrel{~}{T}_\delta ,c_{\delta ,\sigma }(v)v𝕊.$$
### 1.6
Let $`s,t^{}`$. For $`𝐢𝒜_t^s`$ and $`x𝐢\widehat{𝔖}_t`$ set $`x|=e_𝐢\stackrel{~}{T}_a`$ where $`𝐢a=x`$ and $`a𝔖^𝐢`$. The set $`\{x|,x𝐢\widehat{𝔖}_t\}`$ is an $`𝔸`$-basis of the space $`e_𝐢\widehat{𝐇}_t`$. Define a semilinear involution $`u\overline{u}`$ of $`e_𝐢\widehat{𝐇}_t`$ by $`\overline{e_𝐢x}=e_𝐢\overline{x}`$. There exists a unique $`𝔸`$-basis $`\{𝐜_x^{},x𝐢\widehat{𝔖}_t\}`$ of $`e_𝐢\widehat{𝐇}_t`$ such that
$$\mathrm{i})\overline{𝐜_x^{}}=𝐜_x^{},\mathrm{ii})𝐜^{}_x=x|+\underset{y}{}P_{y,x}^{}y|,P_{y,x}^{}v^1[v^1].$$
The polynomials $`P_{y,x}^{}`$ are parabolic affine Kazhdan-Lusztig polynomials introduced by Deodhar \[De\]. These polynomials are (up to a sign) denoted by $`\overline{n}_{a_y,a_x}`$ in Soergel’s notation, where $`a_x,a_y𝔖^𝐢`$ are such that $`x=𝐢a_x`$, $`y=𝐢a_y`$.
## 2 The main result
### 2.1
Let $`\mathrm{\Gamma }`$ be Macdonald’s ring of symmetric polynomial in the variables $`y_i`$, $`i`$, defined over $`𝔸`$ (see \[Mac\]). Let $`\mathrm{\Gamma }_r=𝔸[y_1,\mathrm{},y_r]^{𝔖_r}`$. Let $`s_\lambda \mathrm{\Gamma }_r`$ be the Schur polynomial associated to $`\lambda \mathrm{\Pi }_r`$.
###
Fix some $`n`$ and let $`𝐢𝒜_r^n`$. From $`s_\lambda (X_1^1,\mathrm{},X_r^1)Z(\widehat{𝐇}_r)`$ it follows that $`e_𝐢s_\lambda (X_1^1,\mathrm{},X_r^1)\widehat{𝐇}_{\mathrm{𝐢𝐢}}`$. Define polynomials $`J_{\lambda ,\sigma }^𝐢[v,v^1]`$ by the relation
$$e_𝐢s_\lambda (X_1^1,\mathrm{},X_r^1)=(v)^{(n1)|\lambda |}\underset{\sigma 𝔖_𝐢\backslash \widehat{𝔖}_r/𝔖_𝐢}{}J_{\lambda ,\sigma }^𝐢𝐜_\sigma .$$
In this section we give an explicit expression for $`J_{\lambda ,\sigma }^𝐢`$ involving (parabolic) affine Kazhdan-Lusztig polynomials of type $`A`$.
### Remark.
It is clear that (up to a power of $`v`$) $`J_{\lambda ,\sigma }^𝐢`$ depends only on $`𝔖_𝐢`$ rather than on $`𝐢`$. In particular, any parabolic subgroup $`𝔖_{i_1}\times \mathrm{}\times 𝔖_{i_t}`$ occurs as $`𝔖_𝐢`$ for some $`𝐢𝒜_r^n`$ as soon as $`nt`$.
### 2.2
We first make some preliminary definitions. We will represent a partition $`\lambda `$ by its associated Young diagram in the usual fashion. We will consider diagrams where the $`(i,j)`$-box has content $`ij+r_0modn`$ for some fixed $`r_0/n`$ and call the resulting tableau the partition $`\lambda `$ with residue $`r_0`$. We will say that a box with content $`j/n`$ can be added to the partition $`\lambda `$ with residue $`r_0`$ if there exists a partition $`\lambda ^{}`$ with residue $`r_0`$ such that $`\lambda ^{}/\lambda `$ is a single box with content $`j`$. For example, when $`n=3`$, the partition $`\lambda =(421)`$ with residue $`1`$ is
and the dotted lines correspond to addable boxes.
###
To each $`𝐩(^+)^r`$ and $`𝐢𝒜_r^n`$ we associate a multipartition (with residues) $`_𝐢(𝐩)`$. First, we attach a diagram (not a partition!)
$$D_𝐩=\{(i,j)|\mathrm{\hspace{0.33em}0}<j𝐩_i\}/r\times ^+,$$
where we fill the $`(i,j)`$-box with the content $`𝐢_i+𝐩_ijmodn`$.
###
Example 1. Suppose $`r=n=5`$, $`𝐢=(1,2,3,4,5)`$ and $`𝐩=(4,3,4,3,5)`$. Then $`D_𝐩`$ is
Now consider the horizontal slices $`s_k=D_𝐩(/r\times \{k\})`$ and let $`k_0`$ be maximal such that $`s_{k_0}\mathrm{}`$. We construct the multipartition with residues $`_𝐢(𝐩)`$ by successively adding the boxes from $`s_{k_0},\mathrm{},s_1`$ in the following way. Set $`^{k_0+1}=\mathrm{}`$. Suppose $`^i=(\lambda _i^{(1)},\mathrm{},\lambda _i^{(t)})`$ is known. Then $`^{i1}=(\lambda _{i1}^{(1)},\mathrm{},\lambda _{i1}^{(r)})`$ is obtained from $`^i`$ by adding the boxes from $`s_i`$ (possibly creating new partitions) in such a way that
1. For every $`1vr`$, $`\lambda _{i1}^{(v)}/\lambda _i^{(v)}`$ is a skew tableau with at most one box in each row,
2. $`^{i1}`$ is maximal for the following order :
$`(\lambda _{i1}^{(1)},\lambda _{i1}^{(2)},\mathrm{})(\mu _{i1}^{(1)},\mu _{i1}^{(2)},\mathrm{})ifthereexistswsuchthat`$
$$\lambda _{i1}^{(l)}=\mu _{i1}^{(l)}for\mathrm{\hspace{0.33em}1}l<wand\lambda _{i1}^{(w)}\mu _{i1}^{(w)},$$
where $``$ stands for the usual dominance order of partitions,
3. If several new partitions appear in $`^{i1}`$ then they are in increasing order of their residue.
Set $`_𝐢(𝐩)=^1`$. We note that condition iii) above is not essential for the rest of the paper and here only to fix notations.
###
Examples. i) Let $`r=n`$ and $`𝐢=(1,\mathrm{},r)`$. Suppose that $`𝐩`$ is antidominant up to cyclic permutation, i.e there exists $`i/r`$ such that $`𝐩_i𝐩_{i1}\mathrm{}𝐩_{i+1}`$. Let $`\lambda `$ be the associated partition. Then $`_𝐢(𝐩)`$ consists of the single partition $`\lambda `$ with residue $`i`$.
ii) Consider $`r,n`$, $`𝐢`$ and $`𝐩`$ as in example 1. Then the algorithm for computing $`_𝐢(𝐩)`$ runs as follows :
For $`\sigma 𝔖_𝐢\backslash \widehat{𝔖}_r/𝔖_𝐢`$ we let $`\sigma _0`$ be the longest element in $`\sigma `$ and we set $`𝐢\sigma =𝐢\sigma _0𝐢^r`$. For each $`\sigma `$ such that $`𝐢\sigma (^+)^r`$ we set $`(\sigma )=_𝐢(𝐢\sigma )`$. Write $`(\sigma )=(\sigma ^{(1)},\mathrm{},\sigma ^{(l)})`$ where $`\sigma ^{(l)}\mathrm{}`$, and $`𝐫_\sigma =(r_1,\mathrm{},r_l)`$ where $`r_i/n`$ is the residue of $`\sigma ^{(i)}`$.
### 2.3
Let $`l`$. Let $`(\sigma ^{(1)},\mathrm{},\sigma ^{(l)})`$, $`(\mu ^{(1)},\mathrm{},\mu ^{(l)})`$ be any $`l`$-multipartitions and let $`𝐫=(r_1,\mathrm{},r_l)(/n)^l`$. Choose some $`𝐬=(s_1,\mathrm{},s_l)^l`$ such that $`s_ir_i(\mathrm{mod}n)`$. For $`i=1,\mathrm{}l`$ and $`j`$ we set $`u_j^{(i)}=s_i+\sigma _j^{(i)}+1j`$ and $`v_j^{(i)}=s_i+\mu _j^{(i)}+1j`$. Consider, for $`t0`$
$$𝐮=(u_1^{(1)},\mathrm{},u_t^{(1)},u_1^{(2)},\mathrm{},u_t^{(2)},\mathrm{},u_t^{(l)}),$$
$$𝐯=(v_1^{(1)},\mathrm{},v_t^{(1)},v_1^{(2)},\mathrm{},v_t^{(2)},\mathrm{},v_t^{(l)}).$$
Finally, we put
$$𝐏_{(\mu ^{(1)},\mathrm{},\mu ^{(l)}),(\sigma ^{(1)},\mathrm{},\sigma ^{(l)})}^{,𝐬}=P_{𝐯,𝐮}^{}.$$
Now let $`𝐬`$ be in the asymptotic range $`s_1s_2\mathrm{}s_l`$ and set
$$𝐏_{(\mu ^{(1)},\mathrm{},\mu ^{(l)}),(\sigma ^{(1)},\mathrm{},\sigma ^{(l)})}^{,𝐫}=P_{𝐯,𝐮}^{}.$$
This polynomial is independent of the choices of $`𝐬`$ and $`t`$ in the given asymptotic range (this follows for instance from \[U\] Section 4 and \[S\], Theorem 4.1). These can be thought of as some “stabilization” of polynomials $`P_{\mu +\rho _r,\sigma +\rho _r}^{}`$ of type $`\stackrel{~}{A}_r`$ as $`r`$ tends to infinity (see \[LT\]). Moreover, it is easy to see that when $`l=1`$, $`𝐏^{,𝐫}`$ is independent of $`𝐫`$ and we will omit it.
### 2.4
For any multipartition $`\mu =(\mu ^{(1)},\mathrm{},\mu ^{(l)})`$ and $`𝐫=(r_1,\mathrm{},r_l)(/n)^l`$ we set $`\mu ^{}=((\mu ^{(l)})^{},\mathrm{},(\mu ^{(1)})^{})`$ and $`𝐫^{}=(r_l,\mathrm{},r_1)`$.
### 2.5
The following is the main result of this paper, and will be proved in Section 5.
###### Theorem.
We have
$$e_𝐢s_\lambda (X_1^1,\mathrm{},X_r^1)=(v)^{(n1)|\lambda |}\underset{\sigma |𝐢\sigma (^+)^r}{}J_{\lambda ,\sigma }^𝐢𝐜_\sigma $$
where
$$J_{\lambda ,\sigma }^𝐢=\underset{\stackrel{\nu _1,\mathrm{},\nu _l}{\mu _1,\mathrm{},\mu _l}}{}c_{\mu _1,\mathrm{},\mu _l}^\lambda v^{{\scriptscriptstyle (b1)|\mu _b|}}𝐏_{\nu _1,n\mu _1^{}}^{}\mathrm{}𝐏_{\nu _l,n\mu _l^{}}^{}𝐏_{\nu ,(\sigma )^{}}^{,𝐫_\sigma ^{}}$$
and $`\nu =(\nu _1,\mathrm{},\nu _l)`$. Here $`c_{\mu _1,\mathrm{},\mu _l}^\lambda `$ is the (generalized) Littlewood-Richardson coefficient.
###
Examples. i) Suppose that $`n=1`$. Then $`𝐢=(1^r)`$ and $`𝔖_𝐢=𝔖_r`$. Moreover, $`𝔖_r\backslash \widehat{𝔖}_r/𝔖_r=\mathrm{\Pi }_r`$ and for $`\sigma \mathrm{\Pi }_r`$ we have $`(\sigma )=\sigma `$ and $`l=1`$. Hence the above theorem reduces to $`J_{\lambda ,\sigma }^𝐢=_\nu 𝐏_{\nu ,\lambda ^{}}^{}𝐏_{\nu ,\sigma ^{}}^{}=\delta _{\lambda ,\sigma }`$, i.e
$$\left(\underset{w𝔖_r}{}T_w\right)s_\lambda (X_1^1,\mathrm{},X_r^1)=𝐜_\lambda ,$$
in accordance with \[L1\].
ii) Let $`r=n`$ and $`𝐢=\rho `$ (i.e $`𝔖_𝐢=\{1\}`$). Let $`\lambda =(1^l)`$, $`lr`$ be a minuscule weight. Then in the above expression for $`J_{\lambda ,\sigma }^𝐢`$ the only nonzero terms correspond to the case when $`\mu _i`$ is also minuscule for all $`i`$. We obtain an expression for $`s_\lambda (X_1^1,\mathrm{},X_r^1)`$ analogous to Theorem 1.1 in \[H2\] for $`G=GL(r)`$ (but which involves Kazhdan-Lusztig polynomials rather than $`R`$-polynomials). Note that \[H1\], Proposition 5 also easily follows from the above theorem.
## 3 Hall algebra of a cyclic quiver
### 3.0 Notations.
In this section we fix a positive integer $`n`$. Let $`(ϵ_i)`$, $`i/n`$ be the canonical basis of $`^{/n}`$. For $`i/n`$ and $`l^{}`$, define the cyclic segment $`[i;l)`$ to be the image of the projection to $`/n`$ of the segment $`[i_0,i_0+l1]`$ for any $`i_0i(\mathrm{mod}n)`$. A cyclic multisegment is a linear combination $`𝐦=_{i,l}a_i^l[i;l)`$ of cyclic segments with coefficients $`a_i^l`$. Let $``$ be the set of cyclic multisegments. For $`𝐦`$ we set $`\mathrm{dim}𝐦=a_i^l(ϵ_i+\mathrm{}+ϵ_{i+l1})`$. Note that $``$ is canonically isomorphic to $`\mathrm{\Pi }^n`$ : to $`𝐦=a_i^l[i;l)`$ we associate the multipartition $`(\lambda ^{(1)},\mathrm{},\lambda ^{(n)})`$ with $`\lambda ^{(i)}=(1^{a_i^1}2^{a_i^2}\mathrm{})`$.
### 3.1
Let $`Q`$ be the quiver of type $`\stackrel{~}{A}_{n1}`$, i.e the oriented graph with vertex set $`I=/n`$ and edge set $`\mathrm{\Omega }=\{(i,i+1),iI\}`$. For any $`I`$-graded $`𝔽`$-vector space $`V=_{iI}V_i`$, let $`E_V_{(i,j)\mathrm{\Omega }}\mathrm{Hom}(V_i,V_j)`$ denote the space of nilpotent representations of $`Q`$. The group $`G_V=_{iI}GL(V_i)`$ acts on $`E_V`$ by conjugation. For each $`iI`$ there exists a unique simple $`Q`$-module $`S_i`$ of dimension $`ϵ_i`$, and for each pair $`(i,l)I\times ^{}`$ there exists a unique (up to isomorphism) indecomposable $`Q`$-module $`S_{i;l}`$ of length $`l`$ and tail $`S_i`$. Furthermore, every nilpotent $`Q`$-module $`M`$ admits an essentially unique decomposition
$$M\underset{i,l}{}a_i^lS_{i;l}.$$
(3.1)
We denote by $`\overline{𝐦}`$ the isomorphism class of $`Q`$-modules corresponding (by (3.1)) to the multisegment $`𝐦=_{i,l}a_i^l[i;l)`$. For $`𝐦`$ with $`\mathrm{dim}𝐦=𝐝`$ and $`V_𝐝`$ an $`I`$-graded vector space of dimension $`𝐝`$, we let $`O_𝐦E_{V_𝐝}`$ be the $`G_{V_𝐝}`$-orbit consisting of representations in the class $`\overline{𝐦}`$, and we let $`\mathrm{𝟏}_𝐦_G(V_𝐝)`$ be the characteristic function of $`O_𝐦`$. Finally, we set $`𝐟_𝐦=q^{\mathrm{dim}O_𝐦}\mathrm{𝟏}_𝐦`$. We will write $`𝐦<𝐧`$ if $`𝒪_𝐦\overline{𝒪_𝐧}`$.
### 3.2
Set $`𝐔_n^{}=_𝐝_G(E_{V_𝐝})`$. Note that, by definition, $`(𝐟_𝐦)_𝐦`$ is a $``$-basis of $`𝐔_n^{}`$. The space $`𝐔_n^{}`$ is endowed with the structure of a (Hall) algebra (see \[L1\]). We use the definitions of \[VV\], \[S\]. Moreover, the structure constants for this algebra are polynomials in $`q`$, and one can consider $`𝐔_n^{}`$ as an $`𝔸`$-algebra with $`q=v^1`$. The algebra $`𝐔_n^{}`$ is naturally $`^{/n}`$-graded and we denote by $`𝐔_n^{}[𝐝]`$ the component of degree $`𝐝`$. Let $`𝐔_v(\widehat{𝔰𝔩}_n)`$ denote the Lusztig integral form of the quantum affine algebra of type $`\stackrel{~}{A}_{n1}`$ and let $`e_i^{(l)},k_i,f_i^{(l)}`$, $`iI`$, $`l`$ be the divided powers of the standard Chevalley generators. Let $`𝐔_v^{}(\widehat{𝔰𝔩}_n)`$ be the subalgebra of $`𝐔_v(\widehat{𝔰𝔩}_n)`$ generated by $`f_i^{(l)},iI`$, $`l^{}`$. It is known that the map $`f_i^{(l)}𝐟_{lϵ_i}`$ extends to an embedding of the algebras $`𝐔_v^{}(\widehat{𝔰𝔩}_n)𝐔_n^{}`$.
### 3.3
For $`𝐦`$, set
$$𝐛_𝐦=\underset{i,𝐧}{}v^{i+\mathrm{dim}O_𝐦\mathrm{dim}O_𝐧}\mathrm{dim}_{O_𝐧}^i(IC_{O_𝐦})𝐟_𝐧,$$
(3.2)
where $`_{O_𝐧}^i(IC_{O_𝐦})`$ is the stalk over a point of $`O_𝐧`$ of the ith intersection cohomology sheaf of the closure $`\overline{O}_𝐦`$ of $`O_𝐦`$. Then $`𝐁=\{𝐛_𝐦\}`$ is the canonical basis of $`𝐔_n^{}`$, introduced in \[VV\].
### 3.4
Let $`L,L^{}Y`$ be two $`n`$-step periodic flags in $`𝕃^r`$ satisfying $`L^{}L`$. Following Lusztig (see \[L2\],\[GV\]) we associate to such a pair a nilpotent representation of $`\stackrel{~}{A}_{n1}`$ of graded dimension $`(\mathrm{dim}_𝔽(L_i/L_i^{}))_{\overline{i}/n}`$. Let us denote by $`L/L^{}`$ this $`\stackrel{~}{A}_{n1}`$-module. Set
$$a(L^{},L)=\underset{i=1}{\overset{n}{}}\mathrm{dim}_𝔽(L_i/L_i^{})(\mathrm{dim}_𝔽(L_{i+1}^{}/L_i^{})\mathrm{dim}_𝔽(L_i/L_i^{})).$$
Define a map $`\mathrm{\Theta }:𝐔_n^{}\widehat{𝐒}_{n,r}`$ by
$$\mathrm{\Theta }(f)(L^{},L)=q^{a(L^{},L)}f(L/L^{})\mathrm{if}L^{}L$$
and $`\mathrm{\Theta }(f)(L,L^{})=0`$ if $`L^{}L`$.
###
In order to describe $`\mathrm{\Theta }`$, we consider the following parametrization of the collection of $`𝔾`$-orbits in $`Y\times Y`$. Let $`M_{r,n}`$ be the set of $`\times `$-matrices $`𝐬=(s_{ij})_{i,j}`$ with entries in $``$ such that $`s_{i+n,j+n}=s_{i,j}`$ and $`_j_{i=1}^ns_{ij}=r.`$ To each such $`𝐬M_{r,n}`$ we associate the $`𝔾`$-orbit $`Y_𝐬`$ whose elements are the pairs $`(L,L^{})`$ for which
$$s_{ij}=\mathrm{dim}_𝔽\left(\frac{L_iL_j^{}}{(L_iL_{j1}^{})+(L_{i1}L_j^{})}\right).$$
For $`𝐢,𝐣𝒜_r^n`$ we denote by $`M_{\mathrm{𝐢𝐣}}`$ the set of all $`𝐬`$ such that $`Y_𝐬Y_𝐢\times Y_𝐣`$. It is easy to see that
$$M_{ij}=\{𝐬M_{r,n}|\underset{j}{}s_{ij}=\mathrm{\#}𝐢^1(i),\underset{i}{}s_{ij}=\mathrm{\#}𝐣^1(j)\}.$$
In particular, $`M_{\mathrm{𝐢𝐣}}`$ is naturally identified with $`𝔖_𝐢\backslash \widehat{𝔖}_r/𝔖_𝐣`$.
###
Let us associate to each $`𝐦=a_i^l[i;l)`$ the matrix $`(m_{i,j})_rM_{r,n}`$ with $`m_{i,j}=a_i^{ji+1}`$. The set
$$M^+=\{(m_{i,j})_{i,j}|m_{i+n,j+n}=m_{i,j},i>jm_{i,j}=0\}$$
is then identified with $``$. If $`𝐢𝒜_r^n`$ and $`𝐦M^+`$ we let $`𝐦^𝐢_𝐣M_{\mathrm{𝐢𝐣}}`$ be the matrix whose $`(i,j)`$th entry is
$$\delta _{ij}(\mathrm{\#}𝐢^1(j+1)\underset{kj}{}m_{kj})+(1\delta _{ij})m_{i+1,j}.$$
###### Proposition (\[VV\]).
The map $`\mathrm{\Theta }:𝐔_n^{}\widehat{𝐒}_{n,r}`$ is an algebra morphism satisfying $`\mathrm{\Theta }(\overline{u})=\tau (\mathrm{\Theta }(u))`$ for every $`u𝐔_n^{}`$. Furthermore,
$$\mathrm{\Theta }(𝐟_𝐦)=\underset{𝐢|𝐦^𝐢M^+}{}\stackrel{~}{T}_{𝐦^𝐢},\mathrm{\Theta }(𝐛_𝐦)=\underset{𝐢|𝐦^𝐢M^+}{}𝐜_{𝐦^𝐢}.$$
###
It follows from the above Proposition that $`𝐓_{n,r}`$ is endowed with a canonical $`𝐔_n^{}`$-module structure.
### 3.5
Let $`e_i^{}`$, $`i/n`$ be the adjoint of the left multiplication by $`𝐟_i`$. Set $`𝐑=_i\mathrm{Ker}e_i^{}𝐔_n^{}`$. Let us identify the ring of symmetric polynomials $`\mathrm{\Gamma }_r`$ with $`Z_r^{}`$ by $`y_iX_i^1`$.
###### Theorem (\[S\]).
The vector space $`𝐑`$ is a graded central subalgebra of $`𝐔_n^{}`$ and the multiplication map induces an isomorphism $`𝐔_v^{}(\widehat{𝔰𝔩}_n)_𝔸𝐑\stackrel{}{}𝐔_n^{}`$. Moreover there exists surjective algebra morphisms $`i_r:𝐑Z_r^{}`$ and an algebra isomorphism $`i:𝐑\mathrm{\Gamma }`$ such that
$$\rho _r\mathrm{\Theta }=\sigma _ri_r,i=\underset{}{\mathrm{lim}}i_r.$$
###
Let $`s_\lambda \mathrm{\Gamma }`$ be the Schur polynomial associated to $`\lambda \mathrm{\Pi }`$, and set $`a_\lambda =i^1(s_\lambda )`$. Then $`i_r(a_\lambda )=s_\lambda (X_1^1,\mathrm{},X_r^1)`$ for any $`rl(\lambda )`$. For $`𝐦`$, define polynomials $`J_𝐦^\lambda [v,v^1]`$ by
$$a_\lambda =\underset{𝐦}{}J_{\lambda ,𝐦}𝐛_𝐦.$$
(3.3)
###### Corollary.
For any $`r`$ and $`𝐢𝒜_r^n`$ we have
$$𝐞_𝐢s_\lambda (X_1^1,\mathrm{}X_r^1)=\underset{𝐦|𝐦^𝐢M^+}{}J_{\lambda ,𝐦}𝐜_{𝐦^𝐢}.$$
(3.4)
Proof. This follows by applying $`a_\lambda `$ to $`𝐞_𝐢1𝐓_{n,r}`$, and using Theorem 3.5 and Proposition 3.4.$`\mathrm{}`$
### Remarks.
i) Let us consider the case $`n=1`$ and $`𝐢=(1^r)`$. Then $`=\mathrm{\Pi }`$ and $`𝐔_1^{}=𝐑\stackrel{i}{}\mathrm{\Gamma }`$, and it is known that $`i`$ identifies the Poincaré-Birkhoff-Witt basis element $`𝐟_\lambda `$ with the Hall-Littlewood polynomial $`P_\lambda `$ (see \[Mac\], Chap. III). In particular, $`K_\mu ^\lambda (v)`$ is the Kostka-Foulkes polynomial and from (3.4) we recover the well-known result of Lusztig (\[L1\]) concerning the Satake isomorphism
$$(\underset{\sigma 𝔖_r}{}T_\sigma )s_\lambda (X_1^1,\mathrm{},X_r^1)=\underset{\mu \mathrm{\Pi }}{}K_\mu ^\lambda (v)\stackrel{~}{T}_{𝔖_r\mu 𝔖_r}.$$
ii)Define the following symmetric bilinear form on $`𝐔_n^{}`$ (the Green’s scalar product) :
$$𝐟_𝐦,𝐟_𝐦^{}=v^{2\mathrm{dim}\mathrm{Aut}(𝐦)}\frac{(1v^2)^{|𝐦|}}{|\mathrm{Aut}(𝐦)|}\delta _{𝐦,𝐦^{}},$$
where $`\mathrm{Aut}(𝐦)`$ stands for the group of automorphism of any representation in the orbit $`O_𝐦`$ and $`|a_i^l[i;l)|=_{i,l}la_i^l`$. It is natural to consider the restriction of this scalar product $`(,)`$ on $`𝐔_n^{}`$ to $`𝐑\stackrel{i}{}\mathrm{\Gamma }`$. Let $`^{\mathrm{per}}`$ denote the set of multisegments of the form $`𝐦=a_i^l[i;l)`$ such that $`a_i^l=a_j^l`$ for all $`i,j`$. By \[S\], Proposition 2.4 we have
$$𝐑=\left(\underset{𝐦^{\mathrm{per}}}{}𝔸𝐛_𝐦\right)^{}.$$
Hence the restriction of $`(,)`$ to $`𝐑`$ is nondegenerate. When $`n=1`$ this restriction coincides, up to a constant, with the Hall-Littlewood scalar product. Let $`(p_\mu )_{\mu \mathrm{\Pi }}`$ be the basis of power-sum symmetric functions and let $`z_{(1^{m_1}2^{m_2}\mathrm{})}=_im_i!i^{m_i}.`$
###### Conjecture.
The restriction of Green’s scalar product on $`𝐑𝐔_n^{}`$ is given by
$$(p_\lambda ,p_\mu )=\delta _{\lambda ,\mu }z_\lambda v^{2(n1)|\lambda |}(1v^2)^{n|\lambda |}\underset{i=1}{\overset{l(\lambda )}{}}\frac{1v^{2n\lambda _i}}{(1v^{2\lambda _i})^2}.$$
This scalar product can be seen as a higher-rank analogue of the Hall-Littlewood scalar product.
## 4 Uglov’s Fock spaces
### 4.1
Let $`n,l`$ be positive integers and let $`𝐬_l^l`$. Following \[JMMO\], Uglov attached to this data an integrable $`𝐔_v(\widehat{𝔰𝔩}_n)`$-module $`\mathrm{\Lambda }_{𝐬_l}^{\mathrm{}}`$ equipped with a distinguished $`𝔸`$-basis $`\{|\lambda _l,𝐬_l\}`$, $`\lambda _l\mathrm{\Pi }^l`$ (the higher-level Fock space, see \[U\], Section 1). The Fock space $`\mathrm{\Lambda }_{𝐬_l}^{\mathrm{}}`$ is also endowed with an action of a Heisenberg algebra $``$ generated by operators $`B_m`$, $`m^{}`$ (see \[U\], Sections 4.2, 4.3). Moreover, the $`𝐔_v(\widehat{𝔰𝔩}_n)`$-action and the $``$-action commute.
### Remark.
When $`l=1`$, Uglov’s Fock space coincides with the Fock space $`\mathrm{\Lambda }^{\mathrm{}}`$ introduced in \[KMS\].
### 4.2
We now extend the action of $`𝐔_v^{}(\widehat{𝔰𝔩}_n)`$ on $`\mathrm{\Lambda }_{𝐬_l}^{\mathrm{}}`$ to an action of $`𝐔_n^{}`$. We follow the method of Varagnolo-Vasserot \[VV\], Section 5. Let $`𝐔_{\mathrm{}}^{}`$ be the Hall algebra of the quiver of type $`A_{\mathrm{}}`$. It is known that $`𝐔_{\mathrm{}}^{}=𝐔_v^{}(𝔰𝔩_{\mathrm{}})`$. Let $`f_i`$, $`i`$ be the standard generator corresponding to the vertex $`i`$.
###
We associate to each $`\lambda _l=(\lambda ^{(1)},\mathrm{}\lambda ^{(l)})\mathrm{\Pi }^l`$ an l-tuple of Young tableaux $`(T_1,\mathrm{},T_l)`$ such that
1. $`T_d`$ is of shape $`\lambda ^{(d)}`$ for $`d=1,\mathrm{},l`$,
2. The $`(i,j)`$-box of $`T_d`$ is filled with content $`s_d+ij`$.
If $`\lambda _l`$ and $`\mu _l`$ are two $`l`$-multipartitions such that $`\gamma =\mu _l\backslash \lambda _l`$ corresponds to a box with content $`k`$, we say that $`\gamma `$ is an addable $`k`$-box of $`\lambda _l`$ and a removable $`k`$-box of $`\mu _l`$. Let $`\gamma ,\gamma ^{}`$ be two addable $`k`$-boxes of $`\lambda _l`$. We say that $`\gamma <\gamma ^{}`$ if $`\gamma `$ and $`\gamma ^{}`$ belong to $`T_d`$ and $`T_d^{}`$ respectively and $`d<d^{}`$.
Let $`\lambda _l,\mu _l\mathrm{\Pi }^l`$ be such that $`\mu _l\backslash \lambda _l`$ is a $`k`$-box. Define
$$\begin{array}{cc}\hfill N^>(\mu _l,\lambda _l)=& \mathrm{\#}\{addablekboxes\gamma ^{}of\lambda _lsuchthat\gamma ^{}>\gamma \}\hfill \\ & \mathrm{\#}\{removablekboxes\gamma ^{}of\lambda _lsuchthat\gamma ^{}>\gamma \}.\hfill \end{array}$$
###### Proposition.
The following endows $`\mathrm{\Lambda }_{𝐬_l}^{\mathrm{}}`$ with a structure of a $`𝐔_{\mathrm{}}^{}`$-module :
$$f_k|\lambda _l,𝐬_l=\underset{\mu _l}{}v^{N^>(\mu _l,\lambda _l)}|\mu _l,𝐬_l$$
where the sum ranges over all $`\mu _l`$ for which $`\mu _l\backslash \lambda _l`$ is a $`k`$-box.
Proof. Straightforward. $`\mathrm{}`$
###
Define operators $`𝐤_k\mathrm{End}(\mathrm{\Lambda }_{𝐬_l}^{\mathrm{}})`$, $`k`$ by $`𝐤_k|\lambda _l,𝐬_l=v^{N_k(\lambda _l)}|\lambda _l,𝐬_l`$ where
$$N_k(\lambda _l)=\mathrm{\#}\{addablekboxesof\lambda _l\}\mathrm{\#}\{removablekboxesof\lambda _l\}.$$
###
Now let $`d^{()}`$ and set $`\overline{d}=(\overline{d}_1,\mathrm{},\overline{d}_n)`$ where $`\overline{d}_i=_{ji(\mathrm{mod}n)}d_j`$. Let $`V`$ be a $``$-graded $`𝔽`$-vector space of dimension $`d`$ and let $`\overline{V}`$ be the $`/n`$-graded $`𝔽`$-vector space with $`\overline{V}_i=_{ji}V_j`$. The collection of subspaces $`\overline{V}_i=_{ji}V_j`$ defines a filtration of $`\overline{V}`$ whose associated graded is $`V`$. Set
$$E_{\overline{V},V}=\{xE_{\overline{V}}|x(\overline{V}_i)\overline{V}_{i+1}\mathrm{for}\mathrm{all}i\}.$$
Let $`p:E_{\overline{V},V}E_V`$ be the projection onto the graded. Let $`j:E_{\overline{V},V}E_V`$ be the closed embedding. Following \[VV\], define a map $`\gamma _d:𝐔_n^{}[\overline{d}]𝐔_{\mathrm{}}^{}[d]`$ by
$`\gamma _{d|v=q^1}:_{G_{\overline{V}}}(E_{\overline{V}})`$ $`_{G_V}(E_V)`$
$`f`$ $`q^{h(d)}p_!j^{}(f)`$
where $`h(d)=_{i<j,ij}d_i(d_{j+1}d_j)`$.
###
For all $`\lambda _l\mathrm{\Pi }^l`$ and $`x𝐔_n^{}`$ we put
$$x|\lambda _l,𝐬_l=\underset{d}{}\left(\gamma _d(x)\underset{j<i,ji}{}𝐤_i^{d_j}\right)|\lambda _l,𝐬_l.$$
(4.1)
Then (see \[VV\] Section 6.2, and \[A\])
###### Proposition.
Formula (4.1) defines a representation $`\mathrm{\Xi }:𝐔_n^{}\mathrm{End}(\mathrm{\Lambda }_{𝐬_l}^{\mathrm{}})`$ which extends Uglov’s action of $`𝐔_v^{}(\widehat{𝔰𝔩}_n)`$.
### Remarks.
i) The number $`h(d)`$ has the following interpretation. Let $`_d`$ be the variety of filtrations of $`\overline{V}`$ whose associated graded is of dimension $`d`$. Then $`\mathrm{dim}T^{}_d=\mathrm{dim}G_{\overline{V}}+h(d)`$.
ii) The map $`\gamma _d`$ is “upper triangular” in the following sense. Let $`xE_V`$ and define $`r(x)E_{\overline{V}}`$ by $`r(x)_i=_{ji}x_j`$. Then $`\gamma _d(𝐟_𝐦)(x)0r(x)\overline{𝒪_𝐦}`$.
### 4.3
Let $`^{}`$ denote the subalgebra generated by $`B_m`$, $`m^{}`$. Define an algebra isomorphism $`j:\mathrm{\Gamma }\stackrel{}{}^{}`$ by setting $`j(p_m)=B_m`$, where $`p_m`$ is the power-sum symmetric function. Recall the canonical map $`i:𝐑\stackrel{}{}\mathrm{\Gamma }`$ from Theorem 3.5.
###### Lemma.
We have $`\mathrm{\Xi }_{|𝐑}=ji`$.
Proof (sketch). This is shown in a way similar to \[VV\]. We first consider the “limit” $`^{\mathrm{}}`$ of $`𝐓_{n,r}`$ when $`r\mathrm{}`$ (see \[VV\], Section 10). Then $`\mathrm{\Lambda }_{𝐬_l}^{\mathrm{}}`$ is naturally embedded in a certain quotient of $`^{\mathrm{}}`$ (see \[U\], Section 3.3). In particular, the $`𝐔_n^{}`$-action on $`𝐓_{n,r}`$ induces an action on $`^{\mathrm{}}`$ and on $`\mathrm{\Lambda }_{𝐬_l}^{\mathrm{}}`$. Let $`\mathrm{\Xi }^{}`$ denote this last action. It follows from Theorem 3.5 and \[U\], Section 4 that $`\mathrm{\Xi }_{|𝐑}^{}=ji`$. Finally, an easy extension to the higher-level Fock space of the computation in \[VV\], Lemma 10.1 shows that $`\mathrm{\Xi }^{}=\mathrm{\Xi }`$.$`\mathrm{}`$
## 5 Canonical bases of Fock spaces
### 5.1
We keep the settings of the previous Section. Uglov has defined a semilinear involution $`a\overline{a}`$ on $`\mathrm{\Lambda }_{𝐬_l}^{\mathrm{}}`$ (\[U\], Section 4.4) and two canonical bases $`\{𝐛_{\lambda _l}^\pm \}_{\lambda _l\mathrm{\Pi }^l}`$ characterized by the following properties :
$$\overline{𝐛^\pm }_{\lambda _l}=𝐛_{\lambda _l}^\pm ,$$
$$𝐛_{\lambda _l}^+|\lambda _l+v\underset{\mu _l}{}𝕊|\mu _l,𝐛_{\lambda _l}^{}|\lambda _l+v^1\underset{\mu _l}{}\overline{𝕊}|\mu _l.$$
He furthermore computed the transition matrices $`[𝐛_{\lambda _l}^\pm :|\mu _l,𝐬_l]`$. In particular we have the following result.
###### Theorem (\[U\], 3.26).
$$𝐛_{\lambda _l}^{}=\underset{\mu _l}{}𝐏_{\mu _l,\lambda _l}^{,𝐬_l}|\mu _l,𝐬_l.$$
### Remark.
When $`l=1`$, Uglov’s canonical bases coincide with the canonical bases considered by Leclerc-Thibon (\[LT\]). In that setting, the transition matrices above were first obtained by Varagnolo and Vasserot \[VV\].
### 5.2
Let us now consider the nondegenerate scalar product $`(,)`$ on $`\mathrm{\Lambda }_{𝐬_l}^{\mathrm{}}`$ for which $`\{|\lambda _l,𝐬_l\}`$ is orthonormal. Let $`\{𝐛_{\lambda _l}^+\}`$ be the dual basis to $`\{𝐛_{\lambda _l}^+\}`$ with respect to the scalar product $`(,)`$.
Define a semilinear isomorphism $`\mathrm{\Lambda }_{𝐬_l}^{\mathrm{}}\mathrm{\Lambda }_{𝐬_l^{}}^{\mathrm{}},uu^{}`$ by $`|\lambda _l,𝐬_l^{}=|\lambda _l^{},𝐬_l^{}`$.
###### Proposition (\[U\], 5.14).
We have $`(𝐛_{\lambda _l}^+)^{}=𝐛_{\lambda _l^{}}^{}`$.
### 5.3
Let $`𝐁_{𝐬_l}=\{𝐛_{\lambda _l}^+\}_{\lambda _l\mathrm{\Pi }^l}`$ be the (positive) canonical basis of $`\mathrm{\Lambda }_{𝐬_l}^{\mathrm{}}`$.
###### Theorem.
Let $`𝐦`$. Then $`𝐛_𝐦|0,𝐬_l𝐁_{𝐬_l}\{0\}`$.
Proof. Lemma 4.3 implies that the $`𝐔_n^{}`$-action on $`\mathrm{\Lambda }_{𝐬_l}^{\mathrm{}}`$ is the same as that considered in \[S\], Section 4. The result follows from \[S\], Theorem 4.2. $`\mathrm{}`$
### 5.4
Define a map $`\tau _{𝐬_l}:\mathrm{\Pi }^l\{0\}`$ by $`\tau _{𝐬_l}(𝐦)=0`$ if $`𝐛_𝐦|0,𝐬_l=0`$ and $`𝐛_𝐦|0,𝐬_l=𝐛_{\tau _{𝐬_l}(𝐦)}^+`$ otherwise. This map is not easy to describe for a general $`𝐬_l`$. Nevertheless we have the following result.
###
Let $`𝐦`$ and let $`\lambda =(\lambda ^{(1)},\mathrm{},\lambda ^{(n)})`$ be the associated $`n`$-multipartition. Let $`r=_il(\lambda ^{(i)})`$. Set $`𝐢=(1^{l(\lambda ^{(1)})},2^{l(\lambda ^{(2)})},\mathrm{})𝒜_r^n`$ and
$$𝐩=(\lambda _1^{(1)},\mathrm{},\lambda _{l(\lambda ^{(1)})}^{(1)},\lambda _1^{(2)},\mathrm{})(^+)^r.$$
Finally, let $`_𝐢(𝐩)=(p^{(1)},\mathrm{},p^{(l)})`$ and let $`r_i/n`$ be the residue of $`p^{(i)}`$.
###### Lemma.
Suppose that $`s_1s_2\mathrm{}s_l`$ and that $`s_ir_i(\mathrm{mod}n)`$ for $`i=1,\mathrm{}l`$. Then $`\tau _{𝐬_l}(𝐦)=_𝐢(𝐩)`$.
Proof. See appendix.$`\mathrm{}`$
### 5.5
Proof of Theorem 1. Let $`\sigma 𝔖_𝐢\backslash \widehat{𝔖}_r/𝔖_𝐢`$. It follows from Corollary 3.5 that $`J_{\lambda ,\sigma }^𝐢=J_{\lambda ,𝐦}`$ if there exists $`𝐦`$ such that $`𝐦^𝐢=\sigma `$ and $`J_{\lambda ,\sigma }^𝐢=0`$ otherwise. From Section 3.4 we see that
$$𝐦^𝐢=\sigma 𝐢\sigma (^+)^r\mathrm{and}𝐦=\underset{j=1}{\overset{r}{}}[𝐢_j;(𝐢\sigma )_j).$$
Now we compute $`J_{\lambda ,𝐦}`$. Let $`l,𝐩,(r_i)_{i=1}^l`$ be associated to $`𝐦`$ as in Section 5.4. Let $`𝐬_l=(s_1,\mathrm{},s_l)`$ be in the asymptotic region $`s_1s_2\mathrm{}s_l`$ and satisfy $`s_ir_i(\mathrm{mod}n)`$ for all $`i`$. We evaluate both sides of (3.3) on $`|0,𝐬_l\mathrm{\Lambda }_{𝐬_l}^{\mathrm{}}`$. On the one hand, it follows from Lemma 4.3 and Uglov’s description of the action of the Heisenberg algebra \[U\], Proposition 5.3 that
$$a_\lambda |0,𝐬_l=\underset{\mu _1,\mathrm{},\mu _l}{}c_{\mu _1,\mathrm{},\mu _l}^\lambda v^{{\scriptscriptstyle (b1)|\mu _b|}}\left(\underset{\nu _1,\mathrm{},\nu _l}{}e_{\nu _1,\mu _1}\mathrm{}e_{\nu _l,\mu _l}|(\nu _1,\mathrm{},\nu _l),𝐬_l\right)$$
where $`e_{\nu _i,\mu _i}[v^1]`$ are defined by the relations $`s_{\mu _i}|0=_{\nu _i}e_{\nu _i,\mu _i}|\nu _i`$ in the level l=1 Fock space representation of $`𝐔_n^{}`$. But by \[LT\], Theorem 6.9 we have $`s_{\mu _i}|0=𝐛_{n\mu _i}^{}`$ and thus $`e_{\nu _i,\mu _i}=𝐏_{\nu _i,n\mu _i}^{}`$. On the other hand, from Theorem 5.3 we have
$$\underset{𝐧}{}J_{\lambda ,𝐧}𝐛_𝐧|0,𝐬_l=\underset{𝐧,\tau _{𝐬_l}(𝐧)0}{}J_{\lambda ,𝐧}𝐛_{\tau _{𝐬_l}(𝐧)}^+.$$
In particular, $`J_{\lambda ,𝐦}=((𝐛_{\tau _{𝐬_l}(𝐦)}^+)^{},a_\lambda |0,𝐬_l)`$. But by Lemma 5.4 and Proposition 5.2,
$$(𝐛_{\tau _{𝐬_l}(𝐦)}^+)^{}=(𝐛_{_𝐢(p)}^+)^{}=(𝐛_{(\sigma )}^+)^{}=(𝐛_{(\sigma )^{}}^{})^{}.$$
Using the relations $`(\overline{u},v)=(u^{},\overline{v^{}})`$ for any $`u,v\mathrm{\Lambda }_{𝐬_l}^{\mathrm{}}`$ (\[U\], Proposition 5.13) and $`\overline{a_\lambda |0,𝐬_l}=a_\lambda |0,𝐬_l`$ (\[U\], Proposition 4.2) we get
$$J_{\lambda ,𝐦}=(𝐛_{(\sigma )^{}}^{},a_\lambda |0,𝐬_l^{}).$$
Now, from \[LT\], Theorem 7.13 i) we have $`(𝐛_{n\mu }^{})^{}=(v)^{(n1)|\mu |}𝐛_{n\mu ^{}}^{}`$ in the level $`l=1`$ Fock space. Thus
$$\begin{array}{cc}\hfill a_\lambda & |0,𝐬_l=\hfill \\ \hfill =& (v)^{(n1)|\lambda |}\underset{\mu _1,\mathrm{},\mu _l}{}c_{\mu _1,\mathrm{},\mu _l}^\lambda v^{_{b=1}^l(b1)|\mu _b|}\left(\underset{\nu _1,\mathrm{},\nu _l}{}𝐏_{\nu _1,n\mu _1^{}}^{}\mathrm{}𝐏_{\nu _l,n\mu _l^{}}^{}|\nu ,𝐬_l\right)\hfill \end{array}$$
where $`\nu =(\nu _l,\mathrm{},\nu _1)`$. The theorem follows. $`\mathrm{}`$
## 6 On the center of $`𝐔_n^{}`$
###
In this section we give a simple geometric characterization of the central subalgebra $`𝐑𝐔_n^{}`$ in terms of the maps $`\gamma _d:𝐔_n^{}𝐔_{\mathrm{}}^{}`$ defined in Section 4.2.
### 6.1
Let $`d^{()}`$ such that $`d_i\{0,1\}`$ for all $`i`$. Then $`d`$ is the dimension of a unique (noncyclic) multisegment $`𝐧_d=_{k=1}^t[i_k;l_k)`$ in $``$ satisfying the following condition :
$$j,k[i_k,l_k)[i_j,l_j)isnotasegment.$$
(6.1)
Let $`V_d`$ be a $``$-graded $`𝔽`$-vector space of dimension $`d`$. Set $`l(d)=_k(l_k1)`$. Note that it follows from (6.1) that $`E_{V_d}`$ has a unique open $`G_{V_d}`$-orbit, say $`𝒪_d`$.
###### Lemma.
Suppose that $`i_1i_2\mathrm{}i_t`$ and set $`𝐢_t=(i_1,\mathrm{},i_t)`$. Then for any $`𝐟𝐔_{\mathrm{}}^{}[d]`$ we have
$$𝐟|0,𝐢_t=v^{l(d)}𝐟_{|𝒪_d}|((l_1),\mathrm{},(l_k)),𝐢_t.$$
Proof. Note that $`E_{V_d}=_{k=1}^tE_{V_d(k)}`$ where $`V_d(k)=_{l=0}^{l_k1}𝔽V_{i_k+l}`$. Let $`f_kE_{V_d(k)}`$ for $`k=1,\mathrm{},t`$. From (6.1) and Section 4.2 we deduce that
$$f_1\mathrm{}f_t|0,𝐢_t=\underset{\nu _1,\mathrm{},\nu _t}{}d_1(\nu _1)\mathrm{}d_t(\nu _t)|(\nu _1,\mathrm{},\nu _t),𝐢_t$$
where $`f_k|0,i_k=_\nu d_k(\nu )|\nu ,𝐢_k`$ in the level $`l=1`$ Fock space. But from \[VV\], Proposition 5., it is easy to see that $`f_k|0,𝐢_k=v^{(l_k1)}(f_k)_{|𝒪_d(k)}|(l_k),i_k`$ where $`𝒪_d(k)E_{V_d(k)}`$ is the open orbit.$`\mathrm{}`$
### 6.2
Recall the element $`a_\lambda =i^1(s_\lambda )𝐑`$. For any $`\lambda ,\mu \mathrm{\Pi }`$ let $`K_\mu ^\lambda `$ be the Kostka number.
###### Theorem.
Let $`d^{()}`$ such that $`d_i\{0,1\}`$. Then
$$\gamma _d(a_\lambda )_{|𝒪_d}=v^{l(d)+h(d)}K_{(u_1,\mathrm{},u_t)}^\lambda $$
if there exists $`i_k,u_k`$, $`k=1,\mathrm{},t`$ such that $`𝐧_d=_{k=1}^t[i_k;nu_k)`$, and $`\gamma _d(a_\lambda )_{|𝒪_d}=0`$ otherwise.
Proof. Without loss of generality we may assume that $`𝐧=_{k=1}^t[i_k;l_k)`$ where $`i_1>i_2>\mathrm{}>i_t`$. Choose $`d^{}=_{k=1}^t[i_k^{};l_k)`$ where $`i_k^{}i_k(\mathrm{mod}n)`$ and $`i_1^{}i_2^{}\mathrm{}i_t^{}`$. Let $`\xi :E_{V_d^{}}\stackrel{}{}E_{V_d}`$ be the obvious isomorphism. Then $`\xi \gamma _d^{}=\gamma _d`$. Now let us consider the Fock space $`\mathrm{\Lambda }_{𝐢_t^{}}^{\mathrm{}}`$ where $`𝐢_t^{}=(i_1^{},\mathrm{},i_t^{})`$. Using \[U\], Proposition 5.3 we have
$$\begin{array}{cc}\hfill (a_\lambda |0,𝐢_t^{},|((l_1),\mathrm{},(l_t)),& 𝐢_t^{})\hfill \\ & =\underset{\mu _1,\mathrm{},\mu _t}{}c_{\mu _1,\mathrm{},\mu _t}^\lambda v^{_b(b1)|\mu _b|}𝐏_{(l_1),n\mu _1}^{}\mathrm{}𝐏_{(l_t),n\mu _t}^{}\hfill \\ & =\underset{\mu _1,\mathrm{},\mu _t}{}\delta _{(l_1)=n\mu _1}\mathrm{}\delta _{(l_t)=n\mu _t}c_{\mu _1,\mathrm{},\mu _t}^\lambda v^{_b(b1)|\mu _b|}\hfill \end{array}$$
Note that for any $`u_1,\mathrm{},u_t`$ we have $`c_{(u_1),\mathrm{},(u_t)}^\lambda =K_\mu ^\lambda `$ where $`\mu \mathrm{\Pi }`$ is the partition with parts $`\{u_1,\mathrm{},u_t\}`$.
On the other hand, by Lemma 6.1
$$(a_\lambda |0,𝐢_t^{},|((l_1),\mathrm{},(l_t)),𝐢_t^{})=v^{ϵ(d^{},𝐢_t^{})l(d)}\gamma _d^{}(a_\lambda )_{|𝒪_d^{}}$$
where
$$ϵ(d^{},𝐢_t^{})=\underset{l=1}{\overset{t}{}}\underset{ji_l^{};j<i_l^{}}{}d_j^{}.$$
The result now follows from the easily checked identity
$$ϵ(d^{},𝐢_t^{})=\underset{b}{}(b1)|\mu _b|+h(d^{})$$
when there exists $`u_k`$, $`k=1,\mathrm{},t`$ such that $`d^{}=_k[i_k^{},nu_k)`$ and $`\mu _k=(u_k)`$.$`\mathrm{}`$
### Remark.
It follows from Remark 4.2 ii) that the previous theorem gives a characterization of the central element $`a_\lambda `$.
## 7 Appendix
###
In this appendix we prove Lemma 5.4.
### A.1
As in \[U\], Section 4, define a partial order on $`\mathrm{\Pi }^l`$ (depending on $`𝐬_l`$) as follows. Let $`\mu =(\mu ^{(1)},\mathrm{},\mu ^{(l)})\mathrm{\Pi }^l`$. Set $`k_i^{(d)}=\mu _i^{(d)}+s_d+1i`$ for $`d=1,\mathrm{},l`$ and $`i`$. Let us write $`k_i^{(d)}=c_i^{(d)}nm_i^{(d)}`$ where $`c_i\{1,\mathrm{},n\}`$, and let $`𝐤=(k_1>k_2>\mathrm{})`$ be the ordered sequence whose underlying set is $`\{c_i^{(d)}+n(d1)nlm_i^{(d)}|i,d=1,\mathrm{},l\}`$. Let $`s=s_1+\mathrm{}+s_l`$. It is easy to see that $`k_i=s+1i`$ for $`i0`$ and we denote by $`\zeta (\mu )`$ the partition such that $`\zeta (\mu )_i=k_is+i1`$. Now let $`\mu ,\nu \mathrm{\Pi }^l`$. By definition, we set $`\mu \nu `$ if $`\zeta (\mu )\zeta (\nu )`$.
### A.2
From now on we assume that $`v=1`$.
###
It is more convenient to work with a different basis than $`\{𝐟_𝐧\}`$. Let $`𝐧`$ and let $`x𝒪_𝐧`$. Set $`V_k=\mathrm{Ker}x^k`$ and let $`\alpha ^1,\mathrm{},\alpha ^r^{/n}`$ be such that
$$\mathrm{dim}V_k=\alpha ^1+\mathrm{}+\alpha ^k,k=1,\mathrm{},r$$
and $`\mathrm{dim}V_r=\mathrm{dim}𝐧`$. Let $`𝐟_{\alpha ^i}𝐔_n^{}`$ be the characteristic function of the trivial representation of the quiver $`\stackrel{~}{A}_{n1}`$ on $`V_{\alpha ^i}V_i/V_{i1}`$.
###### Lemma 1 (\[VV\], Section 13).
We have $`𝐟_{\alpha ^1}\mathrm{}𝐟_{\alpha ^r}𝐟_𝐧+_{𝐥<𝐧}𝐟_𝐥`$.
Now let $`𝐧,𝐥`$ such that $`\mathrm{dim}𝐧=\mathrm{dim}𝐥`$. Let $`(\beta ^k)`$ and $`(\gamma ^k)`$ be the sequences of dimensions attached as above to $`𝐧`$ and $`𝐥`$ respectively. If $`𝐮,𝐯/n`$ we write $`𝐮𝐯`$ if $`𝐮_i𝐯_i`$ for all $`i/n`$.
###### Lemma 2.
We have $`𝐧𝐥`$ if and only if
$$\beta ^1+\mathrm{}+\beta ^k\gamma ^1+\mathrm{}+\gamma ^kforallk.$$
(a)
Proof. Straightforward.$`\mathrm{}`$
###
We will write $`(\beta ^k)(\gamma ^k)`$ if (a) holds and if $`_k\beta ^k=_k\gamma ^k`$. Let $`(\alpha ^1,\mathrm{},\alpha ^r)`$ be the sequence attached to $`𝐦`$. We will first prove
$`𝐟_{\alpha ^1}\mathrm{}𝐟_{\alpha ^r}|0,𝐬_l`$ $`^{}|_𝐢(𝐩),𝐬_l+{\displaystyle \underset{\mu _𝐢(𝐩)}{}}|\mu ,𝐬_l`$ (b)
$`𝐟_{\beta ^1}\mathrm{}𝐟_{\beta ^r}|0,𝐬_l`$ $`{\displaystyle \underset{\mu _𝐢(𝐩)}{}}|\mu ,𝐬_lforall(\beta ^k)>(\alpha ^k).`$ (c)
###### Lemma 3.
Let $`\mu =(\mu ^{(1)},\mathrm{},\mu ^{(l)})\mathrm{\Pi }^l`$ and let $`\beta N^{/n}`$. We have
$$𝐟_\beta |\mu ,𝐬_l=\underset{\nu }{}|\nu ,𝐬_l$$
where the sum ranges over all multipartitions $`\nu =(\nu ^{(1)},\mathrm{},\nu ^{(l)})`$ such that
1. $`\nu ^{(i)}\backslash \mu ^{(i)}`$ is a skew diagram with at most one box in each row,
2. The number of boxes in $`_i\nu ^{(i)}\backslash \mu ^{(i)}`$ with content $`j\mathrm{mod}n`$ is $`\beta _j`$.
Proof. Let $`d^{()}`$ such that $`d\beta (\mathrm{mod}n)`$. Then $`\gamma _{d|v=1}(𝐟_\beta )=\stackrel{}{}_if_i^{(d_i)}`$, where $`\stackrel{}{}`$ denotes the ordered product from $`\mathrm{}`$ to $`\mathrm{}`$ (see \[VV\], Remark 6.1) and where $`f_i^{(d_i)}`$ is the divided power. Moreover, for any $`\sigma \mathrm{\Pi }^l`$,
$$f_i|\sigma ,𝐬_l=\underset{\gamma }{}|\gamma ,𝐬_l$$
where the sum ranges over all $`\gamma \mathrm{\Pi }^l`$ such that $`\gamma \backslash \sigma `$ is an $`i`$-box. The Lemma now follows from Section 4.2. $`\mathrm{}`$
###
Finally, recall that $`s_1s_2\mathrm{}s_l`$. It is clear from the definition that for $`\mu ,\lambda \mathrm{\Pi }^l`$,
$$\mu \lambda ksuchthat\mu ^{(i)}=\lambda ^{(i)}fori=1,\mathrm{},k1and\mu ^{(k)}\lambda ^{(k)}.$$
(d)
Note that $`\alpha _i^k`$ is equal to the number of boxes with content $`i`$ in the slice $`s_k`$ of the diagram $`D_𝐩`$ associated to $`𝐩`$. Statements (b) and (c) now easily follow by Lemma 3 and by construction of $`_𝐢(𝐩)`$.
### A.3
By \[U\], Theorem 2.4 it possible to choose $`s_ls_{l+1}\mathrm{}s_t`$ for some $`t0`$ in such a way that $`𝐛_𝐦|0,𝐬_t0`$, where $`𝐬_t=(s_1,\mathrm{},s_t)`$.
###### Lemma 4.
We have $`𝐛_𝐦|0,𝐬_t=𝐛_{\stackrel{~}{}_𝐢(𝐩)}^+`$, where $`\stackrel{~}{}_𝐢(𝐩)=(_𝐢(𝐩),0^{tl})`$.
Proof. By Lemma 1, we have
$$𝐟_{\alpha ^1}\mathrm{}𝐟_{\alpha ^r}|0,𝐬_t|\tau _{𝐬_t}(𝐦),𝐬_t+\underset{\mu <\tau _{𝐬_t}(𝐦)}{}|\mu ,𝐬_t.$$
But from (b) and (d) it is clear that
$$𝐟_{\alpha ^1}\mathrm{}𝐟_{\alpha ^r}|0,𝐬_t|\stackrel{~}{}_𝐢(𝐩),𝐬_t+\underset{\mu \stackrel{~}{}_𝐢(𝐩)}{}|\mu ,𝐬_t.$$
Hence $`\tau _{𝐬_t}(𝐦)=\stackrel{~}{}_𝐢(𝐩)`$. $`\mathrm{}`$
###
In particular,
$$𝐛_𝐦|0,𝐬_t|\stackrel{~}{}_𝐢(𝐩),𝐬_t+\underset{\mu <\stackrel{~}{}_𝐢(𝐩)}{}|\mu ,𝐬_l.$$
Consider the projection $`\pi :\mathrm{\Lambda }_{𝐬_t}^{\mathrm{}}\mathrm{\Lambda }_{𝐬_l}^{\mathrm{}}`$ given by
$$|(\mu ^{(1)},\mathrm{},\mu ^{(t)}),𝐬_t\{\begin{array}{cc}|(\mu ^{(1)},\mathrm{},\mu ^{(l)}),𝐬_l\hfill & if\mu ^{(j)}=0forj>l\hfill \\ 0\hfill & otherwise\hfill \end{array}$$
It is clear from (4.1) that $`\pi (𝐛_𝐦|0,𝐬_t)=𝐛_𝐦|0,𝐬_l`$. Hence
$$𝐛_𝐦|0,𝐬_l|_𝐢(𝐩),𝐬_l+\underset{\mu <_𝐢(𝐩)}{}|\mu ,𝐬_l.$$
This proves Lemma 5.4 $`\mathrm{}`$
Acknowledgements
I am grateful to J. Dat for interesting discussions and for pointing out to me the work \[H1\], \[H2\], to B. Leclerc, A. Schilling and E. Vasserot for valuable advice at various stages of this work. I would like to thank the MIT mathematics department for its hospitality. This research was partially conducted for the Clay Mathematics Institute.
Olivier Schiffmann, MIT, 77 Massachusetts Avenue, CAMBRIDGE 02139, USA;
email: schiffma@math.mit.edu
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# 1 Introduction
## 1 Introduction
Recently there have been tremendous progresses in understanding non-BPS configrations of D-branes and tachyon condensations in them, pioneered by Sen (for a review see ). In Superstring theory most of these systems are realized as either brane-antibrane systems or non-BPS D-branes . The open strings between a Dp-brane and an anti-Dp-brane are projected by the GSO projection opposite to the usual cases and the tachyon survives. A non-BPS D-brane is defined as a D-brane without any GSO projections and the tachyonic instability also occurs. Sen argued that if the condensation of the constant tachyon field stabilizes the system, then the system will finally go down to the vacuum and if the condensation has nontrivial configurations such as kinks or vortexes, then the final object will be D-branes of corresponding codimension . For example, the kink configuration in $`Dp\overline{Dp}`$ brane system is identified as a non-BPS D$`(p1)`$-brane.
Three different approaches have been considered <sup>4</sup><sup>4</sup>4Quite recently a new approach which utilizes the noncommutative field theory description of the world volume theory has been considered in the presence of a large B-field . to analyze those systems. The first one is to use conformal field theory descriptions of string world sheet with a boundary . In this method, the tachyon condensation can be regarded as a marginal deformation of the boundary conformal field theory (BCFT) at a special radius. The second is the K-theory approach by Witten. He argued that the topological charges of lower D-branes in a non-BPS configuration of D-branes are classified by the corresponding K-group . In other words we can say that the topological configurations of tachyon fields one-to-one correspond to the element of the K-group. The third one is the string field theory description . The tachyon potential has been calculated in the case of a D-brane in bosonic string and a non-BPS D-brane in Superstring . The numerical results are in good agreement with the Sen’s conjecture that the system at the minimum of the tachyon potential can be identified as the vacuum. The generations of lower dimensional D-brane charges have been also discussed in this formalism .
In this paper we are interested in the first approach. From the viewpoint of the open strings the BCFT descriptions of tachyon kink condensations have been given in for a brane-antibrane system or a non-BPS D-brane in the presence of the orbifold and orientation projection. In order to discuss the generation of codimension two D-branes, the vortex line configuration of the tachyon field is needed and is realized in as a pair of the vortex and anti-vortex.
On the other hand we can use the boundary state formalism (for example see ), which can give more systematic CFT description of D-branes. In this formalism D-branes are constructed in the closed string Hilbert space. Therefore the couplings of D-branes to NSNS, RR-fields can be written down explicitly. The equivalence between the open string description and the closed string description should be required as in the usual BPS D-branes and this crucial constraint is called Cardy’s condition . The boundary state description of the tachyon kink pair in flat space was first constructed in . Also in the case of the bosonic string the boundary state description was discussed in . In the first half of this paper we extend this construction to the case of tachyon vortex pairs in flat space. We construct explicitly the boundary state which describes the tachyon vortex pair condensation in $`D2\overline{D2}`$ system at the critical radius. The Cardy’s condition is established and the emergence of lower D-brane RR-charges is shown explicitly. At the point where the tachyon condensation is maximum the boundary state of the system is identified as that of a $`D0\overline{D0}`$ system. Many other points correspond to the bound state of $`D2\overline{D2}`$ and $`D0\overline{D0}`$. Also the requirement<sup>5</sup><sup>5</sup>5In the case of a non-BPS D-brane the similar requirement was mentioned in . of the nontrivial “Chan-Paton factor of closed string” is verified in this formalism. Further we generalize these results into the higher codimension cases.
In the latter half we treat the case of $`T^4/𝐙_2`$ orbifold. Remarkably it is shown that the twisted sector boundary state can be written as the untwisted sector boundary state of another fields on the world sheet at the critical radius by performing bosonizations and fermionizations. Using this fact we can describe the tachyon kink condense explicitly starting from a $`D0\overline{D0}`$ system and show that the untwisted RR charge vanishes and the twisted RR charge remains after the condensation, verifying the known identification of the final object as a non-BPS D1-brane between the fixed points. This boundary state approach enables us to generalize the tachyon kink in the orbifold theory into the higher codimention cases such as the decay mode from $`D4\overline{D4}`$ to $`D0\overline{D0}`$, which has not been discussed before. We also discuss the relation between the bose-fermi degeneracy and the “bosonization procedure” in the boundary state description.
The paper is organized as follows. In section 2 we review the known results about the BCFT description of a tachyon vortex pair in a $`D2\overline{D2}`$ system and a tachyon kink pair in a orbifolded $`D0\overline{D0}`$ system. Further we investigate some details. In section 3 we construct the explicit boundary state which describes the condensation of the tachyon vortex pair in a $`D2\overline{D2}`$ system. We show the final object can be identified with $`D0\overline{D0}`$. Next we generalize the results and see that the tachyon condensation generates D-brane charges of higher codimension. In section 4 we study the tachyon kink in $`T^4/𝐙_2`$ orbifold. We give the corresponding boundary state and identify the final object. We also generalize the result into higher codimension cases. In section 5 we summarize the results and draw conclusions. In appendix A we give a breif review of boundary state and show our conventions. In appendix B we prove that the “bosonized” boundary state indeed satisfies the original boundary condition including detailed cocycle factors.
## 2 CFT description of tachyon condensation
In this section we review the descriptions of tachyon condensation from the viewpoint of open strings. These results are needed when we compare the results with those gained in the boundary state formalism. First we see the vortex-type tachyon condensation in the $`D2\overline{D2}`$ system following and investigate some details in a slightly generalized situation. Next we review several known facts about the brane-antibrane system in $`T^4/𝐙_2`$ orbifold and the tachyon condensation in that system. Such a system was first discussed in and also considered in using the T-dualized picture. In this paper we consider type II string theory only in the weak coupling region.
### 2.1 Tachyon condensation in a $`D2\overline{D2}`$ system
We take a parallel $`D2`$-brane and an anti $`D2`$-brane in type IIA string theory along $`x^1,x^2`$ and compactify these directions on a torus of radii <sup>6</sup><sup>6</sup>6In this paper we use $`\alpha ^{}=1`$ unit. $`R_1=1,R_2=1`$. Then we set a $`𝐙_2`$ Wilson line along each circle. There are four types of Chan-Paton factors for the open strings in $`D2\overline{D2}`$ system and are denoted by $`1,\sigma _1,\sigma _2,\sigma _3`$ using Pauli matrices. We use $`1,\sigma _3`$ in order to represent the open strings with both ends on the same brane and the spectrum is determined by the conventional GSO projection. On the other hand $`\sigma _1,\sigma _2`$ correspond to the open strings with two ends on two different branes and follow the opposite GSO projection allowing the tachyon in the spectrum.
We consider the condensation of the following two types<sup>7</sup><sup>7</sup>7There are also other two marginal deformations which represent other tachyon condensations. But these correspond to the shift of the vortex line center and the physical phenomena which occurs by such tachyon condensations do not change if we ignore these. Thus we only consider the tachyon fields (2.1),(2.2) below. of the tachyon field
$`T_{(1)}(x^1,x^2)`$ $`=`$ $`e^{i\frac{1}{2}(x^1+x^2)}e^{i\frac{1}{2}(x^1+x^2)},`$ (2.1)
$`T_{(2)}(x^1,x^2)`$ $`=`$ $`ie^{i\frac{1}{2}(x^1x^2)}+ie^{i\frac{1}{2}(x^1x^2)}.`$ (2.2)
If we switch on only one of these, we get the tachyon kink configuration and a codimension one D-brane or a non-BPS D1-brane will be generated. On the other hand if we condense both at the same time, the tachyon vortex line pair configuration will lead to a pair of codimension two D-branes or a $`D0\overline{D0}`$ system.
The corresponding open string vertex operators in (0)-picture are written as
$`V_{T1}=(\chi ^1+\chi ^2)(e^{i\frac{1}{2}(X^1+X^2)}+e^{i\frac{1}{2}(X^1+X^2)})\sigma _1,`$
$`V_{T2}=(\chi ^2\chi ^1)(e^{i\frac{1}{2}(X^1X^2)}+e^{i\frac{1}{2}(X^1X^2)})\sigma _2,`$ (2.3)
where $`X^i=X_R^i+X_L^i,\chi ^i=\chi _R^i+\chi _L^i(i=1,2)`$ denote the bosonic fields on the string world sheet in NS-R formalism and their superpartners.
Notice that at the radii $`R_1=1,R_2=1`$ (critical radii) the lightest tachyon vertex operators become marginal owing to the Wilson lines and the tachyon condensation corresponding to such operators can be treated as the marginal deformation of CFT.
Now let us rotate the coordinates by $`\frac{\pi }{4}`$
$`Y^1`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(X^1+X^2),Y^2={\displaystyle \frac{1}{\sqrt{2}}}(X^1X^2),`$
$`\psi ^1`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\chi ^1+\chi ^2),\psi ^2={\displaystyle \frac{1}{\sqrt{2}}}(\chi ^1\chi ^2).`$ (2.4)
This procedure enables us to use the method of bosonization and fermionization as follows
$`e^{i\sqrt{2}Y_R^i}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\xi _R^i+i\eta _R^i)\tau _i,e^{i\sqrt{2}Y_L^i}={\displaystyle \frac{1}{\sqrt{2}}}(\xi _L^i+i\eta _L^i)\tau _i,`$
$`e^{i\sqrt{2}\varphi _R^i}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\xi _R^i+i\psi _R^i)\stackrel{~}{\tau }_i,e^{i\sqrt{2}\varphi _L^i}={\displaystyle \frac{1}{\sqrt{2}}}(\xi _L^i+i\psi _L^i)\stackrel{~}{\tau }_i,`$
$`e^{i\sqrt{2}\varphi _R^i}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\eta _R^i+i\psi _R^i)\widehat{\tau }_i,e^{i\sqrt{2}\varphi _L^i}={\displaystyle \frac{1}{\sqrt{2}}}(\eta _L^i+i\psi _L^i)\widehat{\tau }_i,`$
where $`\tau _i,\stackrel{~}{\tau }_i\widehat{\tau }_i(i=1,2)`$ are cocycle factors and we also assume $`\psi _{L,R}^i,\eta _{L,R}^i`$ have the cocycle factor $`\tau _3,\stackrel{~}{\tau }_3`$. To be exact, other kinds of cocycle factors are needed in front of the exponential fields. The latter type of cocycle factors, which we will call second-type cocycle factors below, can not be ignored when we later discuss the bosonizations and fermionizations of boundary states. We leave the details in the appendix B.
The operator product expansions (OPE) among these fields are<sup>8</sup><sup>8</sup>8Note that the factors $`i`$ in the bosonic field OPE’s are due to second-type cocycle factors.
$`Y_R^i(z)Y_R^j(0){\displaystyle \frac{1}{2}}\delta _{ij}\text{ln}z`$ , $`Y_L^i(\overline{z})Y_L^j(0){\displaystyle \frac{1}{2}}\delta _{ij}\text{ln}\overline{z},`$
$`\xi _R^i(z)\xi _R^j(0)\delta _{ij}{\displaystyle \frac{i}{z}}`$ , $`\xi _L^i(\overline{z})\xi _L^j(0)\delta _{ij}{\displaystyle \frac{i}{\overline{z}}},`$
$`\eta _R^i(z)\eta _R^j(0)\delta _{ij}{\displaystyle \frac{i}{z}}`$ , $`\eta _L^i(\overline{z})\eta _L^j(0)\delta _{ij}{\displaystyle \frac{i}{\overline{z}}}.`$ (2.6)
Also the following identities are useful:
$`\eta _R^i\xi _R^i=i\sqrt{2}Y_R^i`$ , $`\eta _L^i\xi _L^i=i\sqrt{2}\overline{}Y_L^i,`$
$`\psi _R^i\xi _R^i=i\sqrt{2}\varphi _R^i`$ , $`\psi _L^i\xi _L^i=i\sqrt{2}\overline{}\varphi _L^i.`$ (2.7)
Now we can express the tachyon vertex operators (2.1) in the following convenient way
$`V_{T1}`$ $`=`$ $`2i\psi ^1\xi ^1\tau _2\sigma _1=2\sqrt{2}\varphi ^1\tau _2\sigma _1,`$
$`V_{T2}`$ $`=`$ $`2i\psi ^2\xi ^2\tau _1\sigma _2=+2\sqrt{2}\varphi ^2\tau _1\sigma _2,`$ (2.8)
where $``$ denotes tangential derivative along the boundary. Then the tachyon condensation is represented as the insertion of the following Wilson lines in terms of the field $`\varphi `$
$`\mathrm{exp}\left({\displaystyle \frac{i\alpha }{2\sqrt{2}}}{\displaystyle \varphi ^1\tau _2\sigma _1}+{\displaystyle \frac{i\beta }{2\sqrt{2}}}{\displaystyle \varphi ^2\tau _1\sigma _2}\right),`$ (2.9)
where $``$ denotes integration along the boundary and $`\alpha ,\beta `$ mean parameters of tachyon condensations. Notice that $`\tau _2\sigma _1`$ commutes with $`\tau _1\sigma _2`$ and the above Wilson line is well defined without path ordering. The open string spectrum in the R sector does not change in the presence of the Wilson line because for the R sector $`\varphi `$ satisfies Neumann boundary condition at one end and Dirichlet boundary condition at the other end and there is no zero mode for $`\varphi `$ . Therefore we will investigate only the NS sector.
Now let us define several projection operators in the following way
$`(1)^F`$ $`:`$ $`|0|0,\psi ^i\psi ^i,(\xi ^i,\eta ^i)(\xi ^i,\eta ^i),`$
$`h_1`$ $`:`$ $`(\xi ^1,\eta ^1)(\xi ^1,\eta ^1),(\xi ^2,\eta ^2)(\xi ^2,\eta ^2),\psi ^i\psi ^i,`$
$`Y_{L,R}^1Y_{L,R}^1+{\displaystyle \frac{\pi }{\sqrt{2}}},Y_{L,R}^2Y_{L,R}^2,`$
$`h_2`$ $`:`$ $`(\xi ^1,\eta ^1)(\xi ^1,\eta ^1),(\xi ^2,\eta ^2)(\xi ^2,\eta ^2),\psi ^i\psi ^i,`$ (2.10)
$`Y_{L,R}^1Y_{L,R}^1,Y_{L,R}^2Y_{L,R}^2+{\displaystyle \frac{\pi }{\sqrt{2}}}.`$
As is clear from the above definition, $`(1)^F`$ is the fermion number on the world sheet and $`h_1,h_2`$ are the translation operators in the direction of $`Y^1,Y^2`$. We also define $`(1)^{F^\varphi },h_1^\varphi ,h_2^\varphi `$ similarly for $`\varphi `$.
Since so far we have implicitly assumed the radius of circle in the direction of $`Y^1,Y^2`$ is $`\sqrt{2}`$, we should have a certain constraint in order to realize the physical periodicity $`X^1X^1+2\pi ,X^2X^2+2\pi `$ taking the effect of the Wilson line into consideration. Such a constraint is given as
$$(1)^Fh_1h_2=(1)^{F^\varphi }h_1^\varphi h_2^\varphi =1,$$
(2.11)
where we used eq.(2.10). There are eight sectors in NS sector which survive this projection as follows
$`11,1\tau _3,\sigma _31,\sigma _3\tau _3,`$
$`\sigma _1\tau _1,\sigma _1\tau _2,\sigma _2\tau _1,\sigma _2\tau _2.`$ (2.12)
Four of these are insensitive to the tachyon condense or equally the insertion of the Wilson lines. But the momenta of $`\varphi `$ in the other four sectors are shifted in proportion to the deformation parameters $`\alpha ,\beta `$. The details are shown in Table 1. Note that $`\alpha ,\beta `$ have periodicity $`\alpha \alpha +2,\beta \beta +2`$ by applying the same argument discussed in .
The main claim in is that if the tachyon condensation develops into the point $`\alpha =1,\beta =1`$, then the system is identified as the $`D0\overline{D0}`$ system where $`D0`$-brane and $`\overline{D0}`$-brane sit at $`(x^1,x^2)=(0,0)`$ and $`(x^1,x^2)=(\pi ,\pi )`$. For example the open string spectrum at $`\alpha =0,\beta =0`$ is shown to be the same as the spectrum at $`\alpha =1,\beta =1`$ because the momentum shift is $`\mathrm{\Delta }P_{\varphi ^1}=\pm \frac{1}{\sqrt{2}},\mathrm{\Delta }P_{\varphi ^2}=\pm \frac{1}{\sqrt{2}}`$ and for each state the value of $`(1)^Fh_1h_2`$ does not change. But in order to prove the claim it is necessary to distinguish $`D0\overline{D0}`$ from its T-dual equivalent $`D2\overline{D2}`$ at the self-dual radii $`R_1=R_2=1`$ and explain<sup>9</sup><sup>9</sup>9For the NSNS-charge the explanation is given in perturbatively by considering a certain disk amplitude where a NSNS vertex operator was inserted. the emergence of NSNS and RR-charge corresponding to $`D0\overline{D0}`$. For these purposes the boundary state description which will be discussed in the following sections is very useful and systematic since the D-branes are represented in the closed string Hilbert space in this description.
### 2.2 Tachyon condensation in $`T^4/𝐙_2`$ orbifold
Let us denote<sup>10</sup><sup>10</sup>10Here we have used not $`x`$ but $`y`$ because later we will identify these coordinates as $`\frac{\pi }{4}`$ rotated ones. $`y^6,y^7,y^8,y^9`$ as the coordinates of $`T^4/𝐙_2`$ with the involution $`I_4:(y^6,y^7,y^8,y^9)(y^6,y^7,y^8,y^9)`$ and assume the radii of the torus are given as $`R^6=\sqrt{2},R^7=R^8=R^9=R`$.
First we consider a fractional $`D0\overline{D0}`$ system where a $`D0`$-brane is sitting on a fixed point $`(y^6,y^7,y^8,y^9)=(0,0,0,0)`$ and a $`\overline{D0}`$-brane is sitting on another fixed point $`(y^6,y^7,y^8,y^9)=(\pi \sqrt{2},0,0,0)`$. Each of them has $`\frac{1}{2}`$ tension and $`\frac{1}{2}`$ RR charge of a bulk D0-brane and can be interpreted as a D2-brane wrapped on the vanishing 2-cycle which corresponds to the fixed point .
Such a system has no tachyonic modes which survive the $`I_4`$ projection and therefore is stable at the critical radius. One of the marginal “tachyon” vertex in (0)-picture which represents a tachyon kink in the $`y^6`$ direction is given as
$$V_T=\psi ^6(e^{i\frac{1}{\sqrt{2}}(Y_R^6Y_L^6)}\pm e^{i\frac{1}{\sqrt{2}}(Y_R^6Y_L^6)})\sigma _1,$$
(2.13)
where $`\pm `$ depend on the relative twisted sector charge of the system or equally the projection $`I_4=\pm 1`$ and below we only consider the case of $`+`$ sign. Using the bosonizations and fermionizations (LABEL:eqn:bos), it is easy to see $`V_T\varphi ^6\sigma _1`$ and the tachyon condensation<sup>11</sup><sup>11</sup>11Note that this marginal deformation is just the opposite to that considered in , where the deformation from non-BPS D1-brane to $`D0\overline{D0}`$ is considered. is described by the following Wilson line<sup>12</sup><sup>12</sup>12If we consider the case of the opposite twisted charge, then we get $`V_T\varphi ^6\sigma _1`$ and $`W=\mathrm{exp}\left(i\frac{\alpha }{2\sqrt{2}}_n\varphi ^6\sigma _1\right)`$, where $`_n`$ denotes derivative in the normal direction.
$$W=\mathrm{exp}\left(i\frac{\alpha }{2\sqrt{2}}\varphi ^6\sigma _1\right).$$
(2.14)
If we condense the tachyon into $`\alpha =1`$, then the system is identified with the non-BPS D1-brane stretching between the fixed points. The justification of this statement will be given later by constructing the boundary state. In this non-BPS D1-brane is identified with a D2-brane wrapped on a non-supersymmetric cycle. Later we will also construct the marginal deformation from $`D4\overline{D4}`$ to $`D0\overline{D0}`$ in this orbifold theory.
There is also a known interesting fact. If we consider a non-BPS D1-brane at the special radius $`R=\frac{1}{\sqrt{2}}`$, then the vacuum amplitude of the system vanishes and the system develops the bose-fermi degeneracy . Later this phenomenon will be discussed in terms of the boundary state description.
## 3 Boundary state description of tachyon condensation
In this section, we construct the boundary state for a $`D2\overline{D2}`$ system and condensate a tachyon vortex pair. Mainly we follow the line of , where a tachyon kink was considered. The crucial difference from that case is the emergence of nontrivial “Chan-Paton factors” in closed string sectors. After the condensation the final object is identified with a $`D0\overline{D0}`$ system as expected. Next we also calculate the vacuum amplitude and investigate the consistency with open string picture. Finally we generalize these results into the higher codimension cases. The definition and brief review of boundary states are given in appendix A.
### 3.1 Bosonization of the boundary state and tachyon condensation
First the boundary state for a D2-brane at $`x^i(i=39)`$ which is extended to $`x^1,x^2`$ without any Wilson lines is given as follows :
$`|D2,x^i`$ $`=`$ $`{\displaystyle \frac{T_{p=2}}{2}}(|D2,x^i_{NSNS}+|D2,x^i_{RR}),`$
$`|D2,x^i_{NSNS}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }\left({\displaystyle \frac{dk}{2\pi }}\right)^7e^{ikx}[|D2,+,k_X^i_{NSNS}|D2,,k_X^i_{NSNS}],`$
$`|D2,x^i_{RR}`$ $`=`$ $`2{\displaystyle }\left({\displaystyle \frac{dk}{2\pi }}\right)^7e^{ikx}[|D2,+,k_X^i_{RR}+|D2,,k_X^i_{RR}],`$ (3.1)
where $`T_p=2^{3p}\pi ^{\frac{7}{2}p}`$ is the normalization<sup>13</sup><sup>13</sup>13This can be determined by computing the cylinder amplitude. of the Dp-brane boundary state and $`k^i(i=39)`$ are the momenta in the direction of $`x^i`$. The explicit forms of $`|D2,\gamma ,k^i_{sector}`$ are given below. The NSNS-sector is
$`|D2,\gamma ,k_X^i_{NSNS}`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}\{(\alpha _X)_n^0(\stackrel{~}{\alpha }_X)_n^0+{\displaystyle \underset{j=3}{\overset{7}{}}}(\alpha _X)_n^j(\stackrel{~}{\alpha }_X)_n^j\}\right]`$ (3.2)
$`\times `$ $`\mathrm{exp}[i\gamma {\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\{\chi _{n+\frac{1}{2}}^0\stackrel{~}{\chi }_{n+\frac{1}{2}}^0+{\displaystyle \underset{j=3}{\overset{7}{}}}\chi _{n+\frac{1}{2}}^j\stackrel{~}{\chi }_{n+\frac{1}{2}}^j\}]`$
$`\times `$ $`\mathrm{exp}[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{1}{n}}(\alpha _X)_n^i(\stackrel{~}{\alpha }_X)_n^i]\mathrm{exp}[+i\gamma {\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}\chi _{n+\frac{1}{2}}^i\stackrel{~}{\chi }_{n+\frac{1}{2}}^i]`$
$`\times |D2,\gamma ,k_X^i_{NSNS}^{(0)},`$
and zero-mode
$`|D2,\gamma ,k_X^i_{NSNS}^{(0)}`$ $`=`$ $`{\displaystyle \underset{\stackrel{}{w}_X𝐙^2}{}}|\stackrel{}{0},\stackrel{}{w}_X|k_X^i|\mathrm{\Omega }_{NSNS}^{(0)},`$ (3.3)
where $`|\mathrm{\Omega }_{NSNS}^{(0)}`$ is the vacuum of world sheet theory and $`|\stackrel{}{n}_X,\stackrel{}{w}_X`$ represents the zero mode part of $`T^2`$ which has momenta $`\stackrel{}{n}_X=(n_X^1,n_X^2)`$ and windings $`\stackrel{}{w}_X=(w_X^1,w_X^2)`$. The RR-sector is given as
$`|D2,\gamma ,k_X^i_{RR}`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}\{(\alpha _X)_n^0(\stackrel{~}{\alpha }_X)_n^0+{\displaystyle \underset{j=3}{\overset{7}{}}}(\alpha _X)_n^j(\stackrel{~}{\alpha }_X)_n^j\}\right]`$ (3.4)
$`\times `$ $`\mathrm{exp}[i\gamma {\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\{\chi _n^0\stackrel{~}{\chi }_n^0+{\displaystyle \underset{j=3}{\overset{7}{}}}\chi _n^j\stackrel{~}{\chi }_n^j\}]`$
$`\times `$ $`\mathrm{exp}[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{1}{n}}(\alpha _X)_n^i(\stackrel{~}{\alpha }_X)_n^i]\mathrm{exp}[+i\gamma {\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}\chi _n^i\stackrel{~}{\chi }_n^i]`$
$`\times |D2,\gamma ,k_X^i_{RR}^{(0)},`$
and zero-mode
$`|D2,\gamma ,k_X^i_{RR}^{(0)}`$ $`=`$ $`{\displaystyle \underset{\stackrel{}{w}_XZ^2}{}}|\stackrel{}{0},\stackrel{}{w}_X|k_X^i|\mathrm{\Omega },\gamma _{RR}^{(0)},`$ (3.5)
where $`|\mathrm{\Omega },\gamma _{RR}^{(0)}`$ is the solution to
$`\{\chi _0^\mu i\gamma \chi _0^\mu \}|\mathrm{\Omega },\gamma _{RR}^{(0)}=\{\chi _0^i+i\gamma \chi _0^i\}|\mathrm{\Omega },\gamma _{RR}^{(0)}=0(\mu =0,1,2).`$ (3.6)
The normalizations of the zero modes are defined as
$`\mathrm{\Omega }|\mathrm{\Omega }_{NSNS}^{(0)}=1,`$ $`\mathrm{\Omega },\gamma |\mathrm{\Omega },\gamma ^{}_{RR}^{(0)}=\delta _{\gamma ,\gamma ^{}},`$
$`\stackrel{}{n}_X,\stackrel{}{w}_X|\stackrel{}{n^{}}_X,\stackrel{}{w^{}}_X=V\delta _{n,n^{}}\delta _{w,w^{}},`$ $`k_X^i|k_X^i=(2\pi )^7\delta ^7(kk^{}),`$ (3.7)
where $`V`$ is the volume of the time direction. Note that this is the solution to the condition (A.2) for a boundary state of a D2-brane. We use $`\gamma =\pm 1`$ to indicate each choice of the open string boundary condition (A.2). Also note that since we have used the light cone formalism , the superscripts of oscillators run from 0 to 7, not to 9. We have divided the oscillator parts into $`X^0,X^3X^7`$ and $`X^1,X^2`$ ($`\chi ^\mu `$ are also divided) in eq(3.2),(3.4). This is because $`X^0,X^3X^7`$ part does not contribute to the later calculations of tachyon condensation importantly. Therefore we abbreviate $`X^0,X^3X^7,\chi ^0,\chi ^3\chi ^7`$ part and $`k^i`$ from now on.
Next, we construct the boundary state for $`D2\overline{D2}`$ system where the position of the $`D2`$-brane and the $`\overline{D2}`$-brane is $`x^i=0`$. This is given by the superposition of the boundary states for a D2-brane as
$`\{\begin{array}{ccc}\hfill |D2\overline{D2},\gamma _{NSNS}& =& |D2,\gamma _{NSNS}+|D2^{},\gamma _{NSNS},\hfill \\ \hfill |D2\overline{D2},\gamma _{RR}& =& |D2,\gamma _{RR}|D2^{},\gamma _{RR}.\hfill \end{array}`$ (3.10)
Here, we have two important points. One point is that we have switched on $`𝐙_2`$ Wilson lines of the second D2-brane and we have expressed such a boundary state as that with a prime. Next point is the second D2-brane is the anti D-brane. Since an anti D-brane has the opposite RR charge to a D-brane, we have added a minus sign to the second boundary state of RR sector.
From eq.(3.10), the oscillator parts of $`|D2\overline{D2},\gamma >_{NSNS,RR}`$ are the same as eq.(3.2), eq.(3.4), and the zero mode parts are given by
$`|D2\overline{D2},\gamma _{NSNS}^{(0)}`$ $`=`$ $`{\displaystyle \underset{\stackrel{}{w}_X𝐙^2}{}}|\stackrel{}{0},\stackrel{}{w}_X_{NSNS}^{(0)}+(1)^{w_X^1+w_X^2}|\stackrel{}{0},\stackrel{}{w}_X_{NSNS}^{(0)},`$
$`|D2\overline{D2},\gamma _{RR}^{(0)}`$ $`=`$ $`{\displaystyle \underset{\stackrel{}{w}_X𝐙^2}{}}|\stackrel{}{0},\stackrel{}{w}_X,\gamma _{RR}^{(0)}(1)^{w_X^1+w_X^2}|\stackrel{}{0},\stackrel{}{w}_X,\gamma _{RR}^{(0)},`$ (3.11)
where the phase factors $`(1)^{w_X^1},(1)^{w_X^2}`$ are due to $`Z_2`$ Wilson lines .
Then let us perform $`\frac{\pi }{4}`$ rotation (2.4) and rewrite<sup>14</sup><sup>14</sup>14Note that here we have regarded the radii of $`Y^1,Y^2`$ direction as $`R^1=R^2=\frac{1}{\sqrt{2}}`$ and we have used the relations $`w_Y^1=w_X^1+w_X^2,w_Y^2=w_X^1w_X^2`$. the boundary state $`|D2\overline{D2},\gamma _{NSNS,RR}`$ by using the fields $`(Y^1,Y^2,\psi ^1,\psi ^2)`$ as follows
$`|D2\overline{D2},\gamma _{NSNS}`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{1}{n}}(\alpha _Y)_n^i(\stackrel{~}{\alpha }_Y)_n^i\right]\mathrm{exp}\left[+i\gamma {\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}\psi _{n+\frac{1}{2}}^i\stackrel{~}{\psi }_{n+\frac{1}{2}}^i\right]`$ (3.12)
$`\times 2{\displaystyle \underset{\stackrel{}{w}_Y𝐙^2}{}}|\stackrel{}{0},2\stackrel{}{w}_Y_{NSNS}^{(0)},`$
$`|D2\overline{D2},\gamma _{RR}`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{1}{n}}(\alpha _Y)_n^i(\stackrel{~}{\alpha }_Y)_n^i\right]\mathrm{exp}\left[+i\gamma {\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}\psi _n^i\stackrel{~}{\psi }_n^i\right]`$ (3.13)
$`\times 2{\displaystyle \underset{\stackrel{}{w}_Y𝐙^2}{}}|\stackrel{}{0},2\stackrel{}{w}_Y+\stackrel{}{1},\gamma _{RR}^{(0)},`$
where we defined $`\stackrel{}{1}=(1,1)`$.
As explained in section 2, we have to change the base of the world sheet fields $`(Y,\psi )`$ into the bosonized ones $`(\varphi ,\eta )`$ in order to describe the tachyon condensation. Then by using $`(\varphi ,\eta )`$ modes how are the boundary states (3.12),(3.13) represented? Generalizing the discussion in , we argue that the following boundary states are equivalent to eq.(3.12) and (3.13) respectively for $`\gamma =+1`$:
$`|D2\overline{D2},+_{NSNS}`$ $`=`$ $`2\mathrm{exp}\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{1}{n}}\varphi _n^i\stackrel{~}{\varphi }_n^i\right]\mathrm{exp}\left[+i{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}\eta _{n+\frac{1}{2}}^i\stackrel{~}{\eta }_{n+\frac{1}{2}}^i\right]`$ (3.14)
$`\times {\displaystyle \underset{\stackrel{}{w}_\varphi 𝐙^2}{}}|\stackrel{}{0},2\stackrel{}{w}_\varphi ^{(0)}_{NSNS},`$
$`|D2\overline{D2},+_{RR}`$ $`=`$ $`2\mathrm{exp}\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{1}{n}}\varphi _n^i\stackrel{~}{\varphi }_n^i\right]\mathrm{exp}\left[+i{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}\eta _n^i\stackrel{~}{\eta }_n^i\right]`$ (3.15)
$`\times {\displaystyle \underset{\stackrel{}{w}_\varphi 𝐙^2}{}}|\stackrel{}{0},2\stackrel{}{w}_\varphi +\stackrel{}{1},+^{(0)}_{RR}.`$
These are obtained simply by replacing $`(\alpha _Y)_n^i,\psi _n^i,w_X^i`$ in eq.(3.12),(3.13) with $`\varphi _n^i,\eta _n^i,w_\varphi ^i`$. In the appendix B we show they indeed satisfy the desirable boundary conditions if we take the detailed (second-type) cocycle factors into consideration :
$`_2Y^i(w,\overline{w})|_{\sigma _2=0}|D2\overline{D2},+_{NSNS,RR}=0,`$ (3.16)
$`(\psi _R^i(w)i\psi _L^i(\overline{w}))|_{\sigma _2=0}|D2\overline{D2},+_{NSNS,RR}=0.`$ (3.17)
Note that eq.(3.16) and (3.17) are not enough<sup>15</sup><sup>15</sup>15This is because these constraints do not determine the detailed structures of the zero modes such as Wilson lines. for the proof of the equivalence. But further we can see that for several closed string states eq.(3.14) and (3.15) have the same overlap as eq.(3.12) and (3.13) do <sup>16</sup><sup>16</sup>16 This is almost the same calculation as that in appendix B of . Thus we omit its detail.. The third evidence is the equality of their partition functions, which we will see in the next subsection. We propose that these three evidences are enough for the proof of the equivalence.
On the other hand, $`|D2\overline{D2},_{NSNS,RR}`$ is given by acting left-moving fermion number operator $`(1)^{\stackrel{~}{F}_Y}`$ of $`(Y,\psi )`$ system :
$`|D2\overline{D2},_{NSNS,RR}=(1)^{\stackrel{~}{F}_Y}|D2\overline{D2},+_{NSNS,RR}.`$ (3.18)
Note that the action of $`(1)^{\stackrel{~}{F}_Y}`$ is given as
$`(1)^{\stackrel{~}{F}_Y}:`$ $`(\stackrel{~}{\alpha }_Y)_n^i(\stackrel{~}{\alpha }_Y)_n^i,\stackrel{~}{\psi }_n^i\stackrel{~}{\psi }_n^i,`$ (3.19)
$`\stackrel{~}{\varphi }_n^i\stackrel{~}{\varphi }_n^i,\stackrel{~}{\eta }_n^i\stackrel{~}{\eta }_n^i,`$
and to zero mode
$`(1)^{\stackrel{~}{F}_Y}:|\stackrel{}{n}_\varphi ,\stackrel{}{w}_\varphi (\text{phase})\times |{\displaystyle \frac{\stackrel{}{w}_\varphi }{2}},2\stackrel{}{n}_\varphi ,`$ (3.20)
where ”phase” comes from cocycle factors. Therefore for example, $`|D2\overline{D2},_{NS}`$ is given by
$`|D2\overline{D2},_{NSNS}`$ $`=`$ $`2\mathrm{exp}\left[+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{1}{n}}\varphi _n^i\stackrel{~}{\varphi }_n^i\right]\mathrm{exp}\left[+i{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}\eta _{n+\frac{1}{2}}^i\stackrel{~}{\eta }_{n+\frac{1}{2}}^i\right]`$ (3.21)
$`\times {\displaystyle \underset{\stackrel{}{w}_\varphi 𝐙^2}{}}(1)^{w_\varphi ^1+w_\varphi ^2}|\stackrel{}{w}_\varphi ,\stackrel{}{0}_{NSNS}^{(0)}.`$
Then it is straightforward to condense the tachyon by using the Wilson lines (2.9). Since closed strings usually do not have Chan-Paton factors, the tachyon condensation in the closed string viewpoint corresponds to the insertion of the trace of (2.9) in front of the boundary state and the trace is given as
$`W_1(\alpha ,\beta )\mathrm{cos}\left({\displaystyle \frac{\pi \alpha w_\varphi ^1}{2}}\right)\mathrm{cos}\left({\displaystyle \frac{\pi \beta w_\varphi ^2}{2}}\right).`$ (3.22)
But, this is not sufficient. In the authors argued that the Wess-Zumino terms in the effective action of $`D2\overline{D2}`$ systems should possess the following coupling.
$`{\displaystyle C_1}dTd\overline{T},`$ (3.23)
where $`T,\overline{T}`$ are complex tachyon fields, and $`C_1`$ is the R-R 1-form which couples to a D0-brane. This implies that $`C_1`$ have Chan-Paton factor $`\sigma _3`$ since $`T,\overline{T}`$ have Chan-Paton factor $`\sigma _1\pm i\sigma _2`$ respectively. At first sight you may think such an idea is not acceptable, but it is required to reproduce the correct spectrum of open strings and satisfy the Cardy’s condition as we will see in the next subsection. Therefore we regard our results as another evidence of such an idea. Similar “Chan-Paton factor for closed string vertex” was discussed in the case of a non-BPS D-brane<sup>17</sup><sup>17</sup>17For example, in the case of the non-BPS D2-brane the Wess-Zumino term is written as $`C_1dT`$ ., where it was argued that the branch cut due to a RR-vertex provides an extra Chan-Paton factor if its one end is on the non-BPS D-brane. We also argue that some of the states in NSNS sector have the Chan-Paton factor $`\sigma _3`$. This fact can also be verified by the Cardy’s condition and will be required due to the supersymmetry of the bulk theory. Therefore in these sectors we should insert $`\sigma _3`$ in the trace.
Then the tachyon condensation switches not only (3.22) but also
$`W_{\sigma _3}(\alpha ,\beta )\mathrm{sin}\left({\displaystyle \frac{\pi \alpha w_\varphi ^1}{2}}\right)\mathrm{sin}\left({\displaystyle \frac{\pi \beta w_\varphi ^2}{2}}\right),`$ (3.24)
which is obtained by inserting $`\sigma _3`$ in the trace.
Then by switching both contributions we obtain the following boundary state <sup>18</sup><sup>18</sup>18Strictly speaking, from the above explanation we can’t decide the relative normalization between the first and the second term. This is determined by the vacuum energy calculation in the next subsection.
$`|B(\alpha ,\beta ),+_{NSNS}`$ $`=`$ $`2\mathrm{exp}\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{1}{n}}\varphi _n^i\stackrel{~}{\varphi }_n^i\right]\mathrm{exp}\left[+i{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}\eta _{n+\frac{1}{2}}^i\stackrel{~}{\eta }_{n+\frac{1}{2}}^i\right]`$ (3.25)
$`\times {\displaystyle \underset{\stackrel{}{w}_\varphi 𝐙^2}{}}[\mathrm{cos}(\pi \alpha w_\varphi ^1)\mathrm{cos}(\pi \beta w_\varphi ^2)`$
$`+\mathrm{sin}(\pi \alpha w_\varphi ^1)\mathrm{sin}(\pi \beta w_\varphi ^2)]|\stackrel{}{0},2\stackrel{}{w}_\varphi _{NSNS}^{(0)},`$
$`|B(\alpha ,\beta ),+_{RR}`$ $`=`$ $`2\mathrm{exp}\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{1}{n}}\varphi _n^i\stackrel{~}{\varphi }_n^i\right]\mathrm{exp}\left[+i{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}\eta _n^i\stackrel{~}{\eta }_n^i\right]`$
$`\times {\displaystyle \underset{\stackrel{}{w}_\varphi 𝐙^2}{}}[\mathrm{cos}\{\pi \alpha (w_\varphi ^1+{\displaystyle \frac{1}{2}})\}\mathrm{cos}\{\pi \beta (w_\varphi ^2+{\displaystyle \frac{1}{2}})\}`$
$`+\mathrm{sin}\{\pi \alpha (w_\varphi ^1+{\displaystyle \frac{1}{2}})\}\mathrm{sin}\{\pi \beta (w_\varphi ^2+{\displaystyle \frac{1}{2}})\}]|\stackrel{}{0},2\stackrel{}{w}_\varphi +\stackrel{}{1},+_{RR}^{(0)}.`$
This satisfies $`|B(0,0),+_{NSNS,RR}=|D2\overline{D2},+_{NSNS,RR}`$.
As explained in section 2.1 the point $`\alpha =\beta =1`$ is expected to be identified as a $`D0\overline{D0}`$ system where a D0-brane and a $`\overline{D0}`$-brane are produced at $`(x_1,x_2)=(0,0),(\pi ,\pi )`$ respectively. The boundary state of this $`D0\overline{D0}`$ system is
$`|D0\overline{D0},+_{NSNS}`$ $`=`$ $`2\mathrm{exp}\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{1}{n}}(\alpha _Y)_n^i(\stackrel{~}{\alpha }_Y)_n^i\right]\mathrm{exp}\left[i{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}\psi _{n+\frac{1}{2}}^i\stackrel{~}{\psi }_{n+\frac{1}{2}}^i\right]`$ (3.27)
$`\times {\displaystyle \underset{\stackrel{}{n}_Y𝐙^2}{}}|\stackrel{}{n}_Y,\stackrel{}{0}^{(0)}_{NSNS},`$
$`|D0\overline{D0},+_{RR}`$ $`=`$ $`2\mathrm{exp}\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{1}{n}}(\alpha _Y)_n^i(\stackrel{~}{\alpha }_Y)_n^i\right]\mathrm{exp}\left[i{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}\psi _n^i\stackrel{~}{\psi }_n^i\right]`$ (3.28)
$`\times {\displaystyle \underset{\stackrel{}{n}_Y𝐙^2}{}}|\stackrel{}{n}_Y+\stackrel{}{{\displaystyle \frac{1}{2}}},\stackrel{}{0},+^{(0)}_{RR}.`$
On the other hand, at this point the zero mode parts of eq.(3.25) and (3.1) become respectively
$`{\displaystyle \underset{\stackrel{}{w}_\varphi 𝐙^2}{}}(1)^{w_\varphi ^1+w_\varphi ^2}|\stackrel{}{0},2\stackrel{}{w}_\varphi _{NSNS}^{(0)},`$
$`{\displaystyle \underset{\stackrel{}{w}_\varphi 𝐙^2}{}}(1)^{w_\varphi ^1+w_\varphi ^2}|\stackrel{}{0},2\stackrel{}{w}_\varphi +\stackrel{}{1}_{RR}^{(0)}.`$ (3.29)
Then just as we have verified that eq.(3.14) and (3.15) are equivalent to eq. (3.12),(3.13), we can also verify that $`|B(1,1),+_{NSNS,RR}`$ is equivalent to eq.(3.27),(3.28) in the same way. For example, eq.(3.27),(3.28) indeed satisfy the following equations which represent the boundary conditions of D0-branes :
$`_1Y^i(w)|_{\sigma _2=0}|B(1,1),+_{NSNS,RR}`$ $`=`$ $`0,`$ (3.30)
$`\left(\psi _R^i(w)+i\psi _L^i(\overline{w})\right)|_{\sigma _2=0}|B(1,1),+_{NSNS,RR}`$ $`=`$ $`0,(i=1,2).`$ (3.31)
In other words the tachyon condensation from $`\alpha =\beta =0`$ to $`\alpha =\beta =1`$ changes the boundary conditions (3.16),(3.17) into (3.30),(3.31) and the crucial difference between them is that the latter has the phase factor $`(1)^{w_\varphi ^1+w_\varphi ^2}`$. At $`\alpha =\beta =0`$ only the first term of eq.(3.1) is nonzero and this corresponds to the RR charge of D2-brane. As the tachyon is condensed the second term also ceases to be zero and this means<sup>19</sup><sup>19</sup>19It is easy to see that if $`\alpha ,\beta `$ are small, then the second term is proportional to $`V_{T1}V_{T2}|D2_{RR}|D0_{RR}`$. that the RR charge of the D0-brane is generated. Finally at $`\alpha =\beta =1`$ only the second term is nonzero and this is the pure D0-brane RR charge. Note that if we ignored the factor (3.24) which corresponds to $`\sigma _3`$ sector, then the RR-sector boundary state would vanish at $`\alpha =\beta =1`$ and be inconsistent. In this way we see explicitly in the closed string formalism that a tachyon kink on a brane-antibrane system produces a codimension two D-brane (see Figure 1).
Let us turn to the other points of $`\alpha ,\beta `$. It is easy to see that at $`(\alpha ,\beta )=(0,1),(1,0)`$ the RR-sector boundary state does vanish and each system corresponds to a non-BPS D1-brane stretching along the direction of $`Y^1`$ or $`Y^2`$ respectively (see Figure 1). Physically this can be interpreted as the statement that a tachyon kink produces a codimension one (non-BPS) D-brane. All of these identifications will be verified further by the calculation of vacuum amplitudes including the detailed normalization.
### 3.2 Calculation of the vacuum amplitude
Here we calculate the vacuum amplitude of $`D2\overline{D2}`$ system for every value of $`\alpha ,\beta `$ and translate it from the viewpoint of open string. As a result it will be shown that the boundary state have the correct normalizations or equally correct NSNS and RR-charge needed for the identification and that the additional NSNS and RR sector discussed before are indeed required in order to satisfy the Cardy’s condition.
First let us define the propagator for closed string as
$$\mathrm{\Delta }=\frac{1}{2}_0^{\mathrm{}}𝑑se^{sH_c}\frac{1}{k^2}+\mathrm{},$$
(3.32)
where $`H_c`$ denotes the closed string Hamiltonian and its explicit form is given as
$`H_c`$ $`=`$ $`{\displaystyle \underset{i=1,2}{}}\left({\displaystyle \frac{(n_\varphi ^i)^2}{2R_i^2}}+{\displaystyle \frac{1}{2}}(w_\varphi ^i)^2R_i^2\right)+{\displaystyle \underset{i=1,2}{}}\{{\displaystyle \underset{n}{}}(\varphi _n^i\varphi _n^i+\stackrel{~}{\varphi }_n^i\stackrel{~}{\varphi }_n^i)+{\displaystyle \underset{r}{}}(\eta _r^i\eta _r^i+\stackrel{~}{\eta }_r^i\stackrel{~}{\eta }_r^i)\}`$ (3.33)
$`+{\displaystyle \underset{i=3}{\overset{9}{}}}{\displaystyle \frac{1}{2}}(k^i)^2+{\displaystyle \underset{i=0,3}{\overset{7}{}}}\{{\displaystyle \underset{n}{}}(\alpha _n^i\alpha _n^i+\stackrel{~}{\alpha }_n^i\stackrel{~}{\alpha }_n^i)+{\displaystyle \underset{r}{}}(\chi _r^i\chi _r^i+\stackrel{~}{\chi }_r^i\stackrel{~}{\chi }_r^i)\}+a,`$
where $`a`$ denotes the zero-energy for each sector and is given as $`a=1`$ for NSNS-sector and $`a=0`$ for RR-sector.
Then the vacuum amplitudes for NSNS and RR sector are
$`Z_{NSNS}`$ $`=`$ $`{\displaystyle \frac{(T_{p=2})^2}{16}}{\displaystyle _0^{\mathrm{}}}ds{\displaystyle }\left({\displaystyle \frac{dk}{2\pi }}\right)^7\left({\displaystyle \frac{dk^{}}{2\pi }}\right)^7\{B(\alpha ,\beta ),+,k|\mathrm{\Delta }|B(\alpha ,\beta ),+,k^{}_{NSNS}`$
$`B(\alpha ,\beta ),+,k|\mathrm{\Delta }|B(\alpha ,\beta ),,k^{}_{NSNS}\},`$
$`=`$ $`{\displaystyle \frac{(T_{p=2})^2V_{D2}}{4}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{1}{(2\pi s)^{\frac{7}{2}}}}[{\displaystyle \underset{w_\varphi ^1,w_\varphi ^2}{}}\{\mathrm{cos}^2(\pi \alpha w_\varphi ^1)\mathrm{cos}^2(\pi \beta w_\varphi ^2)`$
$`+\mathrm{sin}^2(\pi \alpha w_\varphi ^1)\mathrm{sin}^2(\pi \beta w_\varphi ^2)\}q^{(w_\varphi ^1)^2+(w_\varphi ^2)^2}{\displaystyle \frac{f_3(q)^8}{f_1(q)^8}}2{\displaystyle \frac{f_4(q)^6f_3(q)^2}{f_1(q)^6f_2(q)^2}}],`$
$`Z_{RR}`$ $`=`$ $`(T_{p=2})^2{\displaystyle _0^{\mathrm{}}}𝑑s{\displaystyle \left(\frac{dk}{2\pi }\right)^7\left(\frac{dk^{}}{2\pi }\right)^7B(\alpha ,\beta ),+,k|\mathrm{\Delta }|B(\alpha ,\beta ),+,k^{}_{RR}},`$ (3.34)
$`=`$ $`{\displaystyle \frac{(T_{p=2})^2V_{D2}}{4}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{1}{(2\pi s)^{\frac{7}{2}}}}[{\displaystyle \underset{w_\varphi ^1,w_\varphi ^2}{}}\{\mathrm{cos}^2(\pi \alpha (w_\varphi ^1+{\displaystyle \frac{1}{2}}))\mathrm{cos}^2(\pi \beta (w_\varphi ^2+{\displaystyle \frac{1}{2}}))`$
$`+\mathrm{sin}^2(\pi \alpha (w_\varphi ^1+{\displaystyle \frac{1}{2}}))\mathrm{sin}^2(\pi \beta (w_\varphi ^2+{\displaystyle \frac{1}{2}}))\}q^{(w_\varphi ^1+\frac{1}{2})^2+(w_\varphi ^2+\frac{1}{2})^2}{\displaystyle \frac{f_2(q)^8}{f_1(q)^8}}],`$
where $`V_{D2}=(2\pi )^2V`$ is the volume of D2-brane and we defined
$`f_1(q)=q^{\frac{1}{12}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1q^{2n})`$ , $`f_2(q)=\sqrt{2}q^{\frac{1}{12}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1+q^{2n}),`$
$`f_3(q)=q^{\frac{1}{24}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1+q^{2n1})`$ , $`f_4(q)=q^{\frac{1}{24}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1q^{2n1}),`$ (3.35)
with $`q=e^s`$.
Next let us perform modular transformations and interpret this as the open string cylinder amplitude. We define the modulus of the cylinder as $`t=\frac{\pi }{s}`$ and introduce $`\stackrel{~}{q}=e^{\pi t}`$. Then we get the following open string amplitude
$`Z_{NS}^{open}`$ $`=`$ $`2^{\frac{3}{2}}\pi ^1V{\displaystyle _0^{\mathrm{}}}dtt^{\frac{3}{2}}{\displaystyle \underset{n_1,n_2}{}}[{\displaystyle \frac{1}{2}}(\stackrel{~}{q}^{n_1^2+n_2^2}+\stackrel{~}{q}^{(n_1\alpha )^2+(n_2\beta )^2)}){\displaystyle \frac{f_3(\stackrel{~}{q})^8}{f_1(\stackrel{~}{q})^8}}`$
$`{\displaystyle \frac{1}{2}}(1)^{n_1+n_2}(\stackrel{~}{q}^{n_1^2+n_2^2}+\stackrel{~}{q}^{(n_1\alpha )^2+(n_2\beta )^2}){\displaystyle \frac{f_4(\stackrel{~}{q})^8}{f_1(\stackrel{~}{q})^8}}],`$
$`Z_R^{open}`$ $`=`$ $`2^{\frac{3}{2}}\pi ^1V{\displaystyle _0^{\mathrm{}}}𝑑tt^{\frac{3}{2}}{\displaystyle \underset{n_1,n_2}{}}\stackrel{~}{q}^{n_1^2+n_2^2}{\displaystyle \frac{f_2(\stackrel{~}{q})^8}{f_1(\stackrel{~}{q})^8}},`$ (3.36)
where we used the following identities
$`{\displaystyle \underset{n}{}}q^{n^2}=f_1(q)f_3(q)^2`$ , $`{\displaystyle \underset{n}{}}(1)^nq^{n^2}=f_1(q)f_4(q)^2,`$
$`{\displaystyle \underset{n}{}}q^{(n\frac{1}{2})^2}=f_1(q)f_2(q)^2`$ , $`f_2(q)f_3(q)f_4(q)=\sqrt{2},`$ (3.37)
and the modular properties
$`f_1(e^{\frac{\pi }{t}})=\sqrt{t}f_1(e^{\pi t})`$ , $`f_2(e^{\frac{\pi }{t}})=f_4(e^{\pi t}),`$
$`f_3(e^{\frac{\pi }{t}})=f_3(e^{\pi t})`$ , $`f_4(e^{\frac{\pi }{t}})=f_2(e^{\pi t}).`$ (3.38)
Now it is obvious that for each value of $`\alpha ,\beta `$ the open string spectrum is well defined only if we incorporate the additional sector of the boundary state defined in the previous subsection, otherwise the number of open string states for given $`n_1,n_2,H_c`$ would be fractional. This fact will be more clear if we note that this amplitude can be rewritten as
$`Z`$ $`=`$ $`Z_{NS}^{open}+Z_R^{open},`$ (3.39)
$`=`$ $`V_{D2}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t}}Tr_{NSR}\left[{\displaystyle \frac{1+(1)^{F^\varphi }h_1^\varphi h_2^\varphi }{2}}\stackrel{~}{q}^{2H_o}\right],`$
where $`Tr`$ is the trace over the open string Hilbert space including zero-modes and Chan-Paton sectors. Also $`H_o`$ means the open string Hamiltonian and is given as follows
$`H_o`$ $`=`$ $`(p^0)^2+{\displaystyle \underset{i=1,2}{}}R_{i}^{}{}_{}{}^{2}(w_\varphi ^i)^2+{\displaystyle \underset{i=1,2}{}}\{{\displaystyle \underset{n}{}}\varphi _n^i\varphi _n^i+{\displaystyle \underset{r}{}}\eta _r^i\eta _r^i\}`$ (3.40)
$`+{\displaystyle \underset{i=0,3}{\overset{7}{}}}\{{\displaystyle \underset{n}{}}\alpha _n^i\alpha _n^i+{\displaystyle \underset{r}{}}\chi _r^i\chi _r^i\}+a,`$
where $`a`$ denotes the zero-energy and is given as $`a=\frac{1}{2}`$ for NS-sector and $`a=0`$ for R-sector.
This physically important constraint (3.39) is generally called Cardy’s condition . Also notice that the above open string spectrum is consistent with the momentum shift shown in Table 1.
Finally let us verify the identification at particular $`\alpha ,\beta `$. In the case of $`(\alpha ,\beta )=(1,0)`$ or $`(0,1)`$ we get after the modular transformations
$`Z_{\alpha =1,\beta =0}=Z_{\alpha =0,\beta =1}`$ $`=`$ $`2^{\frac{3}{2}}\pi ^1V{\displaystyle _0^{\mathrm{}}}𝑑tt^{\frac{3}{2}}{\displaystyle \underset{n_1,n_2}{}}\stackrel{~}{q}^{n_1^2+n_2^2}{\displaystyle \frac{f_3(\stackrel{~}{q})^8f_2(\stackrel{~}{q})^8}{f_1(\stackrel{~}{q})^8}},`$ (3.41)
$`=`$ $`2\pi \sqrt{2}V{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t}}Tr_{NSR}(\stackrel{~}{q}^{2H_o}).`$
Therefore we can identify the system as a non-BPS D1-brane of which length is $`2\sqrt{2}\pi `$ as expected. Another case is $`(\alpha ,\beta )=(1,1)`$ and the amplitude can be written as
$`Z_{\alpha =1,\beta =1}`$ $`=`$ $`2^{\frac{3}{2}}\pi ^1V{\displaystyle _0^{\mathrm{}}}dtt^{\frac{3}{2}}{\displaystyle \underset{n_1,n_2}{}}[\stackrel{~}{q}^{n_1^2+n_2^2}{\displaystyle \frac{f_3(\stackrel{~}{q})^8f_2(\stackrel{~}{q})^8}{f_1(\stackrel{~}{q})^8}}(1)^{n_1+n_2}\stackrel{~}{q}^{n_1^2+n_2^2}{\displaystyle \frac{f_4(\stackrel{~}{q})^8}{f_1(\stackrel{~}{q})^8}}],`$
$`=`$ $`2V{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t}}{\displaystyle \frac{1}{\sqrt{8\pi ^2t}}}{\displaystyle \underset{m_1,m_2}{}}[\stackrel{~}{q}^{2m_1^2+2m_2^2}{\displaystyle \frac{f_3(\stackrel{~}{q})^8f_4(\stackrel{~}{q})^8f_2(\stackrel{~}{q})^8}{2f_1(\stackrel{~}{q})^8}}`$
$`+\stackrel{~}{q}^{2(m_1+\frac{1}{2})^2+2(m_2+\frac{1}{2})^2}{\displaystyle \frac{f_3(\stackrel{~}{q})^8+f_4(\stackrel{~}{q})^8f_2(\stackrel{~}{q})^8}{2f_1(\stackrel{~}{q})^8}}],`$
$`=`$ $`2V{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t}}Tr_{NSR}\left[{\displaystyle \frac{1+(1)^F}{2}}\stackrel{~}{q}^{2H_c}\right]+2V{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t}}Tr_{NSR}\left[{\displaystyle \frac{1(1)^F}{2}}\stackrel{~}{q}^{2H_c}\right],`$
where we have defined $`m_1=\frac{n_1+n_2}{2},m_2=\frac{n_1n_2}{2}`$. This shows explicitly that the system is equivalent to a D0-brane and an anti D0-brane which are separated from each other by $`\mathrm{\Delta }x_1=\mathrm{\Delta }x_2=\pi `$.
### 3.3 Moduli space of non-supersymmetric D-branes
So far we have discussed the boundary states which describe various tachyon condensations in $`D2\overline{D}2`$ system at the critical radii. The tachyon condensations are parameterized by $`\alpha ,\beta `$ which have periodicity $`\alpha \alpha +2,\beta \beta +2`$. At this particular radius the non-supersymmetric D-brane configrations for all values of $`\alpha ,\beta `$ are realized in conformal invariant manners. Figure 2 shows the moduli space of the non-supersymmetric D-brane configrations. In particular $`(\alpha ,\beta )=(0,0)(1,0)(1,1)`$ can be regarded as a continuous version of the descent relation (see Figure 1).
Realistically we are interested in the $`D2\overline{D2}`$ system at generic radii. If we shift the radius, tadpoles develop as can be seen by the method discussed in or by computing one point functions using the boundary state we constructed. The tadpoles only vanish at $`\mathrm{sin}(\pi \alpha )=\mathrm{sin}(\pi \beta )=0`$, which correspond to $`D2\overline{D2}`$, non-BPS $`D1`$-brane and $`D0\overline{D0}`$.
### 3.4 Tachyon condensation in general $`Dp\overline{Dp}`$ systems
In this subsection, we generalize our construction to the higher codimension cases. It is enough to consider the $`D8\overline{D8}`$ system (In odd codimension case, we have only to consider the decay to the $`D(1)\overline{D(1)}`$ system.). The different points from the analysis of $`D2\overline{D2}`$ system are the slightly complicated choice of the gamma matrices and Chan-Paton factors. We mainly concentrate on these issues in the presentation here.
In particular, we deal with the case where $`2^3`$ $`D8\overline{D8}`$ pairs at critical radii $`(R_1=\mathrm{}=R_8=1)`$ become $`2^3`$ $`D0\overline{D0}`$ pairs via the tachyon condensation as the marginal deformation. These $`2^3`$ soliton - antisoliton pairs represent $`2^3`$ $`D0\overline{D0}`$ pairs, and a single $`D0`$ brane is identified with the single codimension 8 soliton on the $`D8\overline{D8}`$ pair at the point where the tachyon condensation is maximum. Many other points are identified with the bound states of $`D8\overline{D8},D6\overline{D6},D4\overline{D4},D2\overline{D2},D0\overline{D0}`$ at critical radii. We show the emergence of the RR charges of the lower dimensional $`D`$-branes explicitly.
We switched on $`𝐙_2`$ Wilson line to the second anti-brane using the Wilson lines
$`X^1X^2\mathrm{direction}\sigma _3\sigma _311,`$ $`X^3X^4\mathrm{direction}\sigma _3111,`$
$`X^5X^6\mathrm{direction}1\sigma _3\sigma _31,`$ $`X^7X^8\mathrm{direction}11\sigma _3\sigma _3.`$ (3.43)
where we adopt the following representation of the $`SO(8)`$ Clifford algebra
$`\mathrm{\Gamma }_1=1\sigma _1\sigma _2\sigma _2,\mathrm{\Gamma }_2=\sigma _3\sigma _2\sigma _2\sigma _2,\mathrm{\Gamma }_3=\sigma _1\sigma _2\sigma _2\sigma _2,`$
$`\mathrm{\Gamma }_4=\sigma _2\sigma _2\sigma _2\sigma _2,\mathrm{\Gamma }_5=1\sigma _3\sigma _2\sigma _2,\mathrm{\Gamma }_6=11\sigma _1\sigma _2,`$
$`\mathrm{\Gamma }_7=11\sigma _3\sigma _2,\mathrm{\Gamma }_8=111\sigma _1.`$ (3.44)
Above matrices for the Wilson lines have the properties such that the matrix in the $`X^{2k+1}X^{2k+2}`$ direction anticommutes with $`\mathrm{\Gamma }_{2k+1},\mathrm{\Gamma }_{2k+2}`$, and commutes with the other six gamma matrices. Since we have seen that the oscillator parts don’t contribute crucially the tachyon condensation, we omit these parts. The zero mode part of the boundary states of this $`D8\overline{D8}`$ system is given by
$`|D8\overline{D8},\gamma _{NSNS}^{(0)}`$ $`=`$ $`{\displaystyle \underset{\stackrel{}{w}_X𝐙^8}{}}\left(1+(1)^{w_X^1+w_X^2}\right)\left(1+(1)^{w_X^3+w_X^4}\right)`$
$`\times \left(1+(1)^{w_X^5+w_X^6}\right)\left(1+(1)^{w_X^7+w_X^8}\right)|\stackrel{}{0},\stackrel{}{w}_X_{NSNS}^{(0)},`$
$`|D8\overline{D8},\gamma _{RR}^{(0)}`$ $`=`$ $`{\displaystyle \underset{\stackrel{}{w}_X𝐙^8}{}}\left(1(1)^{w_X^1+w_X^2}\right)\left(1(1)^{w_X^3+w_X^4}\right)`$
$`\times \left(1(1)^{w_X^5+w_X^6}\right)\left(1(1)^{w_X^7+w_X^8}\right)|\stackrel{}{0},\stackrel{}{w}_X,\gamma _{RR}^{(0)}.`$
Now we change the variables as follows
$`Y^{2k+1}={\displaystyle \frac{1}{\sqrt{2}}}\left(X^{2k+1}+X^{2k+2}\right),`$ $`Y^{2k+2}={\displaystyle \frac{1}{\sqrt{2}}}\left(X^{2k+1}X^{2k+2}\right),`$
$`\psi _{R,L}^{2k+1}={\displaystyle \frac{1}{\sqrt{2}}}\left(\chi _{R,L}^{2k+1}+\chi _{R,L}^{2k+2}\right),`$ $`\psi _{R,L}^{2k+2}={\displaystyle \frac{1}{\sqrt{2}}}\left(\chi _{R,L}^{2k+1}\chi _{R,L}^{2k+2}\right),`$
where $`k=0,1,2,3`$. In terms of these variables, the zero mode parts of the boundary state for $`D8\overline{D8}`$ system are rewritten as follows
$`|D8\overline{D8},\gamma _{NSNS}^{(0)}`$ $`=`$ $`16{\displaystyle \underset{\stackrel{}{w}_Y𝐙^8}{}}|\stackrel{}{0},2\stackrel{}{w}_Y_{NSNS}^{(0)},`$ (3.47)
$`|D8\overline{D8},\gamma _{RR}^{(0)}`$ $`=`$ $`16{\displaystyle \underset{\stackrel{}{w}_Y𝐙^8}{}}|\stackrel{}{0},2\stackrel{}{w}_Y+\stackrel{}{1},\gamma _{RR}^{(0)}.`$ (3.48)
where $`w_Y^{2k+1}=w_X^{2k+1}+w_X^{2k+2},w_Y^{2k+2}=w_X^{2k+1}w_X^{2k+2}`$.
In order to describe the effect of the tachyon condensation, we change the variables using bosonization techniques. First, by fermionization $`Y^i`$ are represented by fermions $`\xi ^i,\eta ^i`$ and next, by bosonization we introduce free bosons $`\varphi ^i`$. In the following we neglect the cocycle factors, but we can take their roles into account giving gamma matrix $`\stackrel{~}{\mathrm{\Gamma }}_i`$ to $`e^{i\sqrt{2}Y^i}`$, and $`\stackrel{~}{\mathrm{\Gamma }}_{12345678}`$ to $`\psi ^i,\eta ^i`$ etc. Then the $`D8\overline{D8}`$ system at critical radii is described with the following projection as in section 2.1
$$(1)^Fh_1\mathrm{}h_8=(1)^{F^\varphi }h_1^\varphi \mathrm{}h_8^\varphi =1.$$
(3.49)
Using $`(\varphi ,\eta )`$ modes, we can write down the boundary states of $`D8\overline{D8}`$ system for $`\gamma =+1`$ and their zero modes are given by
$`|D8\overline{D8},+_{NSNS}^{(0)}`$ $`=`$ $`16{\displaystyle \underset{\stackrel{}{w}_\varphi 𝐙^8}{}}|\stackrel{}{0},2\stackrel{}{w}_\varphi _{NSNS}^{(0)},`$ (3.50)
$`|D8\overline{D8},+_{RR}^{(0)}`$ $`=`$ $`16{\displaystyle \underset{\stackrel{}{w}_\varphi 𝐙^8}{}}|\stackrel{}{0},2\stackrel{}{w}_\varphi +\stackrel{}{1},+_{RR}^{(0)}.`$ (3.51)
The equivalence of two states written with different variables is provided by the fact that these states satisfy the definition equation of the boundary state for $`D8\overline{D8}`$ system. Also, the boundary states for $`\gamma =1`$ is given as in the previous subsection.
Now we are ready to condense the tachyon. The tachyon condensation is represented as the insertion of the following Wilson line
$$\mathrm{exp}\left[\underset{i=1}{\overset{8}{}}\frac{i\alpha _i}{2\sqrt{2}}\varphi ^i\mathrm{\Gamma }_i\right],$$
(3.52)
where $`\alpha _i`$ represent the parameters of the condensation and have the periodicity $`\alpha _i\alpha _i+2`$. This represents the marginal deformation at the critical radius. These traces with the insertion of the various Chan-Paton factors are given as the following five types
$`W_1(\{\alpha _i\})`$ $`:`$ $`{\displaystyle \underset{i=1}{\overset{8}{}}}\mathrm{cos}\left({\displaystyle \frac{\pi \alpha _iw_\varphi ^i}{2}}\right),`$
$`W_{\mathrm{\Gamma }_{ij}}(\{\alpha _i\})`$ $`:`$ $`\left[{\displaystyle \underset{i=1}{\overset{6}{}}}\mathrm{cos}\left({\displaystyle \frac{\pi \alpha _iw_\varphi ^i}{2}}\right)\right]\mathrm{sin}\left({\displaystyle \frac{\pi \alpha _7w_\varphi ^7}{2}}\right)\mathrm{sin}\left({\displaystyle \frac{\pi \alpha _8w_\varphi ^8}{2}}\right),{}_{8}{}^{}\mathrm{C}_{2}^{}=28\mathrm{combinations},`$
$`W_{\mathrm{\Gamma }_{ijk\mathrm{}}}(\{\alpha _i\})`$ $`:`$ $`\left[{\displaystyle \underset{i=1}{\overset{4}{}}}\mathrm{cos}\left({\displaystyle \frac{\pi \alpha _iw_\varphi ^i}{2}}\right)\right]\left[{\displaystyle \underset{j=5}{\overset{8}{}}}\mathrm{sin}\left({\displaystyle \frac{\pi \alpha _jw_\varphi ^j}{2}}\right)\right],{}_{8}{}^{}\mathrm{C}_{4}^{}=70\mathrm{combinations},`$
$`W_{\mathrm{\Gamma }_{ijk\mathrm{}mn}}(\{\alpha _i\})`$ $`:`$ $`\mathrm{cos}\left({\displaystyle \frac{\pi \alpha _1w_\varphi ^1}{2}}\right)\mathrm{cos}\left({\displaystyle \frac{\pi \alpha _2w_\varphi ^2}{2}}\right)\left[{\displaystyle \underset{i=3}{\overset{8}{}}}\mathrm{sin}\left({\displaystyle \frac{\pi \alpha _iw_\varphi ^i}{2}}\right)\right],{}_{8}{}^{}\mathrm{C}_{2}^{}=28\mathrm{combinations},`$
$`W_{\mathrm{\Gamma }_{12345678}}(\{\alpha _i\})`$ $`:`$ $`{\displaystyle \underset{i=1}{\overset{8}{}}}\mathrm{sin}\left({\displaystyle \frac{\pi \alpha _iw_\varphi ^i}{2}}\right).`$ (3.53)
Again the insertion these operators in front of the boundary state corresponds to the tachyon condensation in the closed string sector.
This can be understood from the following Wess-Zumino coupling
$$_{D8\overline{D8}}C_{RR}\mathrm{STr}e^{},=\left(\begin{array}{cc}F^+T\overline{T}& DT\\ D\overline{T}& F^{}\overline{T}T\end{array}\right)$$
(3.54)
where $`D=d+A^+A^{}`$ and $`T,\overline{T}`$ represent the complex tachyon field. The fields $`A^+,A^{}`$ denote the gauge fields on the brane, anti-brane respectively, which is 0 in this case. Combining with the fact that the tachyon configuration is given by
$$\mathrm{for}\mathrm{example}T(x)\mathrm{\Gamma }_ix_i,\mathrm{at}x_i0,$$
(3.55)
we can speculate that the RR fields $`C_{2k+1}`$ should have the following Chan-Paton factors
$$C_7:\mathrm{\Gamma }_{ij},C_5:\mathrm{\Gamma }_{ijk\mathrm{}},C_3:\mathrm{\Gamma }_{ijk\mathrm{}mn},C_1:\mathrm{\Gamma }_{12345678}.$$
(3.56)
For example, this can be understood as the following expression
$$_{D8\overline{D8}}C_{D0}𝑑x^1dx^2\mathrm{}dx^8\mathrm{\Gamma }_{12345678}.$$
(3.57)
Thus the traces (3.53) correspond to the closed string sector that belongs to the each Chan-Paton factor. Thus the generation of the lower dimensional $`D`$-branes’ charges is induced by the above Wess-Zumino terms which are characteristic of the brane-antibrane systems.
Switching on the above operators, we obtain the following boundary states
$`|B(\{\alpha _i\}),+_{NSNS}^{(0)}`$ $`=`$ $`16{\displaystyle \underset{\stackrel{}{w}_\varphi 𝐙^8}{}}\left[K\left(\{m_i\}=\{w_\varphi ^i\}\right)\right]|\stackrel{}{0},2\stackrel{}{w}_\varphi _{NSNS}^{(0)},`$ (3.58)
$`|B(\{\alpha _i\}),+_{RR}^{(0)}`$ $`=`$ $`16{\displaystyle \underset{\stackrel{}{w}_\varphi 𝐙^8}{}}\left[K\left(\{m_i\}=\{w_\varphi ^i+{\displaystyle \frac{1}{2}}\}\right)\right]|\stackrel{}{0},2\stackrel{}{w}_\varphi +\stackrel{}{1},+_{RR}^{(0)},`$ (3.59)
where
$`K(\{m_i\})`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{8}{}}}\mathrm{cos}\left(\pi \alpha _im_i\right)+\left[{\displaystyle \underset{i=1}{\overset{6}{}}}\mathrm{cos}\left(\pi \alpha _im_i\right)\mathrm{sin}\left(\pi \alpha _7m_7\right)\mathrm{sin}\left(\pi \alpha _8m_8\right)+27\mathrm{terms}\right]`$
$`+\left[{\displaystyle \underset{i=1}{\overset{4}{}}}\mathrm{cos}\left(\pi \alpha _im_i\right){\displaystyle \underset{j=5}{\overset{8}{}}}\mathrm{sin}\left(\pi \alpha _jm_j\right)+69\mathrm{terms}\right]`$
$`+\left[\mathrm{cos}\left(\pi \alpha _1m_1\right)\mathrm{cos}\left(\pi \alpha _2m_2\right){\displaystyle \underset{i=3}{\overset{8}{}}}\mathrm{sin}\left(\pi \alpha _im_i\right)+27\mathrm{terms}\right]+{\displaystyle \underset{i=1}{\overset{8}{}}}\mathrm{sin}\left(\pi \alpha _im_i\right).`$
We have $`|B\left(\{0\}\right),+_{NSNS,RR}=|D8\overline{D8},+_{NSNS,RR}`$.
At $`\alpha _i=1`$ for all the $`i`$, the zero mode parts of the above boundary states become
$`\mathrm{NSNS}\mathrm{sector}`$ $`{\displaystyle \underset{\stackrel{}{w}_\varphi 𝐙^8}{}}(1)^{w_\varphi ^1+\mathrm{}+w_\varphi ^8}|\stackrel{}{0},2\stackrel{}{w}_\varphi _{NSNS}^{(0)},`$ (3.61)
$`\mathrm{RR}\mathrm{sector}`$ $`{\displaystyle \underset{\stackrel{}{w}_\varphi 𝐙^8}{}}(1)^{w_\varphi ^1+\mathrm{}+w_\varphi ^8}|\stackrel{}{0},2\stackrel{}{w}_\varphi +\stackrel{}{1},+_{RR}^{(0)}.`$ (3.62)
The boundary state of this system is rewritten as follows
$`|D0\overline{D0},+_{NSNS}^{(0)}`$ $`=`$ $`16{\displaystyle \underset{\stackrel{}{n}_Y𝐙^8}{}}|\stackrel{}{n}_Y,\stackrel{}{0}_{NSNS}^{(0)},`$ (3.63)
$`|D0\overline{D0},+_{RR}^{(0)}`$ $`=`$ $`16{\displaystyle \underset{\stackrel{}{n}_Y𝐙^8}{}}|\stackrel{}{n}_Y+\stackrel{}{1/2},\stackrel{}{0},+_{RR}^{(0)}.`$ (3.64)
These boundary states satisfy the definition equation (the boundary conditions) for the $`D0\overline{D0}`$ system.
Again, the extra phase factor $`(1)^{w_\varphi ^1+\mathrm{}+w_\varphi ^8}`$ changes the boundary condition from $`D8\overline{D8}`$ to $`D0\overline{D0}`$. This corresponds to the fact that $`D0`$-branes and $`\overline{D0}`$-branes are produced at each choices of
$$(x_{2k+1},x_{2k+2})=(0,0)\mathrm{or}(\pi ,\pi ).$$
(3.65)
Around the above points, there exists a soliton (anti-soliton) if the number of coordinate pairs taking the value $`(\pi ,\pi )`$ is even (odd). Finally evaluating the vacuum amplitudes for NSNS and RR sector with all the ghosts taking into account, it is easy to check the Cardy’s constraint explicitly. Thus we have established in the closed string viewpoint that a tachyonic soliton on the $`D8\overline{D8}`$ system produces a codimension eight $`D0`$-branes.
Next turn to the other points of $`\alpha _i`$. The tadpole cancellation restricts the admissible values to $`\alpha _i=0,1`$. We note the following basic observation.
$$\begin{array}{ccc}& \alpha =0& \alpha =1\\ \mathrm{cos}\left(\pi \alpha w\right)& 1& (1)^w\\ \mathrm{sin}\left(\pi \alpha w\right)& 0& 0\\ \mathrm{cos}\left(\pi \alpha \left(w+\frac{1}{2}\right)\right)& 1& 0\\ \mathrm{sin}\left(\pi \alpha \left(w+\frac{1}{2}\right)\right)& 0& (1)^w\end{array}$$
(3.66)
Then in the case when $`2n`$ of $`\alpha _i`$ is equal to 1, the system corresponds to the $`D(82n)\overline{D(82n)}`$ system. The number of such configurations corresponds to the possible choice of the Chan-Paton factor in the closed string sector. For example, 28 $`\mathrm{\Gamma }_{ij}`$ corresponds to the degree of freedom in order to set the direction of the codimension 2 among 8 directions. On the other hand, when the odd number of $`\alpha _i`$ is equal to 1, the boundary states in RR sector vanish and the system corresponds to a non-BPS $`D`$-branes.
Again we emphasize that the Chan-Paton factors in the closed string sector played the crucial role in our analysis.
## 4 Boundary state description of tachyon condensation in $`T^4/𝐙_2`$ orbifold theory
In this section we construct the boundary state description of tachyon condensations in the orbifold theory. We first discuss the decay mode from a $`D0\overline{D0}`$ system to a non-BPS D1-brane in detail. Next we extend this result to the higher codimension cases. We also discuss the occurrence of the bose-fermi degeneracy in this formalism.
### 4.1 Construction of the boundary state
The boundary state which represents $`D0\overline{D0}`$ in $`T^4/𝐙_2`$ orbifold at the radii $`R_6=\sqrt{2},R_7=R_8=R_9=R`$ is given as follows
$`|B={\displaystyle \frac{T_{p=0}}{2\sqrt{2}}}(|U_{NSNS}+|U_{RR})+{\displaystyle \frac{N}{2\sqrt{2}}}(|T_{NSNS}+|T_{RR}),`$ (4.1)
where $`|U,|T`$ denote the untwisted, twisted part of the boundary state. The normalization for twisted sector is determined by comparing the closed string vacuum amplitude with the open string one. The result is given by $`N=2^3\pi ^{\frac{3}{2}}`$.
The more detailed structure of each sector is written as
$`|U_{NSNS}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \left(\frac{dk}{2\pi }\right)^5\underset{n_6,n_7,n_8,n_9}{}\frac{1+(1)^{n_6}}{(2\pi \sqrt{2})(2\pi R)^3}\left[|U,+,k^i,n_{NSNS}|U,,k^i,n_{NSNS}\right]},`$
$`|U_{RR}`$ $`=`$ $`2{\displaystyle \left(\frac{dk}{2\pi }\right)^5\underset{n_6,n_7,n_8,n_9}{}\frac{1(1)^{n_6}}{(2\pi \sqrt{2})(2\pi R)^3}\left[|U,+,k^i,n_{RR}+|U,,k^i,n_{RR}\right]},`$
$`|T_{NSNS}`$ $`=`$ $`{\displaystyle \left(\frac{dk}{2\pi }\right)^5\left[|T_1,+,k^i_{NSNS}+|T_1,,k^i_{NSNS}+|T_2,+,k^i_{NSNS}+|T_2,,k^i_{NSNS}\right]},`$
$`|T_{RR}`$ $`=`$ $`{\displaystyle \left(\frac{dk}{2\pi }\right)^5\left[|T_1,+,k^i_{RR}+|T_1,,k^i_{RR}+|T_2,+,k^i_{RR}+|T_2,,k^i_{RR}\right]},`$ (4.2)
where we set $`1i5`$ and $`|T_1,|T_2`$ represent the twisted sector boundary states corresponding to two different fixed points. As we will explain briefly in the appendix A, $`|U,\gamma ,k^i,n`$ and $`|T_1,+,k^i,|T_2,+,k^i`$ are defined by the conditions (A.2) expanding the fields $`(Y,\psi )`$ in each Hilbert space.
Next we need to rewrite the above boundary state in terms of $`(\varphi ,\eta )`$ in order to describe the tachyon condensation as discussed in the previous section. For the untwisted sector the procedure is almost the same and the result are as follows (we show below only the relevant modes which correspond to $`x^6`$ direction and omit the superscript 6 in this subsection)
$`|U,+_{NSNS}`$ $`=`$ $`\mathrm{exp}({\displaystyle \underset{n𝐙}{}}{\displaystyle \frac{1}{n}}\alpha _n\stackrel{~}{\alpha }_ni{\displaystyle \underset{r𝐙+\frac{1}{2}}{}}\psi _r\stackrel{~}{\psi }_r)2{\displaystyle \underset{n_Y}{}}|2n_Y,0_{NSNS}^{(0)}`$
$`=`$ $`\mathrm{exp}({\displaystyle \underset{n𝐙}{}}{\displaystyle \frac{1}{n}}\varphi _n\stackrel{~}{\varphi }_n+i{\displaystyle \underset{r𝐙+\frac{1}{2}}{}}\eta _r\stackrel{~}{\eta }_r)2{\displaystyle \underset{w_\varphi }{}}(1)^{w_\varphi }|0,w_\varphi _{NSNS}^{(0)},`$
$`|U,+_{RR}`$ $`=`$ $`\mathrm{exp}({\displaystyle \underset{n𝐙}{}}{\displaystyle \frac{1}{n}}\alpha _n\stackrel{~}{\alpha }_ni{\displaystyle \underset{r𝐙}{}}\psi _r\stackrel{~}{\psi }_r)2{\displaystyle \underset{n_Y}{}}|2n_Y+1,0,+_{RR}^{(0)}`$
$`=`$ $`\mathrm{exp}({\displaystyle \underset{n𝐙}{}}{\displaystyle \frac{1}{n}}\varphi _n\stackrel{~}{\varphi }_n+i{\displaystyle \underset{r𝐙}{}}\eta _r\stackrel{~}{\eta }_r)2{\displaystyle \underset{w_\varphi }{}}(1)^{w_\varphi }|0,w_\varphi +{\displaystyle \frac{1}{2}},+_{RR}^{(0)},`$
Next let us turn to the twisted sector. The twist operator $`\sigma `$ which map the untwisted sector into twisted sector is needed . A candidate for such an operator is given as
$$\sigma =e^{i\frac{1}{\sqrt{2}}(\varphi _R\stackrel{~}{\varphi }_L)},$$
(4.4)
which has the desired singular property as
$$\psi (z)\sigma (0)O(z^{\frac{1}{2}}),X(z)\sigma (0)O(z^{\frac{1}{2}}),$$
(4.5)
This operator leads to the correct boundary condition of twisted sector boundary state. Then we can rewrite the twisted sector boundary state as
$`|T,+_{NSNS}`$ $`=`$ $`\mathrm{exp}({\displaystyle \underset{n𝐙}{}}{\displaystyle \frac{1}{n+\frac{1}{2}}}\alpha _{(n+\frac{1}{2})}\stackrel{~}{\alpha }_{(n+\frac{1}{2})}i{\displaystyle \underset{r𝐙}{}}\psi _r\stackrel{~}{\psi }_r)\times \{|T_1_{NSNS}^{(0)}+|T_2_{NSNS}^{(0)}\}`$
$`=`$ $`\mathrm{exp}({\displaystyle \underset{n𝐙}{}}{\displaystyle \frac{1}{n}}\varphi _n\stackrel{~}{\varphi }_n+i{\displaystyle \underset{r𝐙+\frac{1}{2}}{}}\eta _r\stackrel{~}{\eta }_r){\displaystyle \underset{w_\varphi }{}}(1)^{w_\varphi }|w_\varphi +{\displaystyle \frac{1}{2}}_{NSNS}^{(0)},`$
$`|T,+_{RR}`$ $`=`$ $`\mathrm{exp}({\displaystyle \underset{n𝐙}{}}{\displaystyle \frac{1}{n+\frac{1}{2}}}\alpha _{(n+\frac{1}{2})}\stackrel{~}{\alpha }_{(n+\frac{1}{2})}i{\displaystyle \underset{r𝐙+\frac{1}{2}}{}}\psi _r\stackrel{~}{\psi }_r)\times \{|T_1_{RR}^{(0)}+|T_2_{RR}^{(0)}\}`$
$`=`$ $`\mathrm{exp}({\displaystyle \underset{n𝐙}{}}{\displaystyle \frac{1}{n}}\varphi _n\stackrel{~}{\varphi }_n+i{\displaystyle \underset{r𝐙}{}}\eta _r\stackrel{~}{\eta }_r){\displaystyle \underset{w_\varphi }{}}(1)^{w_\varphi }|w_\varphi _{RR}^{(0)}.`$
Notice that this transformation or “bosonization” procedure can be verified by showing the bosonized boundary state does indeed satisfy the boundary condition of the original one as in section 3. It is also easy to see that the vacuum amplitude doesn’t change by the bosonization using the relations (3.37).
### 4.2 Tachyon condensation
Since we have constructed the boundary state of $`D0\overline{D0}`$ system in terms of $`(\varphi ,\eta )`$, it is straightforward to determine the boundary state which describe the tachyon condensation process in that system. The Wilson line corresponding to the tachyon condensation discussed in section 2 can be written as
$$W=\mathrm{Tr}\mathrm{exp}(\frac{i}{2\sqrt{2}}\alpha \varphi \sigma _1)=\mathrm{cos}(\pi w_\varphi \alpha ).$$
(4.7)
Then the effect of the tachyon condensation appears at the coefficients in front of the zero mode parts as in section 3. If we consider the point $`\alpha =1`$, which corresponds to the maximal condensation, then the untwisted RR-sector and the twisted NSNS-sector vanish. The untwisted NSNS-sector and the twisted RR-sector become as follows
$`|U,+_{NSNS}`$ $`=`$ $`\mathrm{exp}({\displaystyle \underset{n𝐙}{}}{\displaystyle \frac{1}{n}}\varphi _n\stackrel{~}{\varphi }_n+i{\displaystyle \underset{r𝐙+\frac{1}{2}}{}}\eta _r\stackrel{~}{\eta }_r)2{\displaystyle \underset{w_\varphi }{}}|w_\varphi _{NSNS}^{(0)}`$
$`=`$ $`\mathrm{exp}({\displaystyle \underset{n𝐙}{}}{\displaystyle \frac{1}{n}}\alpha _n\stackrel{~}{\alpha }_n+i{\displaystyle \underset{r𝐙+\frac{1}{2}}{}}\psi _r\stackrel{~}{\psi }_r)2{\displaystyle \underset{w_Y}{}}|w_Y_{NSNS}^{(0)},`$
$`|T,+_{RR}`$ $`=`$ $`\mathrm{exp}({\displaystyle \underset{n𝐙}{}}{\displaystyle \frac{1}{n}}\varphi _n\stackrel{~}{\varphi }_n+i{\displaystyle \underset{r𝐙}{}}\eta _r\stackrel{~}{\eta }_r){\displaystyle \underset{w_\varphi }{}}|w_\varphi _{RR}^{(0)}`$
$`=`$ $`\mathrm{exp}({\displaystyle \underset{n𝐙}{}}{\displaystyle \frac{1}{n+\frac{1}{2}}}\alpha _{(n+\frac{1}{2})}\stackrel{~}{\alpha }_{(n+\frac{1}{2})}+i{\displaystyle \underset{r𝐙+\frac{1}{2}}{}}\psi _r\stackrel{~}{\psi }_r)\times \{|T_1_{RR}^{(0)}+|T_2_{RR}^{(0)}\},`$
where we have “rebosonized” the expression using the basis $`(Y,\psi )`$. Note that the boundary condition is changed into that of D1-brane because of the extra phase $`\mathrm{cos}(\pi w_\varphi )=(1)^{w_\varphi }`$.
Now it is obvious<sup>20</sup><sup>20</sup>20If we start a $`D0\overline{D0}`$ which has the different relative twisted charge, then we can show by using almost the same procedure that the final object is a non-BPS D1-brane with a $`𝐙_2`$ Wilson line after the tachyon condensation. that the above boundary state is the same as that of a non-BPS D1-brane stretching between the fixed points.
In this way the tachyon condensation process from $`D0\overline{D0}`$ to a non-BPS D1-brane (and also its reverse if we replace $`\alpha `$ with $`1\alpha `$) is explicitly shown by using boundary state formalism. It would be an interesting fact that the twisted sector of $`(Y,\psi )`$ can be expressed by using the untwisted sector of another field basis$`(\varphi ,\eta )`$ and this is crucial in the above discussion of the tachyon condensation in the orbifold theory. This fact will also become important if we consider tachyon condensation processes in other orbifold theories.
### 4.3 Generalization to the higher codimension case
Then we will be interested in the higher codimension cases. As we will show below, such generalizations are not so difficult in our boundary state formalism and the results remain almost the same as in section 3. Therefore the discussion is short.
To make things clear we consider the tachyon condensation that changes two $`D4\overline{D4}`$ pairs into two $`D0\overline{D0}`$ pairs (codimension four). Here $`D4\overline{D4}`$ system has appropriate $`𝐙_2`$ Wilson lines in the same sense of section 3. This process includes the decay modes into two $`D2\overline{D2}`$ pairs. First let us define the coordinates of $`T^4`$ as $`(x^6,x^7,x^8,x^9)`$ and their $`\frac{\pi }{4}`$ rotated coordinates as $`(y^6,y^7,y^8,y^9)`$. We also take the radii of $`T^4`$ as $`R^6=R^7=R^8=R^9=1`$ in terms of the coordinates $`(x^6,x^7,x^8,x^9)`$ as in the previous discussion in flat space. It is important to note that this $`D4\overline{D4}`$ system can be described in terms of $`(y^6,y^7,y^8,y^9)`$ as a $`D4\overline{D4}`$ system on $`T^4/𝐙_2`$ of which radii are all $`\sqrt{2}`$ with the projection $`(1)^Fh_6h_7h_8h_9=1`$ as in the case of flat space. At this radius we can change the basis $`(Y^i,\psi ^i)`$ into $`(\varphi ^i,\eta ^i)`$ by the bosonization procedures, which are trivial generalizations of eq.(4.1) and (4.1). Then we can describe the tachyon condensation processes and let us denote the corresponding four parameters as $`\alpha _1,\alpha _2,\alpha _3,\alpha _4`$. Notice that in order to get four parameters<sup>21</sup><sup>21</sup>21If we started with one $`D4\overline{D4}`$, then we would only get the decay modes into a $`D2\overline{D2}`$ and the codimension four configuration is impossible. corresponding to the marginal deformation in the four directions we should start with not one but two pieces of $`D4\overline{D4}`$. At the point $`\alpha _1=\alpha _2=\alpha _3=\alpha _4=1`$ both the untwisted and the twisted boundary states gain the same extra phase factor $`(1)^{w_6+w_7+w_8+w_9}`$ and the boundary conditions along $`(y^6,y^7,y^8,y^9)`$ are reversed. Then we get two $`D0`$-branes and two anti $`D0`$-branes which sit at $`(x^6,x^7,x^8,x^9)=(0,0,0,0),(\pi ,\pi ,\pi ,\pi )`$ and $`(\pi ,\pi ,0,0),(0,0,\pi ,\pi )`$ respectively. In this way we find that the tachyon condensation of brane-antibrane system in $`T^4/𝐙_2`$ orbifold can be treated almost in the same way as in flat space except the treatment of the twisted sector.
### 4.4 Comments on bose-fermi degeneracy
Finally let us discuss the relation between the boundary state description in this section and the bose-fermi degeneracy . First we compute the vacuum amplitude of the system (4.1),(4.1) with the insertion of the Wilson line (4.7) using $`(\varphi ,\eta )`$ field representation for $`x^6`$ direction. The result is
$`B(\alpha )|\mathrm{\Delta }|B(\alpha )`$
$`={\displaystyle \frac{V}{16}}{\displaystyle }ds({\displaystyle \frac{1}{2\pi s}})^{\frac{5}{2}}[{\displaystyle \frac{\pi ^32^{\frac{5}{2}}}{R^3}}\{{\displaystyle \underset{w_\varphi ,\stackrel{}{n}_Y}{}}q^{w_{\varphi }^{}{}_{}{}^{2}+\frac{n_Y^2}{2R^2}}\mathrm{cos}^2(\pi \alpha w_\varphi ){\displaystyle \frac{f_3(q)^8}{f_1(q)^8}}{\displaystyle \underset{\stackrel{}{n}_Y}{}}q^{\frac{n_Y^2}{2R^2}}{\displaystyle \frac{\sqrt{2}f_3(q)f_4(q)^7}{f_2(q)f_1(q)^7}}`$
$`{\displaystyle \underset{w_\varphi ,\stackrel{}{n}_Y}{}}\mathrm{cos}^2(\pi \alpha (w_\varphi +{\displaystyle \frac{1}{2}}))q^{(w_\varphi +\frac{1}{2})^2+\frac{n_Y^2}{2R^2}}{\displaystyle \frac{f_2(q)^8}{f_1(q)^8}}\}`$ (4.9)
$`+2^6\pi ^3\{{\displaystyle \underset{w_\varphi }{}}\mathrm{cos}^2(\pi (w_\varphi +{\displaystyle \frac{1}{2}})\alpha )q^{(w_\varphi +\frac{1}{2})^2}{\displaystyle \frac{f_3(q)^5f_2(q)^3}{\sqrt{2}f_1(q)^5f_4(q)^3}}{\displaystyle \underset{w_\varphi }{}}\mathrm{cos}^2(\pi w_\varphi \alpha )q^{w_\varphi ^2}{\displaystyle \frac{f_2(q)^5f_3(q)^3}{\sqrt{2}f_1(q)^5f_4(q)^3}}\}],`$
where $`\stackrel{}{n}_Y=(n_Y^7,n_Y^8,n_Y^9)𝐙^3`$ are momenta in the directions of $`(y^7,y^8,y^9)`$. We can see that this amplitude does vanish if $`R=\frac{1}{\sqrt{2}},\alpha =1`$ and this phenomenon of non-BPS D1-brane is called bose-fermi degeneracy . Below we would like to discuss this from the viewpoint of the boundary state.
The particular radii of torus $`R=\frac{1}{\sqrt{2}}`$ enable us to perform further bosonization procedures in the direction of $`x^7,x^8,x^9`$. The result is as follows (we only show the zero modes and oscillators which correspond to $`x^6,x^7,x^8,x^9`$)
$`|U,+_{NSNS}`$ $`=`$ $`\mathrm{exp}[{\displaystyle \underset{n𝐙}{}}{\displaystyle \frac{1}{n}}(\alpha _n^6\stackrel{~}{\alpha }_n^6+{\displaystyle \underset{i=7,8,9}{}}\alpha _n^i\stackrel{~}{\alpha }_n^i)`$
$`i{\displaystyle \underset{r𝐙+\frac{1}{2}}{}}(\psi _r^6\stackrel{~}{\psi }_r^6+{\displaystyle \underset{i=7,8,9}{}}\psi _r^i\stackrel{~}{\psi }_r^i)]2{\displaystyle \underset{w_Y^6,\stackrel{}{n}_Y}{}}|w_Y^6,n_Y^7,n_Y^8,n_Y^9_{NSNS}^{(0)},`$
$`=`$ $`\mathrm{exp}[{\displaystyle \underset{n𝐙}{}}{\displaystyle \frac{1}{n}}(\varphi _n^6\stackrel{~}{\varphi }_n^6{\displaystyle \underset{i=7,8,9}{}}\varphi _n^i\stackrel{~}{\varphi }_n^i)`$
$`i{\displaystyle \underset{r𝐙+\frac{1}{2}}{}}(\eta _r^6\stackrel{~}{\eta }_r^6{\displaystyle \underset{i=7,8,9}{}}\eta _r^i\stackrel{~}{\eta }_r^i)]2{\displaystyle \underset{w_\varphi }{}}(1)^{w_\varphi ^7+w_\varphi ^8+w_\varphi ^9}|w_\varphi ^6,2w_\varphi ^7,2w_\varphi ^8,2w_\varphi ^9_{NSNS}^{(0)},`$
$`|T,+_{RR}`$ $`=`$ $`\mathrm{exp}[{\displaystyle \underset{n𝐙}{}}{\displaystyle \frac{1}{n+\frac{1}{2}}}(\alpha _{n\frac{1}{2}}^6\stackrel{~}{\alpha }_{n\frac{1}{2}}^6+{\displaystyle \underset{i=7,8,9}{}}\alpha _{n\frac{1}{2}}^i\stackrel{~}{\alpha }_{n\frac{1}{2}}^i)`$
$`i{\displaystyle \underset{r𝐙+\frac{1}{2}}{}}(\psi _r^6\stackrel{~}{\psi }_r^6+{\displaystyle \underset{i=7,8,9}{}}\psi _r^i\stackrel{~}{\psi }_r^i)]\{|T_1_{RR}^{(0)}+|T_2_{RR}^{(0)}\},`$
$`=`$ $`\mathrm{exp}[{\displaystyle \underset{n𝐙}{}}{\displaystyle \frac{1}{n}}(\varphi _n^6\stackrel{~}{\varphi }_n^6{\displaystyle \underset{i=7,8,9}{}}\varphi _n^i\stackrel{~}{\varphi }_n^i)`$
$`i{\displaystyle \underset{r𝐙}{}}(\eta _r^6\stackrel{~}{\eta }_r^6{\displaystyle \underset{i=7,8,9}{}}\eta _r^i\stackrel{~}{\eta }_r^i)]2{\displaystyle \underset{w_\varphi }{}}(1)^{w_\varphi ^7+w_\varphi ^8+w_\varphi ^9}|w_\varphi ^6,2w_\varphi ^7,2w_\varphi ^8,2w_\varphi ^9,+_{RR}^{(0)}.`$
This expression shows that the twisted RR-sector in terms of $`(Y,\psi )`$ is rewritten to have the same form as the untwisted RR-sector of D4-brane in terms of $`(\varphi ,\eta )`$<sup>22</sup><sup>22</sup>22This expression also implies that the original non-BPS D1-brane in $`T^4/𝐙_2`$ can be thought as a BPS D4-brane in terms of the field $`(\varphi ,\eta )`$ with a “wrong GSO projection” $`(1)^F=I_{4}^{}{}_{}{}^{\varphi }(1)^{F^\varphi }=1`$ , though the essence of this interpretation is not clear.. If we use the basis $`(Y,\psi )`$ for NSNS-sector and $`(\varphi ,\eta )`$ for RR-sector, each open string vacuum amplitude of NS-sector and R-sector cancels each other and the occurrence of the bose-fermi degeneracy is explicitly shown. Therefore we can say that the bosonization procedures at critical radius are crucial in the bose-fermi degeneracy.
## 5 Conclusions
In this paper we have shown explicitly in the boundary state formalism that the tachyon condensation in $`2^{k1}`$ pieces of $`D(p+2k)\overline{D(p+2k)}`$ at critical radii produces $`2^{k1}`$ pieces of $`Dp\overline{Dp}`$. Locally this means that a codimension $`2k`$ soliton of the tachyon field configuration corresponds to a $`Dp`$-brane (or $`\overline{Dp}`$). We have also verified this results in $`T^4/𝐙_2`$ orbifold theory. Note that in these cases there are no gauge fields on the world volume. But the generations of lower D-brane charges indeed occur due to the Wess-Zumino terms which are peculiar to brane-antibrane systems . In the boundary state description we have succeeded to see these phenomena explicitly.
In the process of the explicit calculations we have found two remarkable facts. The first is that the consistency with the open string picture (or Cardy’s condition) requires the closed string sectors should have nontrivial Chan-Paton factors. This somewhat strange phenomenon only occurs if we discuss interactions of closed strings with brane-antibrane systems or non-BPS D-branes. These Chan-Paton factors also ensure the Wess-Zumino coupling proposed in .
The second one is the fact in the case of $`T^4/𝐙_2`$ orbifold we can treat the twisted sector boundary state in the same way as the untwisted one by changing the field basis (or by “bosonization” procedure). This enables us to construct the boundary state which describe the tachyon condensation in the orbifold theory. Another application of this fact is the investigation of bose-fermi degeneracy . At the point where the degeneracy occurs the boundary state of a non-BPS D1-brane becomes very much like that of a BPS D-brane by using the bosonization procedure. Naively it seems that a sort of a symmetry is enhanced at this particular moduli, but it is difficult to see this explicitly even in our formalism. We leave this as a future problem. So far the tachyon condensation in four dimensional orbifold theories other than $`T^4/𝐙_2`$ have not been discussed. If one try to construct the marginal deformations of BCFT in them, something like the previous bosonization procedures of the boundary state will be required.
Acknowledgments
T.T. would like to be grateful to Y. Matsuo for useful discussions and remarks. The work of M.N. and T.T. is supported by JSPS Research Fellowships for Young Scientists.
## Appendix A Structure of boundary state
Here we give our CFT conventions and a short review of boundary states for general Dp-branes with or without the orbifold projection. Remember that we have used the light cone formulation and ignored<sup>23</sup><sup>23</sup>23In the case of the orbifold theory discussed in section 4, we ignored the non-zero modes of $`(X^4,X^5,\psi ^5,\psi ^6)`$. the non-zero modes of the fields $`(Y^8,Y^9,\psi ^8,\psi ^9)`$ in the case of $`D2\overline{D2}`$ system. We use almost the same conventions as Sen’s except the detailed normalizations.
### A.1 CFT conventions
We define $`z=e^{i\sigma _1+\sigma _2}`$ as the cylindrical coordinate of the world sheet and $`w=\sigma _1+i\sigma _2`$ as its radial plane coordinate. First we list the mode expansions of $`(Y,\psi )`$ fields :
$`Y_R^i(z)=y_R^i{\displaystyle \frac{i}{2}}p_{YR}^i\mathrm{ln}z+{\displaystyle \frac{i}{\sqrt{2}}}{\displaystyle \underset{n0}{}}{\displaystyle \frac{1}{n}}{\displaystyle \frac{(\alpha _Y)_n^i}{z^n}},`$
$`Y_L^i(\overline{z})=y_L^i{\displaystyle \frac{i}{2}}p_{YL}^i\mathrm{ln}\overline{z}+{\displaystyle \frac{i}{\sqrt{2}}}{\displaystyle \underset{n0}{}}{\displaystyle \frac{1}{n}}{\displaystyle \frac{(\stackrel{~}{\alpha }_Y)_n^i}{\overline{z}^n}},`$
$`\psi _R^i(z)=i^{\frac{1}{2}}{\displaystyle \underset{r𝐙+\nu }{}}{\displaystyle \frac{\psi _r^i}{z^{r+\frac{1}{2}}}},\psi _L^i(\overline{z})=i^{\frac{1}{2}}{\displaystyle \underset{r𝐙+\nu }{}}{\displaystyle \frac{\stackrel{~}{\psi }_r^i}{\overline{z}^{r+\frac{1}{2}}}},`$ (A.1)
$`(\widehat{p}_{YR}^i=\sqrt{2}(\alpha _Y)_0^i,\widehat{p}_{YL}^i=\sqrt{2}(\stackrel{~}{\alpha }_Y)_0^i),`$
where $`\nu =\frac{1}{2}`$ represents NS-sector and $`\nu =0`$ R-sector and we set $`i=09`$.
Then the OPE relations (2.1) are equivalent to the following (anti)commutation relations for modes :
$`[y_L^i,p_{YL}^j]=[y_R^i,p_{YR}^j]=i\eta ^{ij},`$
$`[(\alpha _Y)_m^i,(\alpha _Y)_n^j]=[(\stackrel{~}{\alpha }_Y)_m^i,(\stackrel{~}{\alpha }_Y)_n^j]=m\delta _{m,n}\eta ^{ij},`$
$`\{\psi _r^i,\psi _s^j\}=\{\stackrel{~}{\psi }_r^i,\stackrel{~}{\psi }_s^j\}=\delta _{r,s}\eta ^{ij},`$ (A.2)
where $`\eta ^{ij}=\delta ^{ij}`$ for $`i=19`$ and $`\eta ^{ij}=\delta ^{ij}`$ for $`i=0`$. The vacuum of these modes is defined as $`|\mathrm{\Omega }_{NSNS}^{(0)}`$. If we compactify the coordinates $`y^i`$ on torus (radii $`R^i`$) then the momenta are quantized as follows
$`p_{YR}^i={\displaystyle \frac{n_Y^i}{R^i}}+R^iw_Y^i,`$ $`p_{YL}^i={\displaystyle \frac{n_Y^i}{R^i}}R^iw_Y^i,`$ (A.3)
where $`n_Y^i`$ and $`w_Y^i`$ denote K.K. modes and winding modes.
We have also used the fields $`(X,\chi )`$ and $`(\varphi ,\eta )`$ as another bases. The mode expansions and commutation relations of these fields are defined in the same way.
### A.2 Definition of boundary states
A boundary state of Dp-brane is defined by the following boundary conditions<sup>24</sup><sup>24</sup>24Of course the conditions remain the same if we replace $`(Y,\psi )`$ with $`(X,\chi )`$, because this procedure does not mix the Neumann and Dirichlet conditions. in the closed string Hilbert space :
$`_2Y^\mu (w)|_{\sigma _2=0}|Dp,\gamma `$ $`=`$ $`0,(\mu =0p),`$
$`_1Y^i(w)|_{\sigma _2=0}|Dp,\gamma `$ $`=`$ $`0,(i=p+17),`$
$`(\psi _R^\mu (w)i\gamma \psi _L^\mu (\overline{w}))|_{\sigma _2=0}|Dp,\gamma `$ $`=`$ $`0,(\mu =0p),`$
$`(\psi _R^i(w)+i\gamma \psi _L^i(\overline{w}))|_{\sigma _2=0}|Dp,\gamma `$ $`=`$ $`0,(i=p+17),`$ (A.4)
where $`\gamma =+,`$ is the spin structure on the boundary and the GSO projection of the closed string determines the correct linear combination of these spin structures. If we expand the left-hand side of eq. (A.2), we get
$`((\alpha _Y)_n^\mu +(\stackrel{~}{\alpha }_Y)_n^\mu )|Dp,\gamma `$ $`=`$ $`0,(\mu =0p),`$
$`((\alpha _Y)_n^i(\stackrel{~}{\alpha }_Y)_n^i)|Dp,\gamma `$ $`=`$ $`0,(i=p+17),`$
$`(\psi _r^\mu i\gamma \stackrel{~}{\psi }_r^\mu )|Dp,\gamma `$ $`=`$ $`0,(\mu =0p),`$
$`(\psi _r^i+i\gamma \stackrel{~}{\psi }_r^i)|Dp,\gamma `$ $`=`$ $`0,(i=p+17).`$ (A.5)
These conditions are easy to solve by using the commutation relations (A.1),(A.1). Notice that for a BPS D-brane the boundary state consists of the NSNS-sector and RR-sector and the correct linear combination of them should be determined by comparing its cylinder amplitude with that of open string (see ). For example in the case of a (BPS) D2-brane the boundary state is given as eq.(3.1). Also note that for a non-BPS D-brane there is no RR-sector.
Finally let us see the orbifold case briefly. In general the orbifold theories have twisted sectors in the closed string Hilbert space and therefore it is necessary to add twisted sector boundary states $`|T`$ to the untwisted one. The twisted sector boundary states are defined by the same equation (A.2), but the mode expansion is different from (A.1) because of the twisted boundary condition. In the case of $`T^4/𝐙_2`$ orbifold discussed in section 4, the mode expansion of $`(Y^i,\psi ^i)(i=6,7,8,9)`$ is shifted by half integer. For example, the twisted sector boundary state of $`D0\overline{D0}`$ is given as eq. (4.1), where we showed only the modes of $`(Y^6,\psi ^6)`$. The correct linear combination of the twisted sector boundary states and the untwisted one is also determined by the calculations of the cylinder amplitude and this is called the Cardy’s condition .
## Appendix B Equivalence of a boundary state and its bosonized version
In section 3, 4 we have used bosonized (and fermionized) descriptions of boundary states at special radii. In the authors calculate several one point functions in the codimension one case and show that the results are the same as those before the bosonization. As a further evidence of the equivalence here we prove that the bosonized boundary states discussed in section 3.1 satisfy the correct boundary conditions in the case of the tachyon condensation in $`D2\overline{D2}`$ system. The other cases appeared in this paper can be treated almost in the same way.
### B.1 Cocycle factors
In order to prove the correct boundary conditions the detailed cocyle factors should be given explicitly. For example, the fermionization relations (LABEL:eqn:bos) are written incorporating the cocyle factors as
$`\tau _iC_Y^i(\pm \sqrt{2},0)e^{\pm i\sqrt{2}Y_R^i}(z)`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\xi _R^i\pm i\eta _R^i)(z),`$
$`\tau _iC_Y^i(0,\pm \sqrt{2})e^{\pm i\sqrt{2}Y_L^i}(\overline{z})`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\xi _L^i\pm i\eta _L^i)(\overline{z}),`$ (B.1)
where $`\tau _i,C_Y^i(k_R,k_L)`$ are both called cocycle factors . $`\tau _i(i=1,2,3)`$ are $`2\times 2`$ Pauli matrices. $`C_Y^i(k_R,k_L)`$ are defined by (for example see )
$`C_Y^i(k_R,k_L)\mathrm{exp}\left[{\displaystyle \frac{1}{4}}\pi i(k_{YR}^ik_{YL}^i)(\widehat{p}_{YR}^i+\widehat{p}_{YL}^i)\right]`$ (B.2)
In bosonization procedure, they are needed to guarantee correct (anti)commutation relations between various fields.
Next step is the rebosonization of two fermions,
$`{\displaystyle \frac{1}{\sqrt{2}}}(\xi _R^i\pm i\psi _R^i)(z)`$ $``$ $`\stackrel{~}{\tau }_iC_\varphi ^i(\pm \sqrt{2},0)e^{\pm i\sqrt{2}\varphi _R^i}(z)`$
$`{\displaystyle \frac{1}{\sqrt{2}}}(\xi _L^i\pm i\psi _L^i)(\overline{z})`$ $``$ $`\stackrel{~}{\tau }_iC_\varphi ^i(0,\pm \sqrt{2})e^{\pm i\sqrt{2}\varphi _L^i}(\overline{z})`$ (B.3)
In this way we accomplished changing variables from $`(Y^i,\psi ^i)`$ to $`(\varphi ^i,\eta ^i)`$(see Figure 3).
### B.2 Proof of the correct boundary conditions
Now let us prove the facts that the bosonized boundary states satisfy the correct boundary conditions of the original ones and give a evidence that they are equivalent. We take the example of $`D2\overline{D2}`$ system discussed in section 3. Then two types of equivalence should be proved. The first is that eq.(3.14) and (3.15) are equivalent to eq.(3.12) and (3.13) respectively. We can verify this by showing that eq.(3.14) and (3.15) satisfy the boundary conditions eq.(3.16),(3.17). The second case is that $`|B(\alpha =1,\beta =1),+`$ (see eq.(3.25),(3.1)) is equivalent to the boundary state of $`D0\overline{D0}`$ system. We can also prove this in the same way by showing eq.(3.30),(3.31). Since these four equations can be proven in the same way, we show the proof of (3.16)below.
First let us note that eq.(3.14),(3.15) satisfy
$`(\eta _R^i(w)i\eta _L^i(\overline{w}))|_{\sigma _2=0}|D2\overline{D2},+_{NSNS,RR}=0,`$ (B.4)
and that we can replace $`_2Y^i(w)|_{\sigma _2=0}`$ with $`(\varphi ,\eta )`$ variables. Then eq.(3.16) can be rewritten as
$`_2Y^i(w)|_{\sigma _2=0}|D2\overline{D2},+_{NSNS,RR}`$
$`=\stackrel{~}{\tau }_i{\displaystyle \frac{\sqrt{z}}{2i}}\eta _R^i(z)[\sqrt{z}\{C_\varphi ^i(\sqrt{2},0)e^{i\sqrt{2}\varphi _R^i}(z)+C_\varphi ^i(\sqrt{2},0)e^{i\sqrt{2}\varphi _R^i}(z)\}`$
$`+\sqrt{\overline{z}}\{C_\varphi ^i(0,\sqrt{2})e^{i\sqrt{2}\varphi _L^i}(\overline{z})+C_\varphi ^i(0,\sqrt{2})e^{i\sqrt{2}\varphi _L^i}(\overline{z})\}]|D2\overline{D2},+_{NSNS,RR}.`$
The detail of the exponential is given as
$`:e^{i\sqrt{2}\varphi _R^i}(z):|_{\sigma _2=0}`$
$`=`$ $`\mathrm{exp}\left[i\sqrt{2}\varphi _R\right]\mathrm{exp}\left[{\displaystyle \frac{i}{\sqrt{2}}}p_R\sigma _1\right]\mathrm{exp}\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}\varphi _ne^{in\sigma _1}\right]\mathrm{exp}\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}\varphi _ne^{in\sigma _1}\right].`$
If we note that eq.(3.14),(3.15) satisfy
$`(\varphi _n+\stackrel{~}{\varphi }_n)|D2\overline{D2},+_{NSNS,RR}=0,`$ (B.7)
then it is easy to see that the first and the third, the second and the fourth term in eq. (B.2) cancel respectively. The proof is almost the same as in the case of $`|B(\alpha =1,\beta =1),+`$ except that the $`𝐙_2`$-phases $`(1)^{w_\varphi ^1+w_\varphi ^2}`$ of eq. (3.1) play an important role for changing boundary conditions of $`Y^1,Y^2`$.
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# References
IC/00/65
UCRHEP-T277
May 2000
Neutrino Masses and the Gluino Axion Model
D. A. Demir<sup>1</sup>, Ernest Ma<sup>2</sup>, and Utpal Sarkar<sup>2,3</sup>
<sup>1</sup> The Abdus Salam International Centre for Theoretical Physics, Trieste 34100, Italy
<sup>2</sup> Department of Physics, University of California, Riverside, California 92521, USA
<sup>3</sup> Physical Research Laboratory, Ahmedabad 380 009, India
## Abstract
We extend the recently proposed gluino axion model to include neutrino masses. We discuss how the canonical seesaw model and the Higgs triplet model may be realized in this framework. In the former case, the heavy singlet neutrinos are contained in superfields which do not have any vacuum expectation value, whereas the gluino axion is contained in one which does. We also construct a specific renormalizable model which realizes the mass scale relationship $`M_{SUSY}f_a^2/M_U`$, where $`f_a`$ is the axion decay constant and $`M_U`$ is a large effective mass parameter.
A new axionic solution to the strong CP problem was recently proposed. Instead of coupling to ordinary matter as in the DFSZ model or to unknown matter as in the KSVZ model, this new axion couples to the gluino as well as all other supersymmetric particles. The instanton-induced CP violating phase of quantum chromodynamics is then canceled by the dynamical phase of the gluino mass, as opposed to that of the quarks in the DFSZ model and that of the unknown colored fermions in the KSVZ model. This means that CP violation is absent in the strong-interaction sector and experimental observables, such as the neutron electric dipole moment, are subject only to weak-interaction contributions.
What sets the gluino axion model apart from all other previous models is its identification of the Peccei-Quinn global symmetry $`U(1)_{PQ}`$ with the $`U(1)_R`$ symmetry of superfield transformations. Under $`U(1)_R`$, the scalar components of a chiral superfield transform as $`\varphi e^{i\theta R}\varphi `$, whereas the fermionic components transform as $`\psi e^{i\theta (R1)}\psi `$. In the Minimal Supersymmetric Standard Model (MSSM), the quark and lepton superfields $`\widehat{Q}`$, $`\widehat{u}^c`$, $`\widehat{d}^c`$, $`\widehat{L}`$, $`\widehat{e}^c`$ have $`R=+1`$ whereas the Higgs superfields $`\widehat{H}_u`$, $`\widehat{H}_d`$ have $`R=0`$. The superpotential
$$\widehat{W}=\mu \widehat{H}_u\widehat{H}_d+h_u\widehat{H}_u\widehat{Q}\widehat{u}^c+h_d\widehat{H}_d\widehat{Q}\widehat{d}^c+h_e\widehat{H}_d\widehat{L}\widehat{e}^c$$
(1)
has $`R=+2`$ except for the $`\mu `$ term (which has $`R=0`$). Hence the resulting Lagrangian breaks $`U(1)_R`$ explicitly, leaving only a discrete remnant, i.e. the usual $`R`$ parity: $`R=(1)^{3B+L+2J}`$. The gluino axion model replaces $`\mu `$ with a singlet composite superfield of $`R=+2`$ so that the resulting supersymmetric Lagrangian is invariant under $`U(1)_R`$. It also requires all supersymmetry breaking terms to be invariant under $`U(1)_R`$, the spontaneous breaking of which then produces the axion and solves the strong CP problem.
In the MSSM, neutrinos are massless. However, in view of the recent experimental evidence for neutrino oscillations, it is desirable to incorporate into any realistic model naturally small Majorana neutrino masses. In the following we will discuss how the canonical seesaw model and the Higgs triplet model may be realized in the framework of the gluino axion model. In the case of the seesaw model, there are in fact proposals that the axion scale is the same as that of the singlet neutrino masses.
Consider first the Higgs triplet model. Add to the gluino axion model two triplet superfields:
$`\widehat{\xi }_1`$ $`=`$ $`(\xi _1^{++},\xi _1^+,\xi _1^0):R=0,`$ (2)
$`\widehat{\xi }_2`$ $`=`$ $`(\xi _2^0,\xi _2^{},\xi _2^{}):R=+2,`$ (3)
then the superpotential (which is required to have $`R=+2`$) has the following additional terms:
$$\mathrm{\Delta }\widehat{W}=m_\xi \widehat{\xi }_1\widehat{\xi }_2+f_{ij}\widehat{\xi }_1\widehat{L}_i\widehat{L}_j+h\widehat{\xi }_2\widehat{H}_u\widehat{H}_u.$$
(4)
Note that the term $`\widehat{\xi }_1\widehat{H}_d\widehat{H}_d`$ is forbidden. The resulting scalar potential has the term $`|m_\xi \xi _1+hH_uH_u|^2`$, hence the desired trilinear scalar interaction $`hm_\xi \xi _1^{}H_uH_u+h.c.`$ is there to combine with the Yukawa interaction $`f_{ij}\xi _1L_iL_j+h.c.`$ to form the well-known dimension-5 effective operator which generates the neutrino masses:
$$(m_\nu )_{ij}=2f_{ij}h\frac{H_u^2}{m_\xi }.$$
(5)
If the intermediate scale $`m_\xi `$ is assumed to be of order the $`U(1)_R`$ breaking scale, i.e. $`10^{11}`$ GeV or so, then $`m_\nu `$ of order 1 eV is obtained if $`f_{ij}h`$ is of order $`10^2`$.
Consider next the canonical sesaw model. Add to the gluino axion model the singlet superfield $`\widehat{N}`$ with $`R=+1`$, then the superpotential is supplemented by
$$\mathrm{\Delta }\widehat{W}=m_N\widehat{N}\widehat{N}+f_i\widehat{L}_i\widehat{N}\widehat{H}_u,$$
(6)
which generates the well-known seesaw neutrino mass
$$(m_\nu )_{ij}=f_if_j\frac{H_u^2}{m_N}.$$
(7)
Since both $`\widehat{N}`$ and $`\widehat{S}`$ have the same $`U(1)_R`$ charge, it is tempting to identify them as one, so that its scalar component has a large vacuum expectation value (VEV) and contains the axion, while its fermionic component is the heavy neutrino singlet of mass $`m_N`$. However, the resulting scalar potential will now contain the term $`|2m_N\stackrel{~}{N}+f_i\stackrel{~}{L}_iH_u|^2`$, so that the scalar bilinear term $`\stackrel{~}{L}_iH_u`$ (which violates lepton number) has the huge coefficient $`2f_im_N\stackrel{~}{N}`$ which is clearly unacceptable. To prevent $`\widehat{N}`$ from picking up any VEV, we introduce the discrete symmetry $`L`$ parity, under which $`\widehat{L}`$, $`\widehat{e}^c`$, and $`\widehat{N}`$ are odd and all other superfields are even, including $`\widehat{S}`$.
In proposing the gluino axion model, the composite operator $`\mu (\widehat{S})(\widehat{S})^2/M_{Pl}`$ with $`R=+2`$ is used. The couplings of $`\mu (S)`$ to the supersymmetric particles of the MSSM are required to be invariant under $`U(1)_R`$. Hence the supersymmetry of the MSSM is broken by $`\mu _{eff}=S^2/M_{Pl}`$. In the following we consider an alternative scheme, using the fundamental singlet superfields $`\widehat{S}_2`$, $`\widehat{S}_1`$, and $`\widehat{S}_0`$, with $`R=2,1,0`$ respectively. We impose the discrete symmetry $`Z_3`$ with $`\omega ^3=1`$ on all superfields as follows:
$`1:`$ $`\widehat{u}^c,\widehat{d}^c,\widehat{e}^c,\widehat{N},`$ (8)
$`\omega :`$ $`\widehat{Q},\widehat{L},\widehat{S}_1,\widehat{S}_0,\widehat{\xi }_1`$ (9)
$`\omega ^2:`$ $`\widehat{H}_u,\widehat{H}_d,\widehat{S}_2,\widehat{\xi }_2.`$ (10)
We see then that Eqs. (4) and (6) are allowed in addition to Eq. (1) except for the $`\mu `$ term. The superpotential involving $`\widehat{S}_2`$, $`\widehat{S}_1`$, and $`\widehat{S}_0`$ is required to have $`R=+2`$ also:
$$\widehat{W}=m_2\widehat{S}_2\widehat{S}_0+f\widehat{S}_1\widehat{S}_1\widehat{S}_0+h\widehat{S}_2\widehat{H}_u\widehat{H}_d.$$
(11)
The resulting scalar potential is
$$V=|m_2S_0+hH_uH_d|^2+|2fS_1S_0|^2+|m_2S_2+fS_1S_1|^2.$$
(12)
Let $`v_iS_i`$, then $`V=0`$ has the solution
$$v_0=0,v_2=\frac{fv_1^2}{m_2}.$$
(13)
The problem now is of course the indeterminate value of $`v_1`$. To fix $`v_1`$ and maintain the above seesaw structure while keeping $`v_0`$ zero, we add the following soft terms:
$`V^{}=m_{1}^{}{}_{}{}^{2}|S_1|^2[\lambda m_2S_2^{}S_1^2+h.c].`$ (14)
The equations of constraint for $`V+V^{}`$ to be a minimum are
$`0`$ $`=`$ $`m_2^2v_2+(f\lambda )m_2v_1^2,`$ (15)
$`0`$ $`=`$ $`m_{1}^{}{}_{}{}^{2}+4f^2v_0^2+2(f\lambda )m_2v_2+2f^2v_1^2,`$ (16)
$`0`$ $`=`$ $`v_0(m_2^2+4f^2v_1^2).`$ (17)
From Eq. (15), we find
$$v_2=\frac{(\lambda f)v_1^2}{m_2},$$
(18)
which indeed preserves the expected seesaw structure. From Eq. (17), we see that $`v_0=0`$ is still a solution, and from Eq. (16), taking into account Eq. (18), we find
$$v_1^2=\frac{m_{1}^{}{}_{}{}^{2}}{2\lambda (2f\lambda )},$$
(19)
where the denominator must be positive for $`V+V^{}`$ to be a minimum. The discrete $`Z_3`$ symmetry is broken spontaneously by $`v_1`$, hence a possible domain wall problem may appear. However, the Majorana fermion singlet $`\stackrel{~}{S}_1`$ may be given a mass $`m_1`$ which breaks the $`Z_3`$ symmetry softly but explicitly, thus avoiding such a problem.
Note that the scalars $`S_0`$ and $`S_2`$ remain heavy with mass $`m_2`$, but their VEV’s are zero or very small. The global $`U(1)_R`$ symmetry is broken by $`S_1`$ and $`S_2`$, hence the resulting Nambu-Goldstone boson is given by
$$\frac{(v_1)\sqrt{2}ImS_1+(2v_2)\sqrt{2}ImS_2}{\sqrt{v_1^2+4v_2^2}}.$$
(20)
In the couplings of $`S_2`$ to the superparticles of the MSSM, the axion enters as $`S_2`$ is replaced by
$$v_2e^{2i\phi }=v_2e^{2i\phi }\mathrm{exp}\left(\frac{ia\sqrt{2}}{v}\right),$$
(21)
where $`v=\sqrt{v_1^2+4v_2^2}`$ and the axion $`a`$ is given by
$$a=(\sqrt{2}v)[\phi \phi ],$$
(22)
with $`\phi =\theta _{QCD}/6`$. Thus the axion decay constant $`f_a`$ is $`\sqrt{2}v\sqrt{2}v_1`$ but $`M_{SUSY}`$ of the MSSM is $`v_2`$. This is analogous to the DFSZ model with $`M_{SUSY}`$ replaced by $`M_W`$.
In this model, the seesaw condition of Eq. (18) implies that $`M_{SUSY}f_a^2/M_U`$, where $`M_U`$ is a large effective mass parameter, i.e. $`2m_2/h(\lambda f)`$. The allowed range of values for $`f_a`$ from astrophysics and cosmology is between $`10^9`$ and $`10^{12}`$ GeV. Hence $`M_U`$ is between $`10^{15}`$ and $`10^{21}`$ GeV. Neutrino masses are given by either Eq. (5) in the Higgs triplet model, or Eq. (7) in the canonical seesaw model. There is no a priori connection between $`f_a`$ and $`m_\xi `$ or $`m_N`$. However, if they are of the same order of magnitude, then $`m_\nu `$ is inversely proportional to $`f_a`$ as proposed in the models of Ref..
The laboratory detection of axions depends on the $`a\gamma \gamma `$ coupling, which is proportional to
$$\frac{E}{N}\left(\frac{2}{3}\right)\frac{4+m_u/m_d+m_u/m_s}{1+m_u/m_d+m_u/m_s}=\frac{E}{N}1.92\pm 0.08,$$
(23)
where $`N`$ and $`E`$ are coefficients proportional to the color and electromagnetic anomalies of the axion. For the gluino axion model, $`N=6`$ but $`E=0`$ without or with neutrino mass from either the canonical seesaw or the Higgs triplet mechanism. This comes from the fact that, except for the gluino, every left-handed fermion has a right-handed partner of the same $`R`$.
In conclusion, we have incorporated neutrino masses (through the canonical seesaw or Higgs triplet mechanism) into the gluino axion model, using the superpotentials of Eq. (1) \[without the $`\mu `$ term\] and Eq. (11) with either Eq. (6) or Eq. (4). The $`\mu `$ term is replaced by $`h\widehat{S}_2`$, so that
$$\mu _{eff}=hS_2=hv_2e^{i\theta _{QCD}/3}.$$
(24)
Assuming that the intermediate scales, i.e. $`v_1`$ and $`m_N`$ or $`m_\xi `$ are of the same order, we then have
$$M_{SUSY}|\mu _{eff}|\frac{f_a^2}{M_U},m_\nu f_{eff}^2\frac{M_W^2}{f_a},$$
(25)
where $`f_{eff}`$ is a dimensionless coupling.
ACKNOWLEDGEMENT
We thank A. Yu. Smirnov for reading the manuscript. The research of E.M. was supported in part by the U. S. Department of Energy under Grant No. DE-FG03-94ER40837.
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# Higher-order nonlinear modes and bifurcation phenomena due to degenerate parametric four-wave mixing.
## I Introduction and Model
Recently parametric wave mixing in Kerr media has attracted significant attention (see, e.g., Refs. where continuous wave (CW) interaction and parametric self-trapping were investigated). This theoretical activity has been backed up by experimental advances, e.g., a novel scheme for quasi-phase matched third harmonic generation (THG) has been suggested . However, in previous works devoted to spatial solitary waves due to THG in planar waveguides, only families of fundamental self-trapped beams were considered. In the Letter , for example, where solitons due to third harmonic generation were considered for a bulk medium geometry, higher-order modes were not discussed in detail. By higher-order we refer to beam shapes whose transverse intensity typically has a multi-peaked structure and has higher energy than the single-peaked fundamental state.
In this work we analyze in some detail the structure and bifurcation phenomena of higher-order bright spatially localized modes or ‘solitons’, which we do not use in a strict mathematical sense, since the models in question are not integrable. The spatial configuration is assumed to be such that there is a well defined propagation direction and the beams are localized in $`n`$ transverse directions, with $`n=1`$ representing a planar waveguide and $`n=2`$ a bulk medium. Specifically we study models representing (1+1)-dimensional and (2+1)-dimensional, weakly-anisotropic media with cubic nonlinearity, under the phase-matched condition that the fundamental wave is resonantly coupled to its third harmonic. This is a particular degenerate case of solitons supported by the four-wave mixing processes , which is not completely understood yet in full generality. We assume that the interaction between the fundamental and third-harmonic waves includes the effects of parametric four-wave mixing, self-phase modulation, and cross-phase modulation.
We closely follow the derivation procedure of Ref. , assuming that the fundamental and the third-harmonic beams have the same linear polarization. The result is the following normalized (dimensionless) system of coupled equations:
$$\begin{array}{c}i\frac{u}{z}+^2uu+\left(\frac{1}{9}|u|^2+2|w|^2\right)u+\frac{1}{3}u^2w=0,\hfill \\ \\ i\sigma \frac{w}{z}+^2w\alpha w+(9|w|^2+2|u|^2)w+\frac{1}{9}u^3=0,\hfill \end{array}$$
(1)
where $`u`$ and $`w`$ are the fundamental and third harmonics, respectively. Also for the case of spatial beams $`^2^2/x^2+^2/y^2`$ in the (2+1)-dimensional case, or $`^2^2/x^2`$ in the (1+1)D case. The parameter $`\alpha `$ measures the shift in the propagation constant, which is induced by the nonlinearity and is also dependent on the quality of wave-vector matching between the harmonics, with $`\alpha =3\sigma `$ corresponding to exact matching, and $`z`$ is the propagation distance. For the spatial soliton case the dimensionless parameter $`\sigma `$ is the ratio of the wave numbers of the harmonics and is equal to $`3`$. Note that the system (1) may also describe temporal pulse propagation of resonantly interacting fundamental and third harmonics in optical fibers. For this physical situation $`^2^2/t^2`$ ($`t`$ is the retarded time variable) and $`\sigma `$ is the absolute value of the ratio of second-order group velocity dispersions for the first and the third harmonics and may be any positive number.
Radially symmetric stationary beams are described by real solutions $`u(r)`$ and $`w(r)`$ which are defined by the following system:
$$\begin{array}{c}\frac{d^2u}{dr^2}+\frac{s}{r}\frac{du}{dr}u+\left(\frac{1}{9}u^2+2w^2\right)u+\frac{1}{3}u^2w=0,\hfill \\ \frac{d^2w}{dr^2}+\frac{s}{r}\frac{dw}{dr}\alpha w+(9w^2+2u^2)w+\frac{1}{9}u^3=0.\hfill \end{array}$$
(2)
Here $`r\sqrt{x^2+y^2}`$ and $`s=1`$ for the (2+1)D case, whereas $`rx`$ and $`s=0`$ for the (1+1)D case. These localized solutions depend only on a single dimensionless parameter $`\alpha `$. Analysis of the linear part of Eqs. (2) in the limit $`r\mathrm{}`$ shows that conventional bright solitons (with exponentially decaying tails) can exist only for $`\alpha >0`$.
By ‘bright symmetric’ in the remainder of this paper we shall mean (2+1)D solutions when the intensity of each localized harmonic reaches a maximum at $`r=0`$ and (1+1)D solutions with $`u(r)=u(r)`$ and $`w(r)=w(r)`$. Thus, we shall only seek these solutions on the interval $`0r\mathrm{}`$ even in the $`(1+1)D`$ case. Note further that Eqs. (2) have odd symmetry, that is, if $`[u(r),w(r)]`$ is a solution then so is $`[u(r),w(r)]`$. Thus all solutions must come in pairs, the second solution being simply a change in sign (a phase shift of $`\pi `$) of both harmonics. For the case $`s=0`$, it is additionally possible to have solutions which are odd in both harmonics, or which are neither odd nor even. The latter type of solutions we shall refer to as being ‘bright asymmetric’. In this paper we shall consider mainly the solitons of bright symmetric type, but shall also present some results about bright asymmetric (1+1)D solitons. Dark solitons (localized solutions with nonzero asymptotics) are out of the scope of this paper.
The case $`\alpha <0`$ also has physical meaning, but there one should expect to find quasisolitons, which are almost localized stationary states that have small periodic oscillations in their tails. See e.g. for the definition, examples and for some issues surrounding them. Quasisolitons in this model will form the subject of another work. Here we shall concentrate almost exclusively on the case $`\alpha >0`$.
Using a direct analogy with the theory of $`\chi ^{(2)}`$ (second-harmonic generation or SHG) solitons (e.g. Ref. ), we start our analysis from the so-called cascading limit when $`\alpha 1`$. In this limit $`wu^3/(9\alpha )`$ and the equation for $`u`$ approaches the cubic-quintic Nonlinear Schrödinger (NLS) equation. This scalar equation possesses a familiar class of fundamental bright solitons consisting of a simple bell-shape \[there are also higher-order families in the (2+1)D case\]. These fundamental solitons can then be used as a starting point in the search for families of stationary solutions using numerical methods. These methods comprise a standard shooting method at fixed $`\alpha `$, and a continuation method allied to solution using a relaxation method for solving an appropriately defined two-point boundary-value problem for Eqs. (2). This latter technique can trace paths of solutions as $`\alpha `$ varies. We choose to characterize these solitons by the value of normalized total power which is one of the conserved quantities of the system (1)
$$P_{\mathrm{tot}}=_A(|u|^2+3\sigma |w|^2)𝑑A.$$
(3)
Here the integration extends over the appropriate one or two-dimensional infinite cross-section $`A`$.
The dependence of $`P_{\mathrm{tot}}`$ on $`\alpha `$ for a branch of solitons is usually, at least in the case of a fundamental solution, closely related to its stability. A necessary condition for stability in the case of fundamental multi-component solitons is typically given by a generalized Vakhitov-Kolokolov (VK) criterion , which often also appears to be a sufficient condition for soliton stability (see, e.g., Refs. ). However, the complexity of Eqs. (1) which, for example, possess collapse-type dynamics in the (2+1)D case, may lead to instability of fundamental solitons even for branches which are supposed to be stable according to the VK criterion . Thus, below we use power versus-$`\alpha `$ dependence only for classification of soliton families, leaving a full-scale stability analysis for future consideration.
## II Results for bulk media
First we present the results for the (2+1)D case. Figure 1 shows the variation of the normalized total power, $`P_{\mathrm{tot}}`$, with the normalized mismatch parameter $`\alpha `$, for different types of one-wave and two-wave localized solutions of the system (1) with $`s=1`$. The corresponding soliton profiles at various points along the presented branches are given in Figs. 27.
The first class of localized solutions of the system (1) consists of one-frequency soliton families for the third harmonic $`w_0`$, which exist for all $`\alpha >0`$ and represent scalar Kerr solitons described by the standard cubic NLS equation which follows from the second of Eqs. (1) at $`u=0`$:
$$\begin{array}{c}\frac{d^2w_0}{dr^2}+\frac{1}{r}\frac{dw_0}{dr}\alpha w_0+9w_0^3=0.\hfill \end{array}$$
(4)
These single frequency solitons differ from each other by the number of zero crossings in their radial profiles so that we denote the corresponding families as $`T_0`$ (no crossing), $`T_1`$ (one crossing), $`T_2`$ (two crossings), etc. Examples of one-wave solitons belonging to different $`T_j`$ families are shown in Fig. 2. Note that the normalized power $`P_{\mathrm{tot}}`$ is constant for each of the $`T_i`$ families. For example, for the fundamental one-wave soliton family $`T_0`$ (which are, in fact, Townes solitons of Ref. ) we have $`P_{tot}11.70`$ for all $`\alpha >0`$.
The second class of solutions to Eqs. (1) are genuinely two-wave bright solitons, described by families of localized beams with coupled fundamental and third harmonics. The simplest way to obtain such solutions numerically is to follow the two-wave soliton families from the cascading limit (large $`\alpha `$) as $`\alpha `$ decreases. In this work we concentrate on the result of following the lowest order two-wave soliton branch whose profiles have a simple one-hump form in the cascading limit. For this family, painstaking numerical continuation reveals a highly complex solution path involving restructuring of the soliton profile while the corresponding $`P(\alpha )`$ curve undergoes several loops (see Figs. 1 and LABEL:fig1-1). Inherent in each loop is a touch with one of the $`T_i`$ families. Such a touch corresponds to a transcritical bifurcation from the pure $`w`$-solution, and note \[for example from Fig. 3(c),(d) which correspond to points C and D on Fig. 1(a)\] that the two different bifurcating branches have opposite signs of their $`u`$-component. The fact that these bifurcations take place further illustrates the severity of the restructuring of the soliton profiles that must take place; in the cascading limit the branch is approximately of pure $`u`$-type, whereas at each bifurcation with $`T_i`$ it is composed of purely a third harmonic component $`w`$. Figures 27 illustrate the complete restructuring process by depicting the soliton profiles in the vicinity of each bifurcation and turning point of the $`P(\alpha )`$ curve. Note finally that the two-wave soliton family also includes the simplest so-called self-similar \[for which $`u(r)w(r)`$\] solution (Fig. 1, point M) at $`\alpha =1`$, see Ref. for the details and also Ref. , where its (1+1)-dimensional counterpart was also been considered.
The position of the bifurcation point from the $`T_0`$ branch can be approximately calculated analytically. Linearization of Eqs. (1) around the solution $`w_0(r)`$ gives the eigenvalue equation
$$\begin{array}{c}\frac{d^2u_1}{dr^2}+\frac{1}{r}\frac{du_1}{dr}+2w_0^2(r)u_1=\lambda u_1,\hfill \end{array}$$
(5)
together with appropriate boundary conditions. Bifurcations occur when $`\lambda =1`$. This may also be viewed as the problem of existence of localized states in the potential $`U(r)=2w_0^2(r)`$ with eigenvalue $`\lambda `$. Due to the lack of a closed form analytical expression for $`w_0^2(r)`$, solutions of (5) may be approximated by feeding in the numerical data for $`w_0`$ or by analytical techniques based on a variational approximation. Using the latter, based on a simple exponential trial function, gives the result $`w_0=\sqrt{8\alpha }/3e^{r\sqrt{\alpha }}`$. Substituting this into Eq. (5) and assuming a similar form of trial function for $`u_1`$, one can use a Rayleigh-Ritz method to obtain $`\alpha _{bif}^{(var)}=105.8`$. This agrees to within 2% with the numerical result $`\alpha _{bif}^{(num)}=104.2`$. Calculation of bifurcation points along the higher-order $`T_i`$ branches may in principle be carried out by the same method. However, this is less straightforward technically because it requires the use of complicated forms of trial functions, and it is perhaps more instructive to rely on numerical detection of the bifurcation points. Symmetry arguments dictate that at each bifurcation point $`\alpha _{bif}`$ there will exist two different bifurcating solutions of Eq. (5): $`u_1(r)`$ and $`u_1(r)`$. Moreover, each bifurcation is a transcritical and gives rise to a pair of two-wave branches $`(w_0(r),\pm \epsilon u_1(r))`$, where $`\epsilon `$ is proportional to $`|\alpha \alpha _{bif}|`$.
The third class of solutions to Eqs. (1) are the aforementioned quasisolitons which exist in the region of negative $`\alpha `$. We do not discuss quasi-solitons here in any detail. A full analysis will appear elsewhere. We simply make the comment that the branch SOP bifurcating from $`T_3`$ can be continued up to the boundary $`\alpha =0`$ separating regular from quasi-solitons. On the other side of the boundary a similar quasi-soliton state can be found with tiny oscillations in its tail \[see Fig. 1 and Fig. 7(t,s)\].
## III Results for planar waveguides
It is interesting to compare the (2+1)D results discussed above with those for the corresponding (1+1)D case. The bifurcation diagram related to the (1+1)D case is presented in Fig. 8 and the corresponding examples of soliton profiles are given in Fig. 914. We now highlight how, together with many obvious differences in comparison to the diagram for the (2+1)D case in Fig. 1, there are also some striking similarities as well. Note that in some respects the model for the (1+1)D case is simpler since the corresponding stationary system (2) with $`s=0`$ does not depend explicitly on $`r`$ and hence represents an autonomous dynamical system in four dimensions. Finding solitons is then reduced to finding homoclinic trajectories in this 4D phase space.
The first class of (1+1)D localized waves of system (1) consists of one-frequency soliton families for the third harmonic $`w_0`$, which exist for all $`\alpha >0`$ and represent scalar Kerr solitons described by the standard cubic (1+1)D NLS equation which follows from the second of Eqs. (1) at $`u=0`$:
$$\begin{array}{c}\frac{d^2w_0}{dx^2}\alpha w_0+9w_0^3=0.\hfill \end{array}$$
(6)
It can be readily solved exactly giving the well-known unique single soliton solution:
$$w_0(x)=\frac{\sqrt{2\alpha }}{3}\text{sech}(\sqrt{\alpha }x),P_{tot}=4\sqrt{\alpha }.$$
(7)
In contrast to the (2+1)D case strictly speaking there are no other one-wave localized solutions. However, it will be helpful in what follows to consider formal multi-soliton states consisting of a different number of infinitely separated single solitons (7), families of which we denote by $`S_1`$ (single soliton), $`S_2`$ (two solitons), $`S_3`$ (three solitons), etc. In this work we are mainly interested in families with an odd number of separated solitons: $`S_{2i+1},i=1,2,3,\mathrm{}`$, but we also investigate ‘bifurcations’ from $`S_2`$. Note that, for $`i>1`$, $`S_i`$ in fact denotes more than a single one-wave family, because each single pulse that is glued together can be either positive or negative.
The second class of (1+1)-dimensional localized solutions of Eqs. (1) consists of two-wave bright symmetric solitons and is described by families of localized beams with coupled fundamental and third harmonics. The simplest way to obtain the lowest order two-wave soliton family is again to continue numerically from solitons of the cascading limit ($`\alpha 1`$) given approximately by the expression:
$$\begin{array}{c}u(x)\frac{6}{\sqrt{1+B\mathrm{cosh}\mathrm{\hspace{0.33em}2}x}},wu^3/(9\alpha ),\hfill \end{array}$$
(8)
where $`B=\sqrt{1+16/\alpha }`$. The first expression for $`u(x)`$ in Eq. (8) is the solution of the corresponding cubic-quintic NLS-type equation.
The results of our numerical continuation from this limiting solution, upon decreasing $`\alpha `$ is that, like in the (2+1)D case, this branch also traces a convoluted path in the $`(P,\alpha )`$-plane, involving four ‘bifurcations’ from one-wave soliton families (from the families $`S_1`$, $`S_3`$, $`S_5`$, and $`S_7`$). As in the (2+1)D case, this branch connects to a self-similar solution at $`\alpha =1`$ (the point O in Fig. 8(b)). In this case, the self-similar solution is expressible in closed analytical form as
$$u(x)=a\mathrm{sech}x,w(x)=bu(x),$$
(9)
where the parameter $`b`$ is the real root of the cubic equation $`63b^33b^2+17b+1=0`$, and $`a^2=18/(18b^2+3b+1)`$. However, it is here that the similarity with the $`(2+1)`$-case ends, as we shall now explain.
First, let us try to motivate what is happening at each of the ‘bifurcations’ from $`S_j`$; for which at first sight it seems remarkable that each one occurs precisely at $`\alpha =9`$. Standard bifurcation analysis (e.g. as in Ref. ) allows us to find the position of the single bifurcation point from the one-wave soliton family $`S_1`$ (7) at $`\alpha =9.0`$ (point C in Fig. 8(a)). As in the (2+1)D case the bifurcation is a transcritical with one branch emerging to the left of the bifurcation point and one to the right. This structure is confirmed by the inset to Fig. 8(a) which shows that the branch emerging to the left undergoes a fold (at point B), so that on a larger scale both branches appear to bifurcate to the right.
Now it seems that this ‘local’ bifurcation from $`S_1`$ causes a topological change in the four-dimensional phase space so that a global event must also happen at this parameter value. This global event is the possibility of gluing together several copies of the $`S_1`$ back to back and forming a new branch of solitons with several large peaks that bifurcate from $`\alpha =9`$. Phenomenologically this is similar to what happens in the SHG case when the parameter equivalent to $`\alpha `$ passes through 1 . A key observation here is that in order to get a symmetric (even) solution, only an odd number of copies of the $`S_1`$ may be taken to form solitons in this way. As a convenient short-hand for this global bifurcation of multi-peaked solutions at $`\alpha =9`$, we have refereed to it as a local ‘bifurcation’ from $`S_{2i+1}`$ where $`i=1,2,3\mathrm{}`$, although this is strictly a misnomer.
Numerical continuation beyond point G of Fig. 8$`(a)`$ shows that two-wave soliton branch approaches $`\alpha =9.0`$ from the left, where it bifurcates from the $`S_3`$ asymptotic one-wave family that has alternative phase between each single-soliton component. However, we find that this is only one of a total of four symmetric two-wave solitons that come out of $`S_3`$. There are 8 in total if you include the change of sign of both $`u`$ and $`w`$. The second bifurcates to the left from the same (alternating phase) $`S_3`$ family and differs only in that the first harmonic has the opposite sign. A representative of this branch, corresponding to point $`H`$ in Fig. 8$`(a)`$ is shown in Fig. 10$`(h)`$. The two other branches exist for $`\alpha >9`$ and bifurcate from the $`S_3`$ family where all peaks are in phase (positive), and representatives are shown in Fig. 14$`(u,v)`$. With the increase of $`\alpha `$ (cascading limit) these complex multi-humped solitons keep their general structure intact, but become more localized. These two branches are not shown in the bifurcation diagram (Fig. 8) but their $`P(\alpha )`$ curves lie very close to each other and to the $`T_3`$ curve to the right of the bifurcation point.
A similar bifurcation picture is observed at $`\alpha =9.0`$ for bifurcations from $`S_5`$ and $`S_7`$ one-wave families. However, because of the increase of the number in possible one-wave multi-soliton families themselves, the number of the corresponding bifurcated two-component branches also increases. For the even solitons considered in this work we have the following formula to calculate the number of two-wave sub-families bifurcating from one-wave $`S_i`$ family: $`N_i=2^{(i+1)/2}`$ (double that if we count the change signs of $`u`$ and $`v`$). For example, there are 16 branches that bifurcate from $`S_7`$ branches which have $`P=84`$ at $`\alpha =9`$. Note that in the bifurcation diagram of Fig. 8, in order to clutter, only branches directly linked to the cascading limit two-wave family are shown. Close to bifurcation points, the third harmonic components of the depicted branches have neighboring humps of alternating sign and first harmonic components have all humps of the same sign. Note that these branches all bifurcate to the left of $`\alpha =9`$. For the branches which bifurcate to the right not all third harmonic neighboring humps alternate in sign.
It is important to note that none of the multi-hump soliton branches bifurcating to the left of $`\alpha =9`$ can be viewed as bound states of single partial solitons. Indeed, single one-hump solitons of Eqs. (1) always have $`u`$ and $`w`$ components in-phase (of the same sign) for $`\alpha <9.0`$, whereas some of the individual humps of the of multi-hump structures bifurcating to the left from $`S_i`$ ($`i>1`$) families have $`u`$ and $`w`$ components of different signs. To illustrate this point we show in Fig. 15 an enlarged bifurcation diagram in the vicinity of $`\alpha =9`$ covering the first three families, $`S_i`$, $`i=1,2,3`$. Some of the corresponding examples of soliton profiles plotted at $`\alpha =8.6`$ are given in Fig. 16. As they approach $`\alpha =9.0`$ the separation between each individual hump (a ‘partial soliton’) increases and the state begins to approach a concatenation of single solitons with slightly overlapping tails. However, some of these partial solitons have out-of-phase $`u`$ and $`w`$ components and hence cannot exist on their own (i.e. without being in superposition with other ‘partial’ solitons).
Fig. 15 shows something even more striking - that there is also a ‘bifurcation’ from the $`S_2`$ family. However, the solitary waves that bifurcate from there are not bright symmetric but in fact are asymmetric solitons, see Fig. 17. Also at least one of these asymmetric solutions is born in a symmetry-breaking (pitchfork) bifurcation from one of the symmetric soliton branches (at the point $`O_{as}`$, see Fig. 15). Thus there is a branch of asymmetric solitons which connects symmetric solitons with a branch of asymptotic antisymetric solitons (the $`S_2`$ family). We conjecture that there are similar asymmetric solitons that ‘bifurcate’ from $`S_j`$ at $`\alpha =9`$ for all even $`j`$.
In contrast to the (2+1)D case, we have found no examples (at least considering all bifurcations from $`S_{2i+1}`$ with $`2i+17`$) of two-wave solitons that survive down to $`\alpha =0`$ where they might form a connection with branches of quasi-solitons existing for $`\alpha <0`$. Instead a representative branch coming from $`T_7`$ bends abruptly (at R) at which point $`\alpha `$ increases through the point S until it reaches T at $`\alpha 3.65`$, where another nonlocal bifurcation occurs. In this process, the third harmonic gradually forms a core with weakly separated wings. At T ,The latter become completely separated one-wave solitons \[see Fig. 13(s,t)\]. The solution at the point T can thus be viewed as a direct sum of two well-separated one-wave solitons and the soliton at point N. Beyond $`T`$ we were unable to find any similar solutions. This non-trivial “jump” bifurcation is indicated by the vertical arrow in Fig. 8.
## IV Conclusion
In conclusion, we have investigated and classified higher-order soliton families and bifurcation phenomena due to resonant parametric interaction of a fundamental frequency wave with its third harmonic.
In the case of (2+1)D solitons the picture is consistent with standard theories, albeit the branch we followed from the cascading limit connects several distinct soliton types in a non-trivial way. Also the structure of the sets of branches we found to approach the limit $`\alpha =0`$ could do with further investigation, perhaps using singular perturbation theory. The relation of these states for positive $`\alpha `$ to quasi-solitons for negative $`\alpha `$ will be addressed elsewhere.
In contrast, in the (1+1)D case the bifurcation diagram is less clear-cut and we have found at least two novel features (i) the non-local bifurcation of multi-humped two-frequency solutions which are a consequence of the local bifurcation from the one-humped one-frequency soliton at $`\alpha =9`$, and (ii) the so-called jump bifurcation at the point $`T`$. The first of these is particularly intriguing since not only are symmetric multi-humped states formed in this way, but also asymmetric ones. Also some of the multi-humped states cannot be viewed as bound states of several distinct one-humped states. The second novel bifurcation, the jump, appears related to, but not the same as, the so-called orbit-flip bifurcation . A dynamical-systems-theory explanation of these new bifurcation events, perhaps using the Lin-Sandstede method as in Ref. , would be most interesting.
Stability of the newly discovered soliton families remains an open question, especially for the (1+1)D case. Although usually higher-order soliton families are subject to one of several types of instability, some exceptions are known (see, e.g., ) and thus a careful stability analysis is worth doing. The promise of detecting stable multi-hump solitons is real indeed because at least some of them cannot be viewed as bound states of two or more single (one-hump) solitons. For such bound state solitons of NLS-type system of equations, there is practically no hope of stability as shown e.g. in Ref. .
The authors acknowledge the use of computing facilities at Optical Sciences Centre, RSPhysSE, the Australian National University. AVB and RAS are indebted to O. Bang, L. Berge, B. A. Malomed, and D. Skryabin for useful discussions and interest in this work. ARC is indebted to the hospitality of Yu. Kivshar at the RSPhysSE and to the UK EPSRC with whom he holds an advanced fellowship.
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# New Mechanism For Mass Generation of Gauge Field This work is suppoeted by China National Natural Science Foundation and China Postdoctorial Science Foundation.
## 1 Introduction
Now, it is generally believed that four kinds of fundamental interactions in Nature are all gauge interactions and can be described by gauge field theory. From theoretical point of view, the principle of local gauge invariance plays a fundamental role in particle’s interaction theory. According to experimental results, some gauge bosons are massive . A usual way to make gauge field obtain non-zero mass is to use spontaneously symmetry breaking and Higgs mechanism , which is well-known in constructing the standard model . But the Higgs mechanism is not the only mechanism that make gauge field massive. In this paper, another mechanism for the mass generation of gauge field is proposed. By introducing two sets of gauge fields, we will introduce the mass term of gauge fields without violating the local gauge symmetry of the Lagrangian. Because the Lagrangian has strict local gauge symmetry, the model is renormalizable.
## 2 The Lagrangian of The Model
Suppose that the gauge symmetry of the theory is $`SU(N)`$ group, $`\psi (x)`$ is a N-component vector in the fundamental representative space of $`SU(N)`$ group and the representative matrices of the generators of $`SU(N)`$ group are denoted by $`T_i(i=1,2,\mathrm{},N^21)`$. They are Hermit and traceless. They satisfy:
$$[T_i,T_j]=if_{ijk}T_k,Tr(T_iT_j)=\delta _{ij}K,$$
$`(2.1)`$
where $`f_{ijk}`$ are structure constants of $`SU(N)`$ group, K is a constant which is independent of indices $`i`$ and $`j`$ but depends on the representation of the group. The representative matrix of a general element of the $`SU(N)`$ group is expressed as:
$$U=e^{i\alpha ^iT_i}$$
$`(2.2)`$
with $`\alpha ^i`$ the real group parameters. In global gauge transformations, all $`\alpha ^i`$ are independent of space-time coordinates, while in local gauge transformations, $`\alpha ^i`$ are functions of space-time coordinates. $`U`$ is a unitary $`N\times N`$ matrix.
In order to introduce the mass term of gauge fields without violating local gauge symmetry, two kinds of gauge fields $`A_\mu (x)`$ and $`B_\mu (x)`$ are needed. $`A_\mu (x)`$ and $`B_\mu (x)`$ are vectors in the canonical representative space of $`SU(N)`$ group. They can be expressed as linear combinations of generators :
$$A_\mu (x)=A_\mu ^i(x)T_i$$
$`(2.3a)`$
$$B_\mu (x)=B_\mu ^i(x)T_i.$$
$`(2.3b)`$
where $`A_\mu ^i(x)`$ and $`B_\mu ^i(x)`$ are component fields of gauge fields $`A_\mu (x)`$ and $`B_\mu (x)`$ respectively.
Corresponding to two kinds of gauge fields, there are two kinds of gauge covariant derivatives in the theory:
$$D_\mu =_\mu igA_\mu $$
$`(2.4a)`$
$$D_{b\mu }=_\mu +i\alpha gB_\mu .$$
$`(2.4b)`$
The strengths of gauge fields $`A_\mu (x)`$ and $`B_\mu (x)`$ are defined as
$$\begin{array}{ccc}A_{\mu \nu }& =& \frac{1}{ig}[D_\mu ,D_\nu ]\hfill \\ & =& _\mu A_\nu _\nu A_\mu ig[A_\mu ,A_\nu ]\hfill \end{array}$$
$`(2.5a)`$
$$\begin{array}{ccc}B_{\mu \nu }& =& \frac{1}{i\alpha g}[D_{b\mu },D_{b\nu }]\hfill \\ & =& _\mu B_\nu _\nu B_\mu +i\alpha g[B_\mu ,B_\nu ].\hfill \end{array}$$
$`(2.5b)`$
respectively. Similarly, $`A_{\mu \nu }`$ and $`B_{\mu \nu }`$ can also be expressed as linear combinations of generators:
$$A_{\mu \nu }=A_{\mu \nu }^iT_i$$
$`(2.6a)`$
$$B_{\mu \nu }=B_{\mu \nu }^iT_i.$$
$`(2.6b)`$
Using relations (2.1) and (2.5a,b), we can obtain
$$A_{\mu \nu }^i=_\mu A_\nu ^i_\nu A_\mu ^i+gf^{ijk}A_\mu ^jA_\nu ^k$$
$`(2.7a)`$
$$B_{\mu \nu }^i=_\mu B_\nu ^i_\nu B_\mu ^i\alpha gf^{ijk}B_\mu ^jB_\nu ^k.$$
$`(2.7b)`$
The Lagrangian density of the model is
$$\begin{array}{ccc}& =& \overline{\psi }(\gamma ^\mu D_\mu +m)\psi \frac{1}{4K}Tr(A^{\mu \nu }A_{\mu \nu })\frac{1}{4K}Tr(B^{\mu \nu }B_{\mu \nu })\hfill \\ & & \frac{\mu ^2}{2K(1+\alpha ^2)}Tr\left[(A^\mu +\alpha B^\mu )(A_\mu +\alpha B_\mu )\right]\hfill \end{array}$$
$`(2.8)`$
where $`\alpha `$ is a constant. In this paper, the space-time metric is selected as $`\eta _{\mu \nu }=diag(1,1,1,1)`$, $`(\mu ,\nu =0,1,2,3)`$. According to relation (2.1), the above Lagrangian density $``$ can be rewritten as:
$$\begin{array}{ccc}& =& \overline{\psi }[\gamma ^\mu (_\mu igA_\mu ^iT_i)+m]\psi \frac{1}{4}A^{i\mu \nu }A_{\mu \nu }^i\frac{1}{4}B^{i\mu \nu }B_{\mu \nu }^i\hfill \\ & & \frac{\mu ^2}{2(1+\alpha ^2)}(A^{i\mu }+\alpha B^{i\mu })(A_\mu ^i+\alpha B_\mu ^i).\hfill \end{array}$$
$`(2.9)`$
An obvious characteristic of the above Lagrangian is that the mass term of the gauge fields is introduced into the Lagrangian and this term does not affect the symmetry of the Lagrangian. We will prove that the above Lagrangian has strict local gauge symmetry in the chapter 4.
Because both vector fields $`A_\mu `$ and $`B_\mu `$ are standard gauge fields, this model is a kind of gauge field model which describes gauge interactions between gauge fields and matter fields.
## 3 Global Gauge Symmetry and Conserved Charges
Now, let’s discuss the gauge symmetry of the Lagrangian density $``$. First, we will discuss the global gauge symmetry and the corresponding conserved charges. In global gauge transformation, the matter field $`\psi `$ transforms as:
$$\psi \psi ^{}=U\psi ,$$
$`(3.1)`$
where $`U`$ is independent of space-time coordinates. That is
$$_\mu U=0.$$
$`(3.2)`$
The corresponding global gauge transformations of gauge fields $`A_\mu `$ and $`B_\mu `$ are
$$A_\mu UA_\mu U^{}$$
$`(3.3a)`$
$$B_\mu UB_\mu U^{}$$
$`(3.3b)`$
respectively. It is easy to prove that
$$D_\mu UD_\mu U^{}$$
$`(3.4a)`$
$$D_{b\mu }UD_{b\mu }U^{}$$
$`(3.4b)`$
$$A_{\mu \nu }UA_{\mu \nu }U^{}$$
$`(3.5a)`$
$$B_{\mu \nu }UB_{\mu \nu }U^{}$$
$`(3.5b)`$
Using all the above transformation relations, it can be strictly proved that all terms in eq(2.8) are gauge invariant. So, the whole Lagrangian density has global gauge symmetry.
Let $`\alpha ^i`$ in eq(2.2) be the first order infinitesimal parameters, then, in the first order approximation, the transformation matrix $`U`$ can be rewritten as:
$$U1i\alpha ^iT^i.$$
$`(3.6)`$
The first order infinitesimal variations of fields $`\psi ,\overline{\psi },A_\mu `$ and $`B_\mu `$ are
$$\delta \psi =i\alpha ^iT^i\psi $$
$`(3.7a)`$
$$\delta \overline{\psi }=i\alpha ^i\overline{\psi }T^i$$
$`(3.7b)`$
$$\delta A_\mu =\alpha ^if^{ijk}A_\mu ^jT^k$$
$`(3.8a)`$
$$\delta B_\mu =\alpha ^if^{ijk}B_\mu ^jT^k$$
$`(3.8b)`$
respectively. From eqs(3.8a,b) and eqs(2.3a,b), we can obtain
$$\delta A_\mu ^k=\alpha ^if^{ijk}A_\mu ^j$$
$`(3.9a)`$
$$\delta B_\mu ^k=\alpha ^if^{ijk}B_\mu ^j.$$
$`(3.9b)`$
The first order variation of the Lagrangian density is
$$\begin{array}{ccc}\delta & =& _\mu \left(\frac{}{_\mu \psi }\delta \psi +\delta \overline{\psi }\frac{}{_\mu \overline{\psi }}+\frac{}{_\mu A_\nu ^k}\delta A_\nu ^k+\frac{}{_\mu B_\nu ^k}\delta B_\nu ^k\right)\hfill \\ & =& \alpha ^i^\mu J_\mu ^i,\hfill \end{array}$$
$`(3.10)`$
where
$$J_\mu ^i=i\overline{\psi }\gamma _\mu T^i\psi f^{ijk}A^{j\nu }A_{\mu \nu }^kf^{ijk}B^{j\nu }B_{\mu \nu }^k.$$
$`(3.11)`$
The conserved current can also be written as
$$\begin{array}{ccc}J_\mu & =& i\overline{\psi }\gamma _\mu T^i\psi T^i+i[A^\nu ,A_{\mu \nu }]+i[B^\nu ,B_{\mu \nu }]\hfill \\ & =& J_\mu ^iT^i.\hfill \end{array}$$
$`(3.12)`$
Because the Lagrangian density $``$ has global gauge symmetry, the variation of $``$ under global gauge transformations vanishes. That is
$$\delta =0.$$
$`(3.13)`$
Because $`\alpha ^i`$ are arbitrary global parameters, from eq(3.10), we can obtain the following conservation equation:
$$^\mu J_\mu ^i=0.$$
$`(3.14)`$
The corresponding conserved charges are
$$\begin{array}{ccc}Q^i& =& d^3xJ^{i0}\hfill \\ & =& d^3x(\psi ^{}T^i\psi +[A_j,A^{j0}]^i+[B_j,B^{j0}]^i).\hfill \end{array}$$
$`(3.15)`$
After quantization, $`Q^i`$ are generators of gauge transformation. An important feature of the above relation is that no matter what the value of parameter $`\alpha `$ is, gauge fields $`A_\mu `$ and $`B_\mu `$ contribute the same terms to the conserved currents and conserved charges.
## 4 Local Gauge Symmetry
If $`U`$ in eq(3.1) depends on space-time coordinates, the transformation of eq(3.1) is a local $`SU(N)`$ gauge transformation. In this case,
$$_\mu U0,_\mu \alpha ^i0$$
$`(4.1)`$
The corresponding transformations of gauge fields $`A_\mu `$ and $`B_\mu `$ are
$$A_\mu UA_\mu U^{}\frac{1}{ig}U_\mu U^{}$$
$`(4.2a)`$
$$B_\mu UB_\mu U^{}+\frac{1}{i\alpha g}U_\mu U^{}$$
$`(4.2b)`$
respectively.
Using above transformation relations, it is easy to prove that
$$D_\mu UD_\mu U^{}$$
$`(4.3a)`$
$$D_{b\mu }UD_{b\mu }U^{}.$$
$`(4.3b)`$
Therefore,
$$A_{\mu \nu }UA_{\mu \nu }U^{}$$
$`(4.4a)`$
$$B_{\mu \nu }UB_{\mu \nu }U^{}$$
$`(4.4b)`$
$$D_\mu \psi UD_\mu \psi $$
$`(4.5)`$
$$A_\mu +\alpha B_\mu U(A_\mu +\alpha B_\mu )U^{}$$
$`(4.6)`$
It can be strictly proved that the Lagrangian density $``$ defined by eq(2.8) is invariant under the above local $`SU(N)`$ gauge transformations. Therefore the model has strict local gauge symmetry.
An obvious characteristics of this gauge field theory is that two different gauge fields $`A_\mu `$ and $`B_\mu `$ which correspond to one gauge symmetry are introduced into the theory. From eq(4.2a,b), we know that both gauge fields $`A_\mu `$ and $`B_\mu `$ are standard gauge fields. But, they have different roles in theory. It is known that, in Yang-Mills theory, gauge field can be regarded as gauge compensatory field of matter fields. In other words, if there were no gauge field, though the Lagrangian could have global gauge symmetry, it would have no local gauge symmetry. In order to make the Lagrangian have local gauge symmetry, we must introduce gauge field and make the variation of gauge field under local gauge transformations compensate the variation of the kinematical terms of matter fields. So, the form of local gauge transformation of gauge field is determined by the form of local gauge transformation of matter fields. Similar case holds in the gauge field theory which is discussed in this paper: gauge field $`A_\mu `$ can regarded as gauge compensatory field of matter fields and gauge field $`B_\mu `$ can be regarded as gauge compensatory field of gauge field $`A_\mu `$. Therefore, in this gauge field theory, the form of the local gauge transformation of gauge field $`A_\mu `$ is determined by the form of local gauge transformation of matter fields, and the form of local gauge transformation of gauge field $`B_\mu `$ is determined by the form of local gauge transformation of gauge field $`A_\mu `$. And because of the compensation of gauge field $`B_\mu `$, the mass term of gauge field can be introduced into the Lagrangian without violating its local gauge symmetry.
## 5 The Masses of Gauge Fields
The mass term of gauge fields can be written as:
$$(A^\mu ,B^\mu )M\left(\begin{array}{c}A_\mu \\ B_\mu \end{array}\right).$$
$`(5.1)`$
where $`M`$ is the mass matrix:
$$M=\frac{1}{1+\alpha ^2}\left(\begin{array}{cc}\mu ^2& \alpha \mu ^2\\ \alpha \mu ^2& \alpha ^2\mu ^2\end{array}\right).$$
$`(5.2)`$
Generally speaking, physical particles generated from gauge interactions are eigenvectors of mass matrix and the corresponding masses of these particles are eigenvalues of mass matrix. $`M`$ has two eigenvalues, they are
$$m_1^2=\mu ^2,m_2^2=0.$$
$`(5.3)`$
The corresponding eigenvectors are
$$\left(\begin{array}{c}\mathrm{cos}\theta \\ \mathrm{sin}\theta \end{array}\right),\left(\begin{array}{c}\mathrm{sin}\theta \\ \mathrm{cos}\theta \end{array}\right),$$
$`(5.4)`$
where,
$$\mathrm{cos}\theta =\frac{1}{\sqrt{1+\alpha ^2}},\mathrm{sin}\theta =\frac{\alpha }{\sqrt{1+\alpha ^2}}.$$
$`(5.5)`$
Define
$$C_\mu =\mathrm{cos}\theta A_\mu +\mathrm{sin}\theta B_\mu $$
$`(5.6a)`$
$$F_\mu =\mathrm{sin}\theta A_\mu +\mathrm{cos}\theta B_\mu .$$
$`(5.6b)`$
It is easy to see that $`C_\mu `$ and $`F_\mu `$ are eigenstates of mass matrix, they describe those particles generated from gauge interactions. The inverse transformations of (5.6a,b) are
$$A_\mu =\mathrm{cos}\theta C_\mu \mathrm{sin}\theta F_\mu $$
$`(5.7a)`$
$$B_\mu =\mathrm{sin}\theta C_\mu +\mathrm{cos}\theta F_\mu .$$
$`(5.7b)`$
Then the Lagrangian density $``$ given by (2.9) changes into:
$$=_0+_I,$$
$`(5.8)`$
where
$$_0=\overline{\psi }(\gamma ^\mu _\mu +m)\psi \frac{1}{4}C_0^{i\mu \nu }C_{0\mu \nu }^i\frac{1}{4}F_0^{i\mu \nu }F_{0\mu \nu }^i\frac{\mu ^2}{2}C^{i\mu }C_\mu ^i.$$
$`(5.9a)`$
$$\begin{array}{ccc}_I& =& ig\overline{\psi }\gamma ^\mu (\mathrm{cos}\theta C_\mu \mathrm{sin}\theta F_\mu )\psi \hfill \\ & & \frac{\mathrm{cos2}\theta }{2\mathrm{c}\mathrm{o}\mathrm{s}\theta }gf^{ijk}C_0^{i\mu \nu }C_\mu ^jC_\nu ^k+\frac{\mathrm{sin}\theta }{2}gf^{ijk}F_0^{i\mu \nu }F_\mu ^jF_\nu ^k\hfill \\ & & +\frac{\mathrm{sin}\theta }{2}gf^{ijk}F_0^{i\mu \nu }C_\mu ^jC_\nu ^k+g\mathrm{sin}\theta f^{ijk}C_0^{i\mu \nu }C_\mu ^jF_\nu ^k\hfill \\ & & \frac{1\frac{3}{4}\mathrm{sin}^22\theta }{4\mathrm{c}\mathrm{o}\mathrm{s}^2\theta }g^2f^{ijk}f^{ilm}C_\mu ^jC_\nu ^kC^{l\mu }C^{m\nu }\hfill \\ & & \frac{\mathrm{sin}^2\theta }{4}g^2f^{ijk}f^{ilm}F_\mu ^jF_\nu ^kF^{l\mu }F^{m\nu }+g^2\mathrm{tg}\theta \mathrm{cos2}\theta f^{ijk}f^{ilm}C_\mu ^jC_\nu ^kC^{l\mu }F^{m\nu }\hfill \\ & & \frac{\mathrm{sin}^2\theta }{2}g^2f^{ijk}f^{ilm}(C_\mu ^jC_\nu ^kF^{l\mu }F^{m\nu }+C_\mu ^jF_\nu ^kF^{l\mu }C^{m\nu }+C_\mu ^jF_\nu ^kC^{l\mu }F^{m\nu }).\hfill \end{array}$$
$`(5.9b)`$
In the above relations, we have used the following simplified notations:
$$C_{0\mu \nu }^i=_\mu C_\nu ^i_\nu C_\mu ^i$$
$`(5.10a)`$
$$F_{0\mu \nu }^i=_\mu F_\nu ^i_\nu F_\mu ^i$$
$`(5.10b)`$
From eq(5.9a), it is easy to see that the mass of field $`C_\mu `$ is $`\mu `$ and the mass of gauge field $`F_\mu `$ is zero. That is
$$m_c=\mu ,m_F=0.$$
$`(5.11)`$
Transformations (5.7a,b) are pure algebraic operations which do not affect the gauge symmetry of the Lagrangian. They can be regarded as redefinitions of gauge fields. The local gauge symmetry of the Lagrangian is still strictly preserved after field transformations. In another words, the symmetry of the Lagrangian before transformations is completely the same as the symmetry of the Lagrangian after transformations. In fact, we do not introduce any kind of symmetry breaking in the whole paper.
Fields $`C_\mu `$ and $`F_\mu `$ are linear combinations of gauge fields $`A_\mu `$ and $`B_\mu `$, so the forms of local gauge transformations of fields $`C_\mu `$ and $`F_\mu `$ are determined by the forms of local gauge transformations of gauge fields $`A_\mu `$ and $`B_\mu `$. Because $`C_\mu `$ and $`F_\mu `$ consist of gauge fields $`A_\mu `$ and $`B_\mu `$ and transmit gauge interactions between matter fields, for the sake of simplicity, we also call them gauge field, just as we call $`W^\pm `$ and $`Z^0`$ gauge fields in electroweak model. Therefor, two different kinds of force-transmitting vector fields exist in this gauge field theory: one is massive and another is massless.
## 6 Equation of Motion
The Euler-Lagrange equation of motion for fermion field can be deduced from eq(5.8):
$$[\gamma ^\mu (_\mu ig\mathrm{cos}\theta C_\mu +ig\mathrm{sin}\theta F_\mu )+m]\psi =0.$$
$`(6.1)`$
If we deduce the Euler-Lagrange equations of motion of gauge fields from eq(5.8), we will obtain very complicated expressions. For the sake of simplicity, we deduce the equations of motion of gauge fields from eq(2.8). In this case, the equations of motion of gauge fields $`A_\mu `$ and $`B_\mu `$ are:
$$D^\mu A_{\mu \nu }\frac{\mu ^2}{1+\alpha ^2}(A_\nu +\alpha B_\nu )=ig\overline{\psi }\gamma _\nu T^i\psi T^i$$
$`(6.2a)`$
$$D_b^\mu B_{\mu \nu }\frac{\alpha \mu ^2}{1+\alpha ^2}(A_\nu +\alpha B_\nu )=0$$
$`(6.2b)`$
respectively. In the above relations, we have used two simplified notations:
$$D^\mu A_{\mu \nu }=[D^\mu ,A_{\mu \nu }]$$
$`(6.3a)`$
$$D_b^\mu B_{\mu \nu }=[D_b^\mu ,B_{\mu \nu }].$$
$`(6.3b)`$
Eqs(6.2a,b) can be expressed in terms of component fields $`A_\mu ^i`$ and $`B_\mu ^i`$:
$$^\mu A_{\mu \nu }^i\frac{\mu ^2}{1+\alpha ^2}(A_\nu ^i+\alpha B_\nu ^i)=ig\overline{\psi }\gamma _\nu T^i\psi +gf^{ijk}A_{\mu \nu }^jA^{k\nu }$$
$`(6.4a)`$
$$^\mu B_{\mu \nu }^i\frac{\alpha \mu ^2}{1+\alpha ^2}(A_\nu ^i+\alpha B_\nu ^i)=\alpha gf^{ijk}B_{\mu \nu }^jB^{k\nu }$$
$`(6.4b)`$
The equations of motion for gauge fields $`C_\mu `$ and $`F_\mu `$ can be easily obtained from eqs(6.4a,b). In other words, cos$`\theta `$ (6.4a) – sin$`\theta `$(6.4b) gives the equation of motion for gauge field $`C_\mu `$, and – sin$`\theta `$ (6.4a) + cos$`\theta `$(6.4b) gives the equation of motion for gauge field $`F_\mu `$.
From eq(6.2a) or (6.2b), we can obtain a supplementary condition. Using eq(6.1), we can prove that
$$[D^\lambda ,ig\overline{\psi }\gamma _\lambda T^i\psi T^i]=0.$$
$`(6.5)`$
Let $`D^\nu `$ act on eq(6.2a) from the left, and let $`D_b^\nu `$ act on eq(6.2b) from the left, applying eq(5.4) and the following two identities:
$$[D^\lambda ,[D^\nu ,A_{\nu \lambda }]]=0$$
$`(6.6a)`$
$$[D_b^\lambda ,[D_b^\nu ,B_{\nu \lambda }]]=0,$$
$`(6.6b)`$
we can obtain the following two equations
$$[D^\nu ,A_\nu +\alpha B_\nu ]=0$$
$`(6.7a)`$
$$[D_b^\nu ,A_\nu +\alpha B_\nu ]=0$$
$`(6.7b)`$
respectively. These two equations are essentially the same, they give a supplementary condition. If we expressed eqs(6.7a,b) in terms of component fields, these two equations will give the same expression:
$$^\nu (A_\nu ^i+\alpha B_\nu ^i)+\alpha gf^{ijk}A_\nu ^jB^{k\nu }=0.$$
$`(6.8)`$
When $`\nu =0`$, eqs(6.4a,b) don’t give dynamical equations for gauge fields, because they contain no time derivative terms. They are just constrains. Originally, gauge fields $`A_\mu ^i`$ and $`B_\mu ^i`$ have $`8(N^21)`$ degrees of freedom, but they satisfy $`2(N^21)`$constrains and have $`(N^21)`$ gauge degrees of freedom, therefore, gauge fields $`A_\mu ^i`$ and $`B_\mu ^i`$ have $`5(N^21)`$ independent dynamical degrees of freedom altogether. This result coincides with our experience: a massive vector field has 3 independent degrees of freedom and a massless vector field has 2 independent degrees of freedom.
## 7 The Case That Matter Fields Are Scalar Fields
In the above discussions, matter fields are spinor fields. Now, let’s consider the case that matter fields are scalar fields. Suppose that there are $`N`$ scalar fields $`\phi _l(x)(l=1,2,\mathrm{}N)`$ which form a multiplet of matter fields:
$$\phi (x)=\left(\begin{array}{c}\phi _1(x)\\ \phi _2(x)\\ \mathrm{}\\ \phi _N(x)\end{array}\right)$$
$`(7.1)`$
All $`\phi (x)`$ form a fundamental representative space of $`SU(N)`$ group. In $`SU(N)`$ gauge transformation, $`\phi (x)`$ transforms as :
$$\phi (x)\phi ^{}(x)=U\phi (x)$$
$`(7.2)`$
The Lagrangian density is
$$\begin{array}{ccc}& =& [(_\mu igA_\mu )\phi ]^+(^\mu igA^\mu )\phi V(\phi )\hfill \\ & & \frac{1}{4K}Tr(A^{\mu \nu }A_{\mu \nu })\frac{1}{4K}Tr(B^{\mu \nu }B_{\mu \nu })\hfill \\ & & \frac{\mu ^2}{2K(1+\alpha ^2)}Tr\left[(A^\mu +\alpha B^\mu )(A_\mu +\alpha B_\mu )\right]\hfill \end{array}$$
$`(7.3)`$
The above Lagrangian density can be expressed in terms of component fields :
$$\begin{array}{ccc}& =& [(_\mu igA_\mu ^iT_i)\phi ]^+(^\mu igA^{i\mu }T_i)\phi V(\phi )\hfill \\ & & \frac{1}{4}A^{i\mu \nu }A_{\mu \nu }^i\frac{1}{4}B^{i\mu \nu }B_{\mu \nu }^i\hfill \\ & & \frac{\mu ^2}{2(1+\alpha ^2)}(A^{i\mu }+\alpha B^{i\mu })(A_\mu ^i+\alpha B_\mu ^i)\hfill \end{array}$$
$`(7.4)`$
The general form for $`V(\phi )`$ which is renormalizable and gauge invariant is
$$V(\phi )=m^2\phi ^+\phi +\lambda (\phi ^+\phi )^2.$$
$`(7.5)`$
It is easy to prove that the Lagrangian density $``$ defined by eq(7.3) has local $`SU(N)`$ gauge symmetry. The Euler-Lagrange equation of motion for scalar field $`\phi `$ is:
$$(^\mu igA^\mu )(_\mu igA_\mu )\phi m^2\phi 2\lambda \phi (\phi ^+\phi )^2=0$$
$`(7.6)`$
If $`N^21`$ scalar fields $`\phi _l(x)(l=1,2,\mathrm{}N^21)`$ form a multiplet of matter fields
$$\phi (x)=\phi _l(x)T_l,$$
$`(7.7)`$
then, the gauge transformation of $`\phi (x)`$ is
$$\phi (x)\phi ^{}(x)=U\phi (x)U^+.$$
$`(7.8)`$
All $`\phi (x)`$ form a space of adjoint representation of $`SU(N)`$ group. In this case, the gauge covariant derivative is
$$D_\mu \phi =_\mu \phi ig[A_\mu ,\phi ],$$
$`(7.9)`$
and the gauge invariant Lagrangian density $``$ is
$$\begin{array}{ccc}& =& \frac{1}{K}Tr[(D^\mu \phi )^+(D_\mu \phi )]V(\phi )\hfill \\ & & \frac{1}{4}A^{i\mu \nu }A_{\mu \nu }^i\frac{1}{4}B^{i\mu \nu }B_{\mu \nu }^i\hfill \\ & & \frac{\mu ^2}{2(1+\alpha ^2)}(A^{i\mu }+\alpha B^{i\mu })(A_\mu ^i+\alpha B_\mu ^i).\hfill \end{array}$$
$`(7.10)`$
## 8 A More General Model
In the above discussions, a gauge field model, which has strict local $`SU(N)`$ gauge symmetry and contains massive gauge bosons, is constructed. In the above model, only gauge field $`A_\mu `$ directly interacts with matter fields $`\psi `$ or $`\phi `$, gauge field $`B_\mu `$ doesn’t directly interact with matter fields. But this restriction is not necessary in constructing the model. In this section, we will construct a more general gauge field model, in which both gauge fields interact with matter fields in the original Lagrangian. As an example, we only discuss the case that matter fields are spinor fields. The case that matter fields are scalar fields can be discussed similarly.
In chapter 4, we have prove that, under local gauge transformations, $`D_\mu `$ and $`D_{b\mu }`$ transform covariantly. It is easy to prove that $`\mathrm{cos}^2\varphi D_\mu +\mathrm{sin}^2\varphi D_{b\mu }`$ is the most general gauge covariant derivative which transforms covariantly under local $`SU(N)`$ gauge transformations
$$\mathrm{cos}^2\varphi D_\mu +\mathrm{sin}^2\varphi D_{b\mu }U(\mathrm{cos}^2\varphi D_\mu +\mathrm{sin}^2\varphi D_{b\mu })U^+,$$
$`(8.1)`$
where $`\varphi `$ is constant. If $`D_\mu `$ in eq(2.8) is replaced by $`\mathrm{cos}^2\varphi D_\mu +\mathrm{sin}^2\varphi D_{b\mu }`$, we can obtain the following Lagrangian:
$$\begin{array}{ccc}& =& \overline{\psi }[\gamma ^\mu (\mathrm{cos}^2\varphi D_\mu +\mathrm{sin}^2\varphi D_{b\mu })+m]\psi \hfill \\ & & \frac{1}{4K}Tr(A^{\mu \nu }A_{\mu \nu })\frac{1}{4K}Tr(B^{\mu \nu }B_{\mu \nu })\hfill \\ & & \frac{\mu ^2}{2K(1+\alpha ^2)}Tr\left[(A^\mu +\alpha B^\mu )(A_\mu +\alpha B_\mu )\right]\hfill \end{array}$$
$`(8.2)`$
Obviously, this Lagrangian has local $`SU(N)`$ gauge symmetry. Let $`_\psi `$ denote the part of the Lagrangian for fermions:
$$_\psi =\overline{\psi }[\gamma ^\mu (\mathrm{cos}^2\varphi D_\mu +\mathrm{sin}^2\varphi D_{b\mu })+m]\psi .$$
$`(8.3)`$
Using eqs(2.4a,b), we can change $`_\psi `$ into
$$_\psi =\overline{\psi }[\gamma ^\mu (_\mu ig\mathrm{cos}^2\varphi A_\mu +i\alpha g\mathrm{sin}^2\varphi B_\mu )+m]\psi .$$
$`(8.4)`$
From the above Lagrangian, we know that both gauge fields $`A_\mu `$ and $`B_\mu `$ directly couple to matter field $`\psi `$. Substitute eqs(5.7a,b) into eq(8.4), we get
$$_\psi =\overline{\psi }[\gamma ^\mu (_\mu ig\frac{\mathrm{cos}^2\theta \mathrm{sin}^2\varphi }{\mathrm{cos}\theta }C_\mu +ig\mathrm{sin}\theta F_\mu )+m]\psi .$$
$`(8.5)`$
The equation of motion for fermion field $`\psi `$ is
$$[\gamma ^\mu (_\mu ig\frac{\mathrm{cos}^2\theta \mathrm{sin}^2\varphi }{\mathrm{cos}\theta }C_\mu +ig\mathrm{sin}\theta F_\mu )+m]\psi =0.$$
$`(8.6)`$
The equations of motion for gauge fields $`A_\mu `$ and $`B_\mu `$ now change into:
$$D^\mu A_{\mu \nu }\frac{\mu ^2}{1+\alpha ^2}(A_\nu +\alpha B_\nu )=ig\mathrm{cos}^2\varphi \overline{\psi }\gamma _\nu T^i\psi T^i$$
$`(8.7a)`$
$$D_b^\mu B_{\mu \nu }\frac{\alpha \mu ^2}{1+\alpha ^2}(A_\nu +\alpha B_\nu )=i\alpha g\mathrm{sin}^2\varphi \overline{\psi }\gamma _\nu T^i\psi T^i.$$
$`(8.7b)`$
If $`\varphi `$ vanish, the Lagrangian density (8.2) will return to the original Lagrangian density (2.8), the equations of motion (8.7a,b) will return to eqs(6.2a,b), and eq(8.6) will return to eq(6.1). So, the model discussed in the above chapters is just a special case of the model we discuss in this chapter.
## 9 U(1) Case
If the symmetry of the model is U(1) group, we will obtain a U(1) gauge field model. We also use $`A_\mu `$ and $`B_\mu `$ to denote gauge fields and $`\psi `$ to denote a fermion fields. In U(1) case, the strengths of gauge fields are
$$A_{\mu \nu }=_\mu A_\nu _\nu A_\mu $$
$`(9.1a)`$
$$B_{\mu \nu }=_\mu B_\nu _\nu B_\mu $$
$`(9.1b)`$
Two gauge covariant derivatives are the same as (2.4a,b) but with different content. The Lagrangian density of the model is:
$$\begin{array}{ccc}& =& \overline{\psi }[\gamma ^\mu (\mathrm{cos}^2\varphi D_\mu +\mathrm{sin}^2\varphi D_{b\mu })+m]\psi \hfill \\ & & \frac{1}{4}A^{\mu \nu }A_{\mu \nu }\frac{1}{4}B^{\mu \nu }B_{\mu \nu }\hfill \\ & & \frac{\mu ^2}{2(1+\alpha ^2)}(A^\mu +\alpha B^\mu )(A_\mu +\alpha B_\mu )\hfill \end{array}$$
$`(9.2)`$
Local U(1) gauge transformations are
$$\psi e^{i\theta }\psi ,$$
$`(9.3a)`$
$$A_\mu A_\mu \frac{1}{g}_\mu \theta $$
$`(9.3b)`$
$$B_\mu B_\mu +\frac{1}{\alpha g}_\mu \theta .$$
$`(9.3c)`$
Then, $`A_{\mu \nu },B_{\mu \nu }`$ and $`A_\mu +\alpha B_\mu `$ are all U(1) gauge invariant. That is
$$A_{\mu \nu }A_{\mu \nu }$$
$`(9.4a)`$
$$B_{\mu \nu }B_{\mu \nu }$$
$`(9.4b)`$
$$A_\mu +\alpha B_\mu A_\mu +\alpha B_\mu $$
$`(9.4c)`$
Using all these results, it is easy to prove that the Lagrangian density given by eq(9.2) has local $`U(1)`$ gauge symmetry.
Substitute eqs(5.7a,b) into eq(9.2), the Lagrangian density $``$ changes into
$$\begin{array}{ccc}& =& \overline{\psi }[\gamma ^\mu (_\mu ig\frac{\mathrm{cos}^2\theta \mathrm{sin}^2\varphi }{\mathrm{cos}\theta }C_\mu +ig\mathrm{sin}\theta F_\mu )+m]\psi \hfill \\ & & \frac{1}{4}C^{\mu \nu }C_{\mu \nu }\frac{1}{4}F^{\mu \nu }F_{\mu \nu }\frac{\mu ^2}{2}C^\mu C_\mu \hfill \end{array}$$
$`(9.5)`$
where,
$$C_{\mu \nu }=_\mu C_\nu _\nu C_\mu $$
$`(9.6a)`$
$$F_{\mu \nu }=_\mu F_\nu _\nu F_\mu $$
$`(9.6b)`$
So, in this model, there is a massive Abel gauge field as well as a massless Abel gauge field. They all have gauge interactions with matter field. In this case, $`U(1)`$ gauge interactions is transmitted by two different kinds of gauge fields.
It is known that QED is a U(1) gauge field theory. According to the model we discuss here, we may guess that there may exist two different kinds of photon, one is massive while another is massless. And if $`\theta `$ is near $`\pi /2`$ and the mass of massive photon is large, the massless photon field couples with charged fields will be much stronger than massive photon.
## 10 Two limits of the model.
Now, we let’s discuss two kinds of limits of this model. The first kind of limits corresponds to very small parameter $`\alpha `$. Let
$$\alpha 0,$$
$`(10.1)`$
then
$$\mathrm{cos}\theta 1,\mathrm{sin}\theta 0.$$
$`(10.2)`$
From eqs(5.6a,b), we know that the gauge field $`A_\mu `$ is just gauge field $`C_\mu `$ and the gauge field $`B_\mu `$ is just the gauge field $`F_\mu `$. That is
$$C_\mu A_\mu ,F_\mu B_\mu .$$
$`(10.3)`$
In this case, the Lagrangian density (2.9) becomes
$$\begin{array}{ccc}& & \overline{\psi }[\gamma ^\mu (_\mu igC_\mu ^iT^i)+m]\psi \hfill \\ & & \frac{1}{4}C^{i\mu \nu }C_{\mu \nu }^i\frac{1}{4}F^{i\mu \nu }F_{\mu \nu }^i\frac{\mu ^2}{2}C^{i\mu }C_\mu ^i.\hfill \end{array}$$
$`(10.4)`$
The massless gauge field do not interact with matter fields in this limit. So, the $`\alpha 0`$ limit corresponds to the case that gauge interactions are mainly transmitted by massive gauge field. So, the above Lagrangian approximately describes those kinds of gauge interactions which are dominated by massive gauge bosons.
The second kind of limits corresponds to very big parameter $`\alpha `$. Let
$$\alpha \mathrm{},$$
$`(10.5)`$
then
$$\mathrm{cos}\theta 0,\mathrm{sin}\theta 1.$$
$`(10.6)`$
From eqs(5.6a,b), we know that:
$$C_\mu B_\mu ,F_\mu A_\mu .$$
$`(10.7)`$
Then, the Lagrangian density (2.9) becomes
$$\begin{array}{ccc}& & \overline{\psi }[\gamma ^\mu (_\mu +igF_\mu ^iT^i)+m]\psi \hfill \\ & & \frac{1}{4}F^{i\mu \nu }F_{\mu \nu }^i\frac{1}{4}C^{i\mu \nu }C_{\mu \nu }^i\frac{\mu ^2}{2}C^{i\mu }C_\mu ^i.\hfill \end{array}$$
$`(10.8)`$
In this case, massive gauge field does not directly interact with matter fields. So, this limit corresponds to the case when gauge interactions are mainly transmitted by massless gauge field.
In the particles’ interaction model which describes the gauge interactions of real world, the parameter $`\alpha `$ is finite,
$$0<\alpha <\mathrm{}.$$
$`(10.9)`$
In this case, both massive gauge field and massless gauge field directly interact with matter fields, and gauge interactions are transmitted by both of them.
## 11 The renormalizability of the theory
The renormalizability of the theory can be very strictly proved. But this proof is extremely long and is not suitable to write it here. We will not discuss it in details in this paper. We only want to discuss some key problems on the renormalizability of the theory.
It is know that, according to the power counting law, a massive vector field model is not renormalizable in most case. The reason is simple. It is known that the propagator of a massive vector field usually has the following form:
$$\mathrm{\Delta }_{F\mu \nu }=\frac{i}{k^2+\mu ^2i\epsilon }(g_{\mu \nu }+\frac{k_\mu k_\nu }{\mu ^2}).$$
$`(11.1)`$
So, when we let
$$k\mathrm{}$$
$`(11.2)`$
then,
$$\mathrm{\Delta }_{F\mu \nu }const.$$
$`(11.3)`$
In this case, there are infinite kinds of divergent Feynman diagrams. According to the power counting law, this theory is a kind of non-renormalizable theory. Though gauge field theory contains massive vector fields, it is renormalizable. The key reason is that the Lagrangian has local gauge symmetry\[8 \].
As we have stated before, this gauge field theory has maximum local $`SU(N)`$ gauge symmetry. When we quantize this gauge field theory in the path integral formulation, we must select gauge conditions first \[17 \]. In order to make gauge transformation degree of freedom completely fixed, we must select two gauge conditions simultaneously: one is for massive gauge field $`C_\mu `$ and another is for massless gauge field $`F_\mu `$. For example, if we select temporal gauge condition for massless gauge field $`F_\mu `$:
$$F_4=0,$$
$`(11.4)`$
there still exists remainder gauge transformation degree of freedom, because temporal gauge condition is unchanged under the following local gauge transformation:
$$F_\mu UF_\mu U^{}+\frac{1}{ig\mathrm{sin}\theta }U_\mu U^{}$$
$`(11.5)`$
where
$$_tU=0,U=U(\stackrel{}{x}).$$
$`(11.6)`$
In order to make this remainder gauge transformation degree of freedom completely fixed, we’d better select another gauge condition for gauge field $`C_\mu `$. For example, we can select the following gauge condition for gauge field $`C_\mu `$:
$$^\mu C_\mu =0.$$
$`(11.7)`$
If we select two gauge conditions simultaneously, when we quantize the theory in path integral formulation, there will be two gauge fixing terms in the effective Lagrangian. The effective Lagrangian can be written as:
$$_{eff}=\frac{1}{2\alpha _1}f_1^af_1^a\frac{1}{2\alpha _2}f_2^af_2^a+\overline{\eta }_1M_{f1}\eta _1+\overline{\eta }_2M_{f2}\eta _2$$
$`(11.8)`$
where
$$f_1^a=f_1^a(F_\mu ),f_2^a=f_2^a(C_\mu )$$
$`(11.9)`$
If we select
$$f_2^a=^\mu C_\mu ^a,$$
$`(11.10)`$
then the propagator for massive gauge field $`C_\mu `$ is:
$$\mathrm{\Delta }_{F\mu \nu }^{ab}(k)=\frac{i\delta ^{ab}}{k^2+\mu ^2i\epsilon }\left(g_{\mu \nu }(1\frac{1}{\alpha _2})\frac{k_\mu k_\nu }{k^2\mu ^2/\alpha _2}\right).$$
$`(11.11)`$
If we let $`k`$ approach infinity, then
$$\mathrm{\Delta }_{F\mu \nu }^{ab}(k)\frac{1}{k^2}.$$
$`(11.12)`$
In this case, according to the power counting law, the gauge field theory which is discussed in this paper is a kind of renormalizable theory. At the same time, the local $`SU(N)`$ gauge symmetry will give a Ward-Takahashi identity which will eventually make the theory renormalizable. A strict proof on the renormalizability of the gauge field theory can be found in the reference \[15 \].
In order to make the gauge field theory renormalizable, it is very important to keep the maximum local $`SU(N)`$ gauge symmetry of the Lagrangian. From the above discussions, we know that, in the renormalization of the gauge field theory, local gauge symmetry plays the following two important roles: 1) to make the propagator of the massive gauge bosons have the renormalizable form; 2) to give a Ward-Takahashi identity which plays a key role in the proof of the renormalizability of the gauge field theory.
## 12 Comments
Up to now, we know that there are two mechanisms that can make gauge field obtain non-zero mass: one is the Higgs mechanism which is well known in constructing the Standard Model; another is the mechanism discussed in this paper. In this new mechanism, the mass term of gauge field is introduced by using another set of gauge field and the mass term of gauge fields does not affect the symmetry of the Lagrangian. We can imagine the new interaction picture as: when matter fields take part in gauge interactions, they emit or absorb one kind of gauge field which is not eigenstate of mass matrix, when we detect this gauge fields in experiments, it will appear in two states which corresponds to two kinds of vector fields, one is massless and another is massive.
Though the Lagrangian of the model contains the mass term of gauge fields, the theory is renormalizable. So, we can use the mechanism to describe gauge interactions of quarks and leptons. If we apply this mechanism to electroweak interactions, we can construct an electroweak model which contains no Higgs particle.
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# Large Scale-Invariant Fluctuations in Normal Blood Cell Counts: A sign of criticality?
## Abstract
All types of blood cells are formed by differentiation from a small self-maintaining population of pluri-potential stem cells in the bone marrow. Despite abundant information on the molecular aspects of division, differentiation, commitment and maturation of these cells, comparatively little is known about the dynamics of the system as a whole, and how it works to maintain this complex “ecology” in the observed normal ranges throughout life. Here we report unexpected large, scale-free, fluctuations detected from the first long-term analysis of the day-to-day variability of a healthy animal’s blood cell counts measured over one thousand days. This scale-invariance cannot be accounted for by current theoretical models, and resembles some of the scenarios described for self-organized criticality.
Running Title: Scale-Invariance in blood cells counts
Address all correspondence to:
Dr. Peter Willshaw, Universidad Favaloro. Solís 453, (1078) Buenos Aires - Argentina
Email: peterw@favaloro.edu.ar Phone: +54-11-43849920 Fax: +54-11-43839732
The existence of scale-invariant fluctuations in nature is well-documented. Scale-invariance refers to the absence of a single time or length scale characteristic of the (temporal or spatial) pattern of change. In such cases fluctuations happen at all scales, unlike the more familiar cases in which there is a simple statistical spread around a defined mean value. Scale-invariance has been shown for fluctuations in heart rate, gait dynamics, DNA nucleotide sequences , and lung inflation. How these fluctuations are generated by regulatory physiological mechanisms and what their relevance is in biology is the focus of much current interest. In addition, it has been argued that scale-invariance in itself can be the fingerprint of a biological system poised at a critical state, resulting in improved dynamics with robust tolerance to error, and at the same time remaining very susceptible to change. Roughly speaking, this idea of self-organized criticality as a homeostatic mechanism is that systems with many interacting non-linear parts could find themselves moving toward the critical state (i.e. “self-organizing”) simply because it would be absolutely impractical to administer it in any centralized way. Here we show that the regulation of the number of the formed elements in blood also has scale-invariant properties. We found that the day-to-day variability in the number of platelets (PLA), red (RBC) and white blood cells (WBC) from two healthy sheep observed for one thousand days is scale-invariant over a range spanning from a few to over two hundred days.
Even in healthy subjects the number of blood cells fluctuates from one day to the next. Health practitioners are very aware of this fact, to the point that these quantities are usually described as ranges or maximal limits and not as means and standard deviations. In other words, for these variables normality could be defined in one sense as “fluctuations within a normal range”. Figure 1 shows the daily values (over 1000 days) for blood cells taken from two healthy sheep. Overall, the values observed are consistent with the normal range reported in the literature for these animals. Simple inspection of the data reveals fast, medium and slow fluctuations, and a certain degree of correlation in the behaviour of the time series between the animals, as if they were responding in a similar way to some common external factor. What is the nature of the extensive fluctuations seen in all the series plotted in Figure 1? According to current understanding, homeostatic mechanisms act to compensate the loss of a given number of cells by speeding up the process of maturation of precursors, i.e., there is a negative feedback loop attempting to maintain constant the number of blood cells. The precise characterization of the observed fluctuations is important in assessing the adequacy of current models, because much is known about the dynamics which regulatory mechanisms of this kind are capable of exhibiting. Thus one could reject or accept classes of models, or propose the need for novel ones, depending on the type of fluctuations found experimentally.
The dynamics of any underlying control process is expected to be reflected in the temporal correlation between the data points; this is the physiologically important information we are interested in. We consider here two extreme (null) hypotheses: the first is to assume a tight feedback mechanism perturbed by some jitter where one should simply see a mean value contaminated with white noise. The second corresponds to the very unlikely case that no control is exerted to compensate for blood cell losses, resulting in a random walk motion, in which the fluctuation is just generated by the day-by-day summation of statistically independent increments or decrements. To compare our observed time series against these two possibilities, surrogate data sets were constructed to mimic the null hypotheses cited above.
In Figure 2, the raw time series of red blood cells of one animal is plotted in the top panel. Panels denoted S1 and S2 depict the surrogate series of the first and second types respectively. The first one (S1) is constructed by re-ordering at random the raw values of the top panel, like shuffling a deck of cards. In this way, any temporal relationship between the counts on one day and the next in the original data is broken, and the resulting time series is white noise. Surrogate S2 is built by (1) generating an accessory file with the day-to-day differences of the raw data in the top panel; (2) randomizing the order of the data points from this accessory file, and (3) integrating the result over the whole length of the file to synthesize a new time series. Surrogate type 2 is therefore a random walk-like signal. By using different random orderings, one can build many examples or realizations of these types of surrogates.
To quantify the characteristics of the temporal correlations in the data we use detrended fluctuation analysis (DFA). This method has been used successfully to analyze biological time series, circumventing some technical difficulties, such as trends and non-stationarities. The DFA algorithm involves the following steps: 1) Denoting $`C_j`$ as the blood cell daily count of interest on the $`j_{th}`$ day, produce a new time-integrated series $`y_i=_{j=1}^i(C_jC_{av})`$, where $`C_{av}`$ is the average of all $`C_j`$. 2) The new series $`y_i`$ is subsequently divided into boxes of equal length $`n`$, and in each box a least-squares line is fitted to the data, representing the trend over the chosen time interval. 3) Define and calculate $`F(n)`$ as the square root of the mean of the squares of the residuals in all boxes of length n for a given value of n. Thus, $`F(n)`$ quantifies the magnitude of the fluctuation of the integrated and detrended time series over increasingly long time intervals. The presence of a straight line of slope $`\alpha `$ in the log-log plot of F(n) vs. n implies a relationship of the type $`F(n)n^\alpha `$ . The range of $`n`$ over which this relationship holds defines the time span over which the fluctuations are scale-invariant. The two extreme cases considered have expected values of $`\alpha =0.5`$ for white noise, and 1.5 for a random-walk motion.
Panel a in Fig. 3 is a log-log plot of F(n) vs. n resulting from the analysis of the red cells data set plotted in Fig. 2. Symbols on the right of the figure indicate to which time series the calculated function belongs, R denoting the raw time series, and S1 and S2 the 30 random realizations of surrogates 1 and 2 respectively. The plot for the raw data is approximately linear over time lags n from 1-200 days, indicating that blood cell count fluctuations are scale-invariant. We found similar scaling behaviour in all series in both sheep. The lines of best fit gave slopes $`\alpha =0.98`$ and 1.00 for RBC, $`\alpha =1.24`$ and 1.11 for PLA, and $`\alpha =0.83`$ and 0.83 for WBC (sheep 77 and 78 respectively). The lowest value of the correlation coefficient r for all raw data series was 0.995. Panel b in Figure 3 illustrates an important consequence of scale-invariance: on increasing the observation time by a factor of $`k`$ the fluctuations are found to be $`k^\alpha `$ larger, that is to say the longer we look the bigger the range of values we will find.
Additional insight can be gained by applying standard spectral analysis techniques to the same series. In this case, it is appropriate to consider the spectral analysis of the daily differences (i.e, $`C_jC_{j1})`$ of the time series, since this is of a stationary nature by construction. We computed the power spectra $`S(f)`$ of each differentiated time series and in all cases we found that the square of their amplitudes scales as $`S(f)f^\beta `$ . The scaling factors $`\alpha `$ and $`\beta `$ are related as $`\beta =32\alpha `$. The values obtained via spectral analysis (for sheep 77 and 78) were $`\beta =0.97`$ and 0.91 for RBC, (the values predicted from $`\alpha `$ are 1.04 and 1.00 respectively); for PLA $`\beta =0.65`$ and 0.70 (predicted values of 0.52 and 0.78 respectively); and for WBC $`\beta =1.14`$ and 1.16 (1.34 and 1.34 respectively). There is clear agreement between the scaling exponents obtained using both methods.
The value of $`\beta `$ provides additional information regarding the nature of the process, because $`0<\beta <1`$ implies that a high value is more likely to be follow by a small value and viceversa, i.e., there is anti-correlation between values collected on successive days. The fact that we found a wide range (1- 200 days) over which this scaling factor holds, indicate that the mechanism responsible for the anti- correlated dynamics operates at all time scales, something that a regulatory system as a simple negative feedback-like loop is incapable of doing.
We must therefore reject the two null hypotheses considered. The fact that the data is not white is surprising, since this is the most likely possibility if there is simple feedback control such as that postulated to relate erythropoietin production and erythrocyte number. The fact that the data is not a random walk does at least suggest that control is present, and that it is non-trivial. Substantial work exists regarding the architecture of haematopoietic regulation, including the seminal work of Mackey but interestingly enough no model has as yet reproduced the scale-invariance reported here. We suspect that an important part of the fluctuations we have observed reflects genuine intrinsic dynamics of the system, which possesses a very large number of parts. To account for these observations, models which preserve this very large number of degrees of freedom are needed. The behaviour here reported for sheep has also been suggested to be present in human neutrophil counts, albeit using a far shorter time series than ours. To show that this type of scale invariance is indeed present in humans without going to the extreme of daily sampling, it could be possible to collect data from many individuals who have been subject to blood sampling and counting on at least two occasions a known number of days, months or years apart and estimating a similar fluctuation function. This could be achieved for example in records of regular blood donors, and will estimate the likelihood of encountering a given difference as a function of time lag. At present no information exists which health practitioners can use to forecast such simple physiological parameters in situations of interest such as bone marrow transplant.
If our results hold in humans, one important consequence of scale invariance is to make it difficult to differentiate between the response to a given treatment and a chance value due to the intrinsic (scale-invariant) nature of this type of fluctuations, especially if the patient is seen at long intervals. Health practitioners generally make clinical decisions on the basis of isolated data points, it being implicit that waiting longer or taking more samples can only reduce the variance. Our findings imply exactly the opposite; waiting longer increases the variance!
Acknowledgements: We are indebted to veterinarians Drs. M. Besancon, P. Iguain, and M. Tealdo for the collection and processing of the blood samples.
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# Dynamical Properties of Spin-Orbital Chains in a Magnetic Field
## I Introduction
There is a wide variety of quasi-one-dimensional spin systems with short-ranged and quasi-long-ranged quantum liquid states. Examples for short-ranged compounds include the spin-Peierls chain, even-leg spin-1/2 Heisenberg ladders, and integer-spin Heisenberg chains. In contrast, half-odd-integer-spin chains and odd-leg spin-1/2 Heisenberg ladders are known to have a quantum critical ground state with quasi-long-range-ordered spin correlations. Many quasi-one-dimensional models, such as the antiferromagnetic spin-1/2 Heisenberg chain with frustrating longer-range exchange interactions, have a phase diagram which contains both types of disordered ground state, depending on the parameters. Particular recent attention has focussed on spin-1/2 Heisenberg chains coupled to orbital degrees of freedom. These systems are known to have a rich phase diagram, containing short-ranged and quasi-long-ranged regions, including an SU(4)-symmetric integrable point at a particular combination of the exchange parameters.
In this work, we study the spin and orbital excitation spectra of spin-orbital chains in a magnetic field. By applying a newly developed dynamical density matrix renormalization group (DMRG) technique, it is found that the admixture of the orbital degrees of freedom has a profound influence on both the field-dependent incommensuration, and the dispersion and spectral weight of the spin triplet modes, although only the spins couple directly to the applied magnetic field. Furthermore, this study is intended to provide useful information for inelastic neutron scattering studies on orbitally degenerate materials such as $`\mathrm{Na}_2\mathrm{Ti}_2\mathrm{Sb}_2\mathrm{O}`$, where the d-electron of the Ti<sup>3+</sup> ion occupies an $`e_g`$ doublet, and $`\mathrm{NaV}_2\mathrm{O}_5`$, where orbital degeneracy appears to arise from the coupling of two types of $`\mathrm{VO}_5`$ chains. As long as experimentally available magnetic fields ($``$ 30 Tesla) are sufficiently large to overcome the respective spin gaps, our results for the dynamical properties of the incommensurate phases in the one-dimensional spin-orbital model may be tested on these or suitably related compounds.
In this paper, the Hamiltonian of the one-dimensional spin-orbital model in a magnetic field is studied. This can be written as
$`H=J_1{\displaystyle \underset{i=1}{\overset{N}{}}}𝐒_i𝐒_{i+1}+J_2{\displaystyle \underset{i=1}{\overset{N}{}}}𝐓_i𝐓_{i+1}+K{\displaystyle \underset{i=1}{\overset{N}{}}}\left(𝐒_i𝐒_{i+1}\right)\left(𝐓_i𝐓_{i+1}\right)h{\displaystyle \underset{i=1}{\overset{N}{}}}S_i^z,`$ (1)
where the spin-1/2 operators $`𝐒_i`$ and the pseudo-spin-1/2 operators $`𝐓_i`$ describe the spins and orbits at a site $`i`$. The first two terms are the Hamiltonians for (pseudo-)spin-1/2 Heisenberg chains with exchange constants $`J_1`$ for the spins and $`J_2`$ for the orbitals. The third term couples these two sets of degrees of freedom with a constant $`K`$. The fourth term is the Zeeman coupling of the spins to an external magnetic field aligned with the $`z`$-direction. In the absence of a magnetic field, this Hamiltonian has an SU(4) symmetry for the particular parameter combination $`J_1=J_2=K/4`$. This symmetry reduces to SU(2)$`\times `$U(1) in an intermediate regime ($`0<h<h_c`$), and to SU(2) at very high magnetic fields where the spins are completely aligned. If the parameter set ($`J_1,J_2,K`$) is chosen away from the SU(4) point, the Hamiltonian remains invariant under SU(2)$`\times `$SU(2). Furthermore, there is a $`Z_2`$ symmetry about $`J_1=J_2`$ for vanishing magnetic fields.
Recent studies have concentrated on the phase diagram and low-energy excitations of $`H`$. At the integrable SU(4) point, the spectrum of excitations has been obtained from the Bethe Ansatz equations. Effective low-energy theories were obtained from bosonization, and approximate solutions from mean-field theory. In addition, a variety of numerical techniques, including series expansions, numerical diagonalization, Monte Carlo, and density matrix renormalization group, has been employed to analyze the ground state and low-energy properties of this Hamiltonian.
In this work, a recently developed density matrix renormalization group technique is applied to study the full spin and orbital dynamics of the one-dimensional spin-orbital model in a magnetic field for arbitrary parameters ($`J_1,J_2,K,h`$). The advantages of this method are that (i) relatively large Hilbert spaces can be studied (typical system size $``$ 80 sites), (ii) the full excitation spectra can be obtained, and (iii) there are no constraints concerning low energies, small deviations from the SU(4) point, or the magnitude of the applied magnetic field. However, the spectra calculated with this method do exhibit some finite-size effects, as expected for low-dimensional quantum systems. Here we concentrate on the characteristic features of the dynamical spin and orbital structure factors for the various regimes of $`H`$. These quantities are defined as
$`S(k,\omega )={\displaystyle \frac{1}{\pi }}\underset{ϵ0}{lim}\mathrm{Im}\mathrm{\Psi }_0|S_k^z{\displaystyle \frac{1}{\omega +iϵH+E_0}}S_k^z|\mathrm{\Psi }_0,`$ (2)
$`T(k,\omega )={\displaystyle \frac{1}{\pi }}\underset{ϵ0}{lim}\mathrm{Im}\mathrm{\Psi }_0|T_k^z{\displaystyle \frac{1}{\omega +iϵH+E_0}}T_k^z|\mathrm{\Psi }_0,`$ (3)
where $`|\mathrm{\Psi }_0`$ is the ground state, $`E_0`$ its energy, $`S_k^z`$ and $`T_k^z`$ are the Fourier transforms of the local spin and orbit operators $`S_i^z`$ and $`T_i^z`$, and a small but finite broadening factor $`ϵ=0.05J_1`$ is chosen to represent the poles in $`S_k^z`$ and $`T_k^z`$.
In general, for the properties of the orbitally degenerate Hamiltonian $`H`$ to be applicable to candidate inorganic materials such as $`\mathrm{Na}_2\mathrm{Ti}_2\mathrm{Sb}_2\mathrm{O}`$ and $`\mathrm{NaV}_2\mathrm{O}_5`$, the orbital energy levels must be nearly degenerate on the scale of the spin dynamics, which have characteristic energies in the range of $`J_1`$ 10-100 meV. While external magnetic fields couple directly to the spin degrees of freedom, applied external pressure could be a “dual” perturbation, primarily affecting the orbitals. It would therefore be of great interest to explore the dynamical inelastic neutron scattering response in candidate spin-orbital compounds subject to a magnetic field as well as to external pressure.
In the following section, we discuss briefly the numerical renormalization group method which is used to obtain the spin dynamics of the spin-orbital chain in a magnetic field. In section III, the field-dependent spectra are studied in the different parameter regimes ($`J_1,J_2,K,h`$). A summary and conclusions are given in section IV.
## II Calculation of Dynamical Properties with the Density Matrix Renormalization Group Technique
The density matrix renormalization group (DMRG) method is known to be a precise numerical technique to calculate the ground state properties of quasi-one-dimensional systems such as chains and ladders. Recently, Hallberg and Kühner and White have proposed extensions of this method to extract dynamical correlation functions, based on continued-fraction expansions of the dynamical Green function. In this paper, we use a variant of this technique to calculate the momentum and frequency-dependent spectrum of the spin-orbital model in one spatial dimension.
The zero-temperature, time-dependent correlation function of a momentum-dependent operator $`\widehat{A}_k`$ is defined as
$`A(k,tt^{})\mathrm{\Psi }_0|\widehat{A}_k^{}(t)\widehat{A}_k(t^{})|\mathrm{\Psi }_0,`$ (4)
where $`|\mathrm{\Psi }_0`$ is the ground state. In the following, $`\widehat{A}_k`$ corresponds either to the spin operator $`S_k^z`$ or to the orbital operator $`T_k^z`$ . The Fourier transform of this correlation function can be expanded in the form of a continued fraction,
$`A(k,\omega )={\displaystyle \frac{1}{\pi }}\mathrm{Im}{\displaystyle \frac{\mathrm{\Psi }_0|\widehat{A}_k^{}\widehat{A}_k|\mathrm{\Psi }_0}{\omega +iϵa_0\frac{b_1^2}{\omega +iϵa_1\mathrm{}}}},`$ (5)
where $`ϵ0`$ is a small broadening factor. The coefficients $`a_n`$ and $`b_n`$ are then obtained from the recursion relation
$`|f_{n+1}=H|f_na_n|f_nb_n^2|f_{n1},`$ (6)
with $`|f_0\widehat{A}_k|\mathrm{\Psi }_0`$ and $`b_00`$. This procedure is very similar to numerical diagonalization calculations. However, there are specific technical problems intrinsic to the DMRG method which need to be addressed.
One complication arises when open boundary conditions (OBC) are used, which is usual in DMRG calculations to yield the highest possible numerical precision. In this case, the momentum is not a well-defined quantum number, and each momentum-dependent operator $`\widehat{A}_k`$ generates wave packets of finite width, centered around $`k`$. To remedy this problem, Hallberg has used periodic boundary conditions (PBC) to obtain a basis with a good momentum quantum number, whereas Kühner and White have used OBCs and the Parzen filter to eliminate the artificial surface excitations on the surface boundaries. In this paper, we adapt OBCs combined with the Parzen filter to obtain the full dynamical spectra associated with $`S_k^z`$ and $`T_k^z`$. OBCs are known to give more precise ground states than PBCs, and therefore larger clusters can be studied, while the width of the corresponding wave packets is greatly reduced.
Further improvement of the dynamical DMRG method may be obtained by using more target states besides the ground state, implemented with “correction vectors”. We find that a rather large cut-off dimension of the density matrix is needed to reach convergence in the spin-orbital chain for the target states at higher energy. These convergence problems tend to be exacerbated for excitations away from the dominant wave vectors. However, with the algorithm used in this work, the full structure of the excitation spectrum can be obtained with rather high accuracy, even if only the ground state is kept as a target state, although the relative overlap (matrix elements) with higher-energy states is reduced.
We have tested the algorithm on the spin-1/2 antiferromagnetic Heisenberg chain in a magnetic field. The Hamiltonian for this system corresponds to the parameter choice ($`J_1=J,J_2=K=0`$) for the spin-orbital model. One observes that the spin-triplet excitation spectrum of this simple model matches very well the exact results from the Bethe Ansatz. A more detailed discussion of this spectrum will be presented in the next section. Chains with $`N=80`$ spins were studied. The cut-off factor, i.e. the number of states kept per block, was taken to be $`m=160`$. We have also used the Lanczos method on lattices of up to 16 sites to check the accuracy of the DMRG results for other non-trivial parameter sets of the spin-orbital chain.
## III Dynamical Spin Structure Factor of the Spin-Orbital Chain
The rich phase diagram of the one-dimensional spin-orbital model, represented by the Hamiltonian $`H`$ in Eq. 1, has been the subject of a number of recent studies. The main result of these works is summarized in Fig. 1, which we have adapted from Ref.. The various phases can be characterized by spin and orbital quantum numbers $`(S,S^z)`$ and $`(T,T^z)`$, and by effective Luttinger parameters in each channel, such as gaps $`(\mathrm{\Delta }_S,\mathrm{\Delta }_O)`$, soliton velocities $`(v_S,v_O)`$, and radii of compactification. For finite values of the spin-orbital interaction $`K`$, the spin and orbital degrees of freedom are coupled, and the excitation spectra generally contain strongly admixed spin and orbital modes.
Let us first briefly describe the phase diagram of $`H`$. Phase I is “ferromagnetic” in the spin and orbital channels, i.e. both degrees of freedom are fully polarized, and there is an infinite degeneracy in the ground state. In phase II, the spin degrees of freedom remain polarized, whereas the orbitals are in a quasi-long-range-ordered antiferromagnetic state with a characteristic orbiton velocity $`v_O`$. In phase III, the roles of spins and orbitals are interchanged. Phase IV is fully gapped in all channels ($`\mathrm{\Delta }_S0,\mathrm{\Delta }_O0`$) because of spontaneous dimerization. In particular, this phase contains a Majumdar-Ghosh point at $`J_1=J_2=3K/4`$ which has a matrix-product ground state. It has recently been shown that phase V contains 3 gapless excitations with non-universal, parameter-dependent soliton velocities. At the SU(4)-symmetric point (filled circle in Fig. 1), these three modes have one common velocity, $`v=\pi J_1/2=\pi J_2/2=\pi K/8`$.
Let us now turn to the excitation spectra of $`H`$ in these regimes, starting with phase III where the orbital degrees of freedom are fully polarized. In Fig. 2, the spectrum traced by the dynamical spin structure factor is shown on a lattice of 80 sites, using the parameter set ($`J_1=J_2=K/2`$). The poles have been given a width of $`ϵ=0.05J_1`$, and the color scheme highlights the areas of high spectral intensity with bright colors. Since in this phase, the orbital degrees of freedom couple only trivially with the spins, the analysis turns out to be particularly simple. The spin excitation spectrum of this system corresponds to a spin-1/2, antiferromagnetic Heisenberg chain with exchange coupling $`J_1`$, which is dominated by a low-energy 2-spinon continuum (see Fig. 2(a)). The elementary excitations of this system are spin-1/2 solitons with dispersion $`\omega _\sigma (q)=3\pi J_1\mathrm{sin}(q)/4`$ for $`q[0,\pi ]`$. Thus the 2-spinon continuum has a lower cut-off at $`\omega _l(k)=3\pi J_1|\mathrm{sin}(k)|/4`$ and an upper cut-off at $`\omega _u(k)=3\pi J_1|\mathrm{sin}(k/2)|/2`$. The dominant singularity of the dynamical spin structure factor at wave vector $`\pi `$ diverges according to $`S(\pi ,\omega )\sqrt{\mathrm{ln}(\omega )}/\omega `$ as $`\omega 0`$.
In a finite magnetic field, this mode becomes incommensurate with nesting vectors $`2k_F=\pi (1\pm 2m)`$, where $`m[0,1/2]`$ is the magnetization. This behavior is confirmed by our dynamical DMRG calculation, as shown in Fig. 2(b). Here, a finite magnetization, $`m=4/40`$, leads to an incommensuration with $`2k_F=12\pi /10`$ and $`8\pi /10`$. Because of the relatively large system size ($`N=80`$), the spectra appear to be almost continuous, as expected in the thermodynamic limit. However, finite-size effects are still visible, causing gaps and uneven spectral weight distributions at higher energies. Furthermore, the momentum is quantized in 80 slices of size $`\pi /40`$. In spite of these limitations, the agreement of the numerical results with the exactly known spectra is encouraging, and gives confidence to proceed to parameter regimes with more complex excitation spectra.
At the SU(4) point ($`J_1=J_2=K/4`$), the spectrum is also known exactly from the Bethe Ansatz. A comparison of the spectra obtained using the DMRG method with the exact results can thus help to further evaluate the accuracy of the numerical approach. In Fig. 3, the spectrum traced by the dynamical spin structure factor is shown on a lattice of 80 sites. From the solution of the Bethe Ansatz equations, it is known that the elementary excitations of the SU(4) chain are 4-fold degenerate spin-1/2 spinons ($`\sigma `$) and pseudo-spin-1/2 orbitons ($`\tau `$), and 6-fold degenerate spin-1 or orbital-1 solitons ($`\nu `$).. The corresponding dispersions of these three elementary excitations are
$`\omega _\sigma (q)`$ $`=`$ $`{\displaystyle \frac{J\pi }{2}}\left[\sqrt{2}\mathrm{cos}(q+3\pi /4)+1\right],q[0,3\pi /2]`$ (7)
$`\omega _\tau (q)`$ $`=`$ $`{\displaystyle \frac{J\pi }{2}}\left[\sqrt{2}\mathrm{cos}(q+\pi /4)1\right],q[0,\pi /2]`$ (8)
$`\omega _\nu (q)`$ $`=`$ $`{\displaystyle \frac{J\pi }{2}}\mathrm{cos}(q+\pi /2),q[0,\pi ]`$ (9)
where $`JJ_1=J_2=K/4`$. The spin-triplet excitation spectrum of an SU(4) chain with $`N=4n`$ sites, shown in Fig. 3(a), contains convolutions of elementary $`\sigma `$ and $`\tau `$ excitations, as well as $`\nu `$ excitations. The $`z`$-component of the dynamical spin structure factor $`S^z(k,\omega )`$ couples directly to the 15-fold multiplet of $`\sigma \tau `$ pairs, the 20-fold multiplet of $`2\nu `$ excitations, and with much smaller matrix elements to the 45-fold $`2\sigma \nu `$ and the 35-fold $`4\sigma `$ multiplets.
Fig. 3(b) illustrates that the spectrum obtained with the dynamical DMRG traces exactly the two-soliton continua, spanned by $`\omega _\sigma (q)+\omega _\tau (q^{})`$ and by $`\omega _\nu (q)+\omega _\nu (q^{})`$. These continua have characteristic low-frequency power-law singularities ($`S^z(k,\omega )\omega ^\alpha `$) at the Fermi vectors $`\pi /2`$ and $`3\pi /2`$, and a weaker divergence at $`\pi `$. The relatively high spectral intensity seen at larger frequencies and centered around the wave vector $`\pi `$, can be attributed to a particle-hole excitation in the $`\nu `$-channel which will be discussed later. We note that the accuracy of the dynamical DMRG method is known to deteriorate at larger energies, and hence the spectral intensities are less reliable in this regime.
When an external magnetic field is applied along the $`z`$-direction, the spectrum of the SU(4) chain also becomes incommensurate. This effect can be understood quantitatively by considering the splitting of spin-up and spin-down bands due to the applied field. The corresponding nesting vectors depend on the magnetization $`m[0,1/2]`$ as $`2k_F=\pi (1+2m)/2`$, $`2k_F=\pi (12m)/2`$, and $`4k_F=4k_F=\pi (12m)`$. This low-frequency behavior is reflected in the finite-size DMRG spectra, shown in Fig. 4, with magnetizations $`m=4/40`$ and $`m=12/40`$. For example, with $`m=4/40`$, the nesting vectors are $`6\pi /10`$, $`4\pi /10`$, and $`8\pi /10`$, as shown in Figs. 4(a) and (c), whereas for $`m=12/40`$ they are $`8\pi /10`$, $`2\pi /10`$, and $`4\pi /10`$ respectively (Figs. 4(b) and (d)).
Furthermore, the overall spectral weight is reduced as the phase space of spin-triplet excitations shrinks with increasing field. At finite magnetic fields, there remains a commensurate soft mode at wave vector $`\pi `$ in the orbital $`\nu `$-channel which does not couple to $`S(k,\omega )`$. It is therefore interesting to examine the orbital excitation spectrum traced out by $`T(k,\omega )`$. The orbital dynamical structure factor of the SU(4) chain in a magnetic field is shown in Figs. 4 (c) and (d). In the absence of a magnetic field, the orbital spectrum is identical to the spin spectrum (Fig. 3). In a finite magnetic field there is one orbital mode, the spin-1 $`\nu `$-channel, which does not become incommensurate, but stays soft at wave vector $`\pi `$ for all magnetic fields. On the other hand, the strongly mixed $`\sigma \tau `$ pairs show incommensuration with the same nesting vectors as for the spin channel. The dominant $`\sigma \tau `$ singularities at low frequencies occur at $`2k_F`$ for the orbitons and at $`2k_F`$ for the spinons.
The excitation spectra at the Majumdar-Ghosh point ($`J_1=J_2=3K/4`$) are shown in Fig. 5. The zero-field ground state at this point in the dimerized phase IV is a doubly degenerate product of spin and orbital singlets with an energy of $`3J_1/4`$ per site and a gap of approximately $`3K/8`$ to the lowest excitations. The spectrum contains a 2-soliton continuum of states with minima at wave vectors $`k=0`$ and $`\pi `$, and maximum spectral weight at wave vector $`\pi `$. In addition, there are massive magnon excitations corresponding to Haldane triplet bound states. These are centered at $`k=\pi /2`$ and $`3\pi /2`$, indicating a doubling of the real-space unit cell. In Fig. 5(a), we have eliminated the artificial low-energy states which arise from surface excitations at the open boundaries of the 80-site cluster. These states have vanishing spectral weight in the thermodynamic limit.
A finite magnetic field leads to the deconfinement of bound spinons. This is manifested in a lowering and eventual vanishing of the corresponding spin gap when the field exceeds a critical value, which for the Majumdar-Ghosh point is at $`h_{c1}3K/8`$. In Fig. 5(b), the spin excitation spectrum is shown in this incommensurate regime ($`h_{c1}<h<h_{c2}`$) at a magnetization $`m=4/40`$. The dominant low-energy singularities of the spin-1 magnons occur at wave vectors $`4k_F=\pi (1\pm 2m)`$. The 2-spinon continuum at higher energies appears to be almost unaffected by the finite magnetic field. Interestingly, a strong low-energy singularity at $`k=\pi `$ in the orbital channel is induced by the applied field, as observed in Fig. 5(c), which is indicative of an effective attractive potential between the incommensurate spin-1 and the commensurate orbital-1 $`\nu `$-type solitons in this case. This magnetic-field-induced two-component Luttinger Liquid behavior, with the orbital channel remaining commensurate and the spin channel becoming incommensurate, has also been observed in recent bosonization studies of the spin-orbital chain.
Finally, let us turn to the gapless phase V. In Figs. 6 and 7, the spin and orbital excitation spectra are shown at $`J_1=J_2=K/8`$ and at $`J_1=J_2=0`$. As $`K0`$, a prominent high-energy excitation emerges, centered at wave vector $`\pi `$ and frequency $`\omega 0.8K`$. This excitation carries the quantum numbers $`(S,T)`$ = (0,1) and (1,0) which indicates a particle-hole-type bound state, involving a spin singlet and an orbital triplet or vice versa. From finite-size scaling, we observe that the energy gap between the ground state and this level increases with system size. If $`J_1=J_2=0`$, the eigenstates of the $`K`$ term can be found exactly on a 4-site cluster. The ground state with energy $`3K/4`$ and the excitons with energy 0 are in this case
$`|GS`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{4}{}}}(1)^i|(t_{i,O}^+t_{i+2,O}^{})(s_{i,S}s_{i+2,S})(s_{i,O}s_{i+2,O})(t_{i,S}^+t_{i+2,S}^{})`$ (10)
$`+`$ $`(t_{i+1,O}^{}t_{i+3,O}^+)(t_{i,S}^+t_{i+2,S}^{})(t_{i,O}^+t_{i+2,O}^{})(t_{i+1,S}^{}t_{i+3,S}^+)`$ (11)
$`|\pi ,S`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{4}{}}}|(t_{i,O}^+t_{i+2,O}^{})(t_{i,S}^+s_{i+2,S})\mathrm{and}|\pi ,O={\displaystyle \underset{i=1}{\overset{4}{}}}|(t_{i,S}^+t_{i+2,S}^{})(t_{i,O}^+s_{i+2,O})`$ (12)
where, for example, $`t_{i,S}^+`$ denotes an up-up triplet state in the spin channel, involving spins at sites $`i`$ and $`i+1`$, and $`s_{i+2,O}`$ denotes a singlet in the orbital channel at sites $`i+2`$ and $`i+3`$. This high-energy bound state is observed in the entire phase V, but is most prominent in the strong-$`K`$ limit.
The magnetic-field dependence of the spectra in phase V can be analyzed analogously to the SU(4) limit. In particular, the zero-energy modes can be understood within the band picture discussed above, with nesting instabilities at $`2k_F`$, $`2k_F`$, and $`4k_F`$=$`4k_F`$. The finite-size gaps in these spectra appear to be most obvious at the critical point, and become less severe as $`J_1,J_20`$.
## IV Summary
In summary, the spin and orbital excitation spectra of the spin-orbital chain were studied in various regimes of the model Hamiltonian $`H`$, given in Eq. 1. The characteristic dynamical response depends strongly on the parameter set ($`J_1,J_2,K`$) and on the applied field $`h`$. In phases II and III with either fully polarized spin or orbital degrees of freedom, the remaining gapless mode has the characteristic spectral response of a (pseudo)-spin-1/2 Heisenberg chain. On the other hand, in the dimerized regime (phase IV), a finite magnetic field is needed to close the spin gap. Interestingly, an additional orbital mode is found to become gapless at the commensurate wave vector $`\pi `$ in the partially spin-polarized regime ($`h_{c1}<h<h_{c2}`$), whereas the Fermi momenta of the incommensurate gapless spin modes depend on the magnetization via $`2k_F=\pi (12m)`$. In the gapless regime (phase V), there are three elementary gapless excitations. The dynamical response consists of various overlapping two-particle continua, arising from pairs of these elementary excitations. At higher energies, a particle-hole bound state whose origin lies in the biquadratic $`K`$-term, is found to dominate the spectrum.
We are greatful to Andreas Honecker, Till Kühner, Bruce Normand, and Edmond Orignac for useful discussions, and we thank Till Kühner for valuable help in developing the dynamical DMRG algorithm. S. H. acknowledges the Zumberge foundation for financial support.
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# 1 |𝐪|,𝜈 plane. The thick solid line separates space- and timelike region. The solid line corresponds to elastic e-p scattering, the dashed lines to constant four-momentum transfer.
Many-body theory interpretation of deep inelastic scattering
O. Benhar<sup>1</sup>, V.R. Pandharipande<sup>2</sup>, I. Sick<sup>3</sup>
<sup>1</sup> INFN, Sezione Roma 1, I-00185 Rome, Italy
<sup>2</sup> Department of Physics, University of Illinois, Urbana, IL 61801, USA
<sup>3</sup> Dept. für Physik und Astronomie, Universität Basel, CH-4056 Basel, Switzerland
ABSTRACT: We analyze data on deep inelastic scattering of electrons from the proton using ideas from standard many-body theory involving bound constituents subject to interactions. This leads us to expect, at large three-momentum transfer $`𝐪`$, scaling in terms of the variable $`\stackrel{~}{y}=\nu |𝐪|`$. The response at constant $`|𝐪|`$ scales well in this variable. Interaction effects are manifestly displayed in this approach. They are illustrated in two examples.
PACS: 13.60.Hb
The cross section for deep inelastic scattering (DIS) of electrons on unpolarized protons (and, mutatis mutandis, neutrons) is usually expressed as :
$$\frac{d^2\sigma }{d\mathrm{\Omega }d\nu }=\sigma _M\left[2W_1(|𝐪|,\nu )tan^2\frac{\theta }{2}+W_2(|𝐪|,\nu )\right],$$
(1)
where $`\sigma _M`$ is the Mott cross section, $`\theta `$ the scattering angle and $`\nu `$ the energy transfer. In this brief note we discuss only $`W_1`$.
$`W_1`$ is generally considered as a function of $`Q^2=|𝐪|^2\nu ^2`$ and Bjorken $`x=Q^2/2m\nu `$ for scattering by protons initially at rest in laboratory frame. The advantages of using the Lorentz scalar variables $`Q^2`$ and $`x`$ are discussed in standard texts . The data show that $`W_1(Q^2,x)`$ obeys Bjorken scaling at large values of $`Q^2`$; it depends primarily on $`x`$. The weak dependence of $`W_1`$ on $`Q^2`$ is well understood via the perturbative QCD theory developed by Gribov, Lipatov, Altarelli and Parisi (GLAP) . The DIS data are usually interpreted by going to the infinite momentum frame where $`x`$ is identified as the fraction of the momentum carried by the quark responsible for the deep inelastic scattering. This interpretation has been very helpful in understanding DIS and in interpreting high energy reactions.
In many-body theory it is natural to study the response of a system in its rest frame at fixed values of the momentum transfer $`𝐪`$ as a function of the energy transfer $`\nu `$. The response to a scalar probe, for example, is viewed as the distribution of the strength of the state $`_ie^{i𝐪𝐫_𝐢}|0`$, created by the probe, among the eigenstates of the system belonging to momentum $`𝐪`$. Fig.1
shows the various domains of the response of a proton in the $`|𝐪|`$, $`\nu `$ plane. The thick line $`|𝐪|=\nu `$ separates the spacelike response above the line, and the time like below the line. The thin line shows the $`ep`$ elastic scattering kinematical condition:
$$\nu _{el}=\sqrt{(|𝐪|^2+m^2)}m.$$
(2)
At large values of $`|𝐪|`$ the $`\nu _{el}=|𝐪|m`$ up to terms of order $`m^2/2|𝐪|`$. There can not be any response above the line $`\nu _{el}(|𝐪|)`$ since none of the target states can have energy less than $`\sqrt{(|𝐪|^2+m^2)}`$. The dashed lines show parabolae: $`|𝐪|^2\nu ^2=Q^2`$ for $`Q^2=5`$ and 10 GeV<sup>2</sup>. They intersect the $`\nu _{el}(|𝐪|)`$ curve at $`x=1`$ and approach the $`\nu =|𝐪|`$ line at $`\nu \mathrm{}`$ or equivalently as $`x0`$. In most of the literature the authors have considered the variation of $`W_1`$ along these parabolae which do not enter the timelike region. Here we study the variation of $`W_1`$ along the dash-dot lines, which have constant $`|𝐪|`$ and enter the timelike region.
In fig.2 we show the proton $`W_1^p(𝐪,\nu )`$, obtained from the MRS(A) fit of ref. to $`e`$-$`p`$
and other cross sections, at several values of $`|𝐪|`$ as a function of $`\nu |𝐪|\stackrel{~}{y}`$. The data show that at large values of $`𝐪`$ the $`W_1^p(𝐪,\nu )`$ depends primarily on $`\stackrel{~}{y}`$. This scaling has a simple interpretation in the many-body theory, related to the well known $`y`$-scaling , as will be discussed below. Since $`\nu |𝐪|=m\xi `$, where $`\xi `$ is the Nachtmann scaling variable, it is closely related to Bjorken scaling as well. The Nachtmann variable, which results from operator product expansion studies, coincides with $`x`$ at large $`Q^2`$ and is generally used to extend the applicability of Bjorken scaling to lower $`Q^2`$. The small scaling violations seen in fig.2 originate from the gluonic radiative corrections as in the standard approach based on GLAP evolution equations . Contrary to the case of $`x`$-scaling, in $`\stackrel{~}{y}`$ scaling $`Q^2`$ has a large variation, from zero at $`|𝐪|=\nu `$ ($`\stackrel{~}{y}=0)`$ to $`2m|𝐪|`$ at $`\nu =|𝐪|m`$ ($`\stackrel{~}{y}=m`$), which does not appear to spoil the quality of scaling.
Many-body theory views the deep inelastic response as follows. The probe creates states $`|i(𝐪+𝐤);(𝐤)`$ by hitting a bound constituent $`i`$ of the system with initial momentum $`𝐤`$ and the residual system in the state $``$ with momentum $`𝐤`$. In the plane wave impulse approximation (PWIA) the final state interaction (FSI) between the struck constituent $`i`$ and the residual system $``$ are neglected. In this approximation the energy of the state $`|i(𝐪+𝐤);(𝐤)`$ is:
$$E(i(𝐪+𝐤);(𝐤))=|𝐪|+k_{}+E((𝐤))+termsoforder\frac{1}{|𝐪|},$$
(3)
where $`k_{}`$ is the projection of $`𝐤`$ in the direction of $`𝐪`$. At large $`|𝐪|`$ the terms of order $`1/|𝐪|`$ can be neglected, and the response due to the excitation of the state $`|i(𝐪+𝐤);(𝐤)`$ occurs at energy transfer:
$$\nu =E(i(𝐤+𝐪);(𝐤))m=|𝐪|+𝐤_{}+E((𝐤))m.$$
(4)
Since the $`\nu |𝐪|`$ is independent of $`𝐪`$ at large $`|𝐪|`$, the response depends only on $`\stackrel{~}{y}`$ in the PWIA; therefore it scales.
The fact that the observed response of the proton, as seen in fig.2, scales with $`\stackrel{~}{y}`$ does not necessarily imply that PWIA is valid. In general the effects of the FSI on the response may not be negligible. Treatments of the FSI for very different cases of inclusive scattering of a probe from a composite system have shown that the main effect of FSI results in a folding of the PWIA response:
$$W_1(|𝐪|,\nu )=𝑑\nu ^{}W_{1,PWIA}(|𝐪|,\nu ^{})f(|𝐪|,\nu ,\nu ^{}).$$
(5)
The scaling of $`W_1(|𝐪|,\nu )`$ with $`\stackrel{~}{y}`$ can occur when the folding function representing the effect of the final state interactions becomes independent of $`|𝐪|`$. For example, in the Glauber approximation, the folding function for the quasi-free scattering of electrons by nuclei becomes independent of $`|𝐪|`$ at $`|𝐪|>2`$ GeV, as has been recently discussed by Benhar . Weinstein and Negele have shown that an analogous $`y`$-scaling in hard sphere Bose gas occurs even though the FSI effects are strong.
By boosting the state $`|(𝐤),i(𝐤)`$ to large velocity, ignoring interaction between $`i`$ and $``$, it can be shown that
$$\xi =\frac{|𝐪|\nu }{m}=1\frac{E((𝐤))+k_{}}{m}$$
(6)
has the usual meaning of the fraction of the momentum carried by the struck particle i.
The complete response is obtained by summing over all the final states. Therefore:
$$W_1(𝐪,\nu )=\underset{i}{}d^3k𝑑e\sigma _1(𝐪,\nu ,𝐤,e)P_i(𝐤,e)\delta (\stackrel{~}{y}k_{}e),$$
(7)
where the spectral function $`P_i(𝐤,e)`$ is given by:
$$P_i(𝐤,e)=\underset{}{}|(𝐤),i(𝐤)|0|^2\delta (mE((𝐤))e).$$
(8)
The $`𝐤,e`$ can be considered as the initial momentum and energy of the struck particle, which is treated as bound, and therefore not on the mass shell. The $`\sigma _1`$ is the transverse cross-section for lepton scattering by a bound spin 1/2 point particle $`i`$, divided by $`\sigma _M`$. In the physical spacelike $`(\stackrel{~}{y}<0)`$ region, $`\sigma _1=q_i^2`$ in the large $`𝐪`$ limit, neglecting the mass $`m_i`$ of $`i`$. Here $`q_i`$ is the charge of $`i`$. This gives:
$$W_1(𝐪,\nu )=W_1(\stackrel{~}{y})=\underset{i}{}q_i^2d^3k𝑑eP_i(𝐤,e)\delta (\stackrel{~}{y}k_{}e).$$
(9)
This Eq. provides the relation between the familiar $`F_1(\xi )=mW_1(\xi )`$ structure function and the $`P_i(𝐤,e)`$ spectral function. The $`\delta `$-function implies that $`\stackrel{~}{y}=k_{}+e`$, closely resembling the scaling variable $`y`$ used in quasi-elastic electron-nucleus scattering, where it is associated with the parallel momentum of the struck nucleon.
This simple picture of the response will be modified by the color confining interactions. The mass of the nucleon contains confinement interaction contribution, while it is omitted in the energy of the struck quark, $`|𝐪|+k_{}`$. It therefore must be included in the energy $`E()`$ of the residual system. We expect that the confinement energy does not change significantly in the time duration of the DIS, and its main influence is via the wave functions $`|0`$ and $`|`$. However, it could also contribute to the FSI folding function (Eq.(5)).
An interesting feature of fig.2 concerns the width of the response, which amounts to only few hundred MeV independent of the value of $`|𝐪|`$. This implies that deep inelastic scattering has an intrinsic energy scale of few hundred MeV. The main part of the transferred energy, of order $`|𝐪|`$, goes into the kinetic energy of the struck constituent, and does not play any interesting role in the dynamics of the target system. Therefore changes in the energy $`E((𝐤))`$ of the residual system of order 100 MeV have significant effect on the $`W_1(\stackrel{~}{y})`$, or equivalently on $`F_1(\xi )`$. In the following we discuss two observed effects of $`E((𝐤))`$ on the response.
The first, studied by Close and Thomas , concerns the difference between the responses due to valence u and d quarks in the proton. Let $`V_u(\stackrel{~}{y})`$ and $`V_d(\stackrel{~}{y})`$ be the contributions of valence u and d quarks to the $`W_1^p(\stackrel{~}{y})`$. When the lepton strikes the valence d quark, the remaining two valence u quarks are left in the residual state $`_1`$ with spin 1. In contrast, when a valence u quark is struck, the residual u-d pair is in states $`_0`$ with spin 0 with probability 0.75, and $`_1`$ with probability 0.25. Therefore $`\chi _1(\stackrel{~}{y})9V_d(\stackrel{~}{y})`$ is the response due to final states $`_1`$ (normalized to unit particle charge), while $`\chi _0(\stackrel{~}{y})1.5(V_u(\stackrel{~}{y})2V_d(\stackrel{~}{y})`$ is that for $`_0`$. The
$$E(_1(𝐤))E(_0(𝐤))\frac{2}{3}(m_\mathrm{\Delta }m),$$
(10)
in perturbation theory, and therefore we expect $`\chi _1(\stackrel{~}{y})`$ to be shifted by $`0.2`$ GeV from $`\chi _0(\stackrel{~}{y})`$. Fig.3
shows that these responses obtained from the MRS(A) parton distributions at $`|𝐪|=10`$ GeV are indeed shifted by $`0.1`$ GeV from each other, at $`\stackrel{~}{y}<0.2`$. In the PWIA, this shift should be independent of $`\stackrel{~}{y}`$, provided the color magnetic interaction can be treated perturbatively. The fact that the shift is only $``$ 0.1 GeV indicates that it has nonperturbative contributions. Differences in FSI can also have an influence.
The second example concerns the modification of the deep inelastic response by nuclear effects as first observed by the EMC collaboration . The EMC ratio $`R_A(x)`$ of the cross section per nucleon, for nucleus with mass number $`A`$ to that for the deuteron, does not show any $`Q^2`$ dependence within the experimental errors. The observed $`R_A(x)`$ has been extrapolated using Local Density Approximation to obtain the ratio $`R_{NM}(x)`$ for uniform nuclear matter .
In fig.4 we show the $`W_1^d(\stackrel{~}{y})`$ for the deuteron calculated from the MRS(A) fits, and the $`W_1^{NM}(\stackrel{~}{y})`$ at $`|𝐪|=10`$ GeV. As we see from this figure, the nuclear matter response is quite similar to that of the deuteron. It is a bit broader, due to the Fermi motion of nucleons in matter, but mainly it is shifted towards higher $`\nu `$ due to nuclear binding. Fig.4 also shows the response of noninteracting nucleons distributed according to the momentum distribution of nucleons in nuclear matter, also calculated using realistic interactions . The observed response is shifted relative the Fermi motion broadened response by $``$ 40–60 MeV, which is comparable to the average nucleon removal energy of $``$ 62 MeV in nuclear matter . For these reasons, the conventional nuclear physics approach is quite successful in describing the EMC ratio for nuclear matter at $`x>0.4`$ .
In the parton model the struck particle is assumed to be on mass-shell before and after the interaction with the electron . In this case
$$\nu =\sqrt{m_i^2+(𝐤+𝐪)^2}\sqrt{m_i^2+𝐤^2}|𝐪|,$$
(11)
and all of the response is at negative $`\stackrel{~}{y}`$, in the spacelike region. The same is valid at the leading twist-two order of the operator product expansion .
In many body theory, however, a timelike response occurs either due to initial state interactions, which can make $`E((𝐤))`$ large enough to give a positive right hand side of eq.(4), or because of FSI. The initial energy of the struck constituent is identified with $`e=mE((𝐤))`$ and not with the on shell energy $`\sqrt{m_i^2+|𝐤|^2}`$ used in eq.(11). This timelike response contributes to various sum rules. For example, the Coulomb sum in quasi-elastic electron nucleus scattering is defined as the integral of the longitudinal response over both space and time like regions. The longitudinal response of deuterium has been calculated with realistic forces ; it extends into the timelike region, and that region has to be included to fulfill the Coulomb sum.
In fact, the shifts in $`W_1(\stackrel{~}{y})`$ illustrated in figs.3 and 4 will move part of the response into the timelike region, barring FSI effects, and thus lead to a violation of sum rules involving $`W_1(\stackrel{~}{y}<0)`$.
In conclusion, we obtain new insights in the deep inelastic response of nucleons by applying standard many-body theory and relate the scaling function to the nucleon spectral function in the lab frame. The natural scaling variable of many body theory, $`\stackrel{~}{y}`$, equals $`m\xi `$ of the conventional approach to DIS. While $`\stackrel{~}{y}`$ scaling is derived assuming bound constituents which are subject to initial and final state interaction, $`\xi `$\- or $`x`$-scaling is obtained assuming free constituents without FSI. The occurrence of scaling thus cannot automatically be taken as evidence for scattering from free constituents.
Acknowledgements: The authors thank A.W. Thomas and D. Beck for interesting discussions. This work was supported by the US and Swiss National Science Foundations.
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# Geometry of River Networks III: Characterization of Component Connectivity
## I Introduction
This is the last paper in a series of three on the geometry of river networks. In the first dodds2000ua we examine in detail the description of river networks by scaling laws maritan96a ; rodriguez-iturbe97 ; dodds99pa ; dodds2000pa and the evidence for universality. Additional introductory remarks concerning the motivation of the overall work are to found in this first paper. In the second article dodds2000ub we address distributions of the basic components of river networks, stream segments and sub-networks. Here, we provide an analysis complementary to the work of the second paper by establishing a description of how river network components fit together. As before, we are motivated by the premise that while relationships of mean quantities are primary in any investigation, the behavior of higher order moments potentially and often do encode significant information.
Our purpose then is to investigate the distributions of quantities which describe the architecture of river networks. The goal is to quantify these distributions and, where this is not possible, to quantify fluctuations. In particular, we center our attention on Tokunaga’s law tokunaga66 ; tokunaga78 ; tokunaga84 which is a statement about network architecture describing the tributary structure of streams. Since Tokunaga’s law can be seen as the main part of a platform from which all other river network scaling laws follow dodds99pa , it is an obvious starting point for the investigation of fluctuations in river network structure. We use data from the Schediegger model of random networks scheidegger67 and the Mississippi river. We find the distributions obtained from these two disparate sources agree very well in form. We are able to write down scaling forms of all distributions studied. We observe a number of distributions to be exponential, therefore requiring only one parameter for their description. As a result, we introduce a dimensionless scale $`\xi _t`$, finding it to be sufficient to describe the fluctuations present in Tokunaga’s law and thus potentially all river network scaling laws. Significantly, we observe the spatial distribution of stream segments to be random implying we have reached the most basic description of network architecture.
Tokunaga’s law is also intimately connected with drainage density, $`\rho `$, a quantity which will be used throughout the paper. Drainage density is a measure of stream concentration or, equivalently, how a network fills space. We explore this connection in detail, showing how simple assumptions regarding drainage density lead to Tokunaga’s law.
The paper is structured as follows. We first outline Horton-Strahler stream ordering which provides the necessary descriptive taxonomy for river network architecture. We then define Tokunaga’s law and introduce a scaling law for a specific form of drainage density. We briefly describe Horton’s laws for stream number and length and some simple variations. (Both stream ordering and Horton’s laws are covered in more detail in dodds2000ub ). We show that the scaling law of drainage density may be taken as an assumption from which all other scaling laws follow. We also briefly consider the variation of basin shapes (basin allometry) in the context of directedness. This brings us to the focal point of the paper, the identification of a statistical generalization of Tokunaga’s law. We first examine distributions of numbers of tributaries (side streams) and compare these with distributions of stream segment lengths. We observe both distributions to be exponential leading to the notion that stream segments are distributed randomly throughout a network. The presence of exponential distributions also leads to the introduction of the characteristic number $`\xi _t`$ and the single length-scale $`\xi _{l^{(\mathrm{s})}}\xi _t`$. We then study the variation of tributary spacing along streams so as to understand fluctuations in drainage density and again find the signature of randomness. This leads us to develop a joint probability distribution connecting the length of a stream with the frequency of its side streams.
## II Definitions
### II.1 Stream ordering
Horton-Strahler stream orderinghorton45 ; strahler57 breaks a river network down into a set of stream segments. The method can be thought of as an iterative pruning. First, we define a source stream as the stream section that runs from a channel head to the first junction with another stream. These source streams are classified as the first-order stream segments of the network. Next, we remove all source streams and identify the new source streams of the remaining network. These are the network’s second-order stream segments. The process is repeated until one stream segment is left of order $`\mathrm{\Omega }`$. The order of the basin is then defined to be $`\mathrm{\Omega }`$ (we will use the words basin and network interchangeably).
In discussing network architecture, we will speak of side streams and absorbing streams. A side stream is any stream that joins into a stream of higher order, the latter being the absorbing stream. We will denote the orders of absorbing and side streams by $`\mu `$ and $`\nu `$ but when referring to an isolated stream or streams where their relative rank is ambiguous, we will write stream order as $`\omega `$, subscripted as seems appealing.
Central to our investigation of network architecture is stream segment length. As in dodds2000ub , we denote this length by $`l_\omega ^{(\mathrm{s})}`$ for a stream segment of order $`\omega `$. We will also introduce a number of closely related lengths which describe distances between side streams. When referring to streams throughout we will specifically mean stream segments of a particular order unless otherwise indicated. This is to avoid confusion with the natural definition of a stream which is the path from a point on a network moving upstream to the most distant source. For an order $`\omega `$ basin, we denote this main stream length by $`l_\omega `$
Note also that we consider river networks in planform, i.e., as networks projected onto the horizontal (or gravitationally flat) plane. This simplification poses no great concern for the analysis of large scale networks such as the Mississippi but must be considered in the context of drainage basins with significant relief.
### II.2 Tokunaga’s law
Defining a stream ordering on a network allows for a number of well-defined measures of connectivity, stream lengths and drainage areas. Around a decade after the Strahler-improved stream ordering of Horton appeared, Tokunaga introduced the idea of measuring side stream statistics tokunaga66 ; tokunaga78 ; tokunaga84 . This technique arguably provides the most useful measurement based on stream ordering but has only recently received much attention dodds99pa ; cui99 ; turcotte98 ; peckham95 . The idea is simply, for a given network, to count the average number of order $`\nu `$ side streams entering an order $`\mu `$ absorbing stream. This gives $`T_{\mu ,\nu }`$, a set of double-indexed parameters for a basin. Note that $`\mathrm{\Omega }\mu >\nu 1`$, so we can view the Tokunaga ratios as a lower triangular matrix. An example for the Mississippi river is shown in Table 1 <sup>1</sup><sup>1</sup>1 The network for the Mississippi was extracted from a topographic dataset constructed from three arc second USGS Digital Elevation Maps, decimated by averaging to approximately 1000 meter horizontal resolution (www.usgs.gov). At this grid scale, the Mississippi was found to be an order $`\mathrm{\Omega }=11`$ basin. The same data is represented pictorially in Figure 1 in what we refer to as a Tokunaga graph.
Tokunaga made several key observations about these side stream ratios. The first is that because of the self-similar nature of river networks, the $`T_{\mu ,\nu }`$ should not depend absolutely on either of $`\mu `$ or $`\nu `$ but only on the relative difference, i.e., $`k=\mu \nu `$. The second is that in changing the value of $`k=\mu \nu `$, the $`T_{\mu ,\nu }`$ must themselves change by a systematic ratio. These statements lead to Tokunaga’s law:
$$T_{\mu ,\nu }=T_k=T_1(R_T)^{k1}.$$
(1)
Thus, only two parameters are necessary to characterize the set of $`T_{\mu ,\nu }`$: $`T_1`$ and $`R_T`$.
The parameter $`T_1>0`$ is the average number of side streams of one order lower than the absorbing stream, typically on the order of 1.0–1.5. Since these side streams of one order less are the dominant side streams of the basin, their number estimates the basin’s breadth. In general, larger values of $`T_1`$ correspond to wider basins while smaller values are in keeping with basins with relatively thinner profiles.
The ratio $`R_T>1`$ measures how the density of side streams of decreasing order increases. It is a measure of changing length scales and has a simple interpretation with respect to Horton’s laws which we describe below. Thus, already inherent in Tokunaga’s law is a generalization of drainage density $`\rho `$. the usual definition of which is given as follows. For a given region of landscape with area $`A`$ with streams totalling in length $`L`$, $`\rho =L/A`$ and has the dimensions of an inverse length scale horton45 . One may think of $`\rho `$ as the inverse of the typical distance between streams, i.e., the characteristic scale beyond which erosion cannot more finely dissect the landscape horton45 . In principle, drainage density may vary from landscape to landscape and also throughout a single region. Below, we will turn this observation about Tokunaga’s law around to show that all river network scaling laws may be derived from an expanded notion of drainage density.
Even though the number of side streams entering any absorbing stream must of course be an integer, Tokunaga’s ratios are under no similar obligation since they are averages. Nevertheless, Tokunaga’s law provides a good sense of the structure of a network albeit at a level of averages. One of our main objectives here is to go further and consider fluctuations about and the full distributions underlying the $`T_{\mu ,\nu }`$.
Finally, a third important observation of Tokunaga is that two of Horton’s laws follow from Tokunaga’s law, which we next discuss.
### II.3 Horton’s laws
We review Horton’s laws horton45 ; schumm56a ; dodds99pa and then show how self-similarity and drainage density lead to Tokunaga’s law, Horton’s laws and hence all other river network scaling laws.
The relevant quantities for Horton’s relations are $`n_\omega `$, the number of order $`\omega `$ streams, and $`l_\omega `$, the average main stream length (as opposed to stream segment length $`l_\omega ^{(\mathrm{s})}`$) of order $`\omega `$ basins. The laws are simply that the ratio of these quantities from order to order remain constant:
$$\frac{n_{\omega +1}}{n_\omega }=1/R_n\text{and}\frac{l_{\omega +1}}{l_\omega }=R_l,$$
(2)
for $`\omega 1`$. Note the definitions are chosen so that all ratios are greater than unity. The number of streams decreases with order while all areas and lengths grow.
A similar law for basin areas horton45 ; schumm56a states that $`a_{\omega +1}/a_\omega =R_a`$ where $`a_\omega `$ is the average drainage area of an order $`\omega `$ basin. However, with the assumption of uniform drainage density it can be shown that $`R_nR_a`$ dodds99pa so we are left with the two independent Horton laws of equation (2).
As in dodds2000ub , we consider another Horton-like law for stream segment lengths:
$$\frac{l_{\omega +1}^{(\mathrm{s})}}{l_\omega ^{(\mathrm{s})}}=R_{l^{(\mathrm{s})}}.$$
(3)
As we will show, the form of the distribution of the variable $`T_{\mu ,\nu }`$ is a direct consequence of the distribution of $`l_\omega ^{(\mathrm{s})}`$.
Tokunaga showed that Horton’s laws of stream number and stream length follow from what we have called Tokunaga’s law, equation (1). For example, the solution of a difference equation relating the $`n_\omega `$ and the $`T_k`$ leads to the result $`R_n=A_T+\left[A_T^22R_T\right]^{1/2}`$ where $`A_T=(2+R_T+T_1)/2`$ for $`\mathrm{\Omega }=\mathrm{}`$ and a more complicated expression is obtained for finite $`\mathrm{\Omega }`$ dodds99pa ; tokunaga78 ; tokunaga84 ; peckham95 . In keeping with our previous remarks on $`T_1`$, this expression for $`R_n`$ shows that an increase in $`T_1`$ will increase $`R_n`$ which, since $`R_aR_n`$, corresponds to a network where basins tend to be relatively broader. Our considerations will expand significantly on this connection between the network descriptions of Horton and Tokunaga.
## III The implications of a scaling law for drainage density
We now introduce a law for drainage density based on stream ordering. We write $`\rho _{\mu ,\nu }`$ for the number of side streams of order $`\nu `$ per unit length of order $`\mu `$ absorbing stream. We expect these densities to be independent of the order of the absorbing stream and so we will generally use $`\rho _\nu `$. The typical length separating order $`\nu `$ side streams is then $`1/\rho _\nu `$. Assuming self-similarity of river networks, we must have
$$\rho _{\nu +1}/\rho _\nu =1/R_\rho $$
(4)
where $`R_\rho >1`$ independent of $`\nu `$.
All river network scaling laws in the planform may be seen to follow from this relationship. Consider an absorbing stream of order $`\mu `$. Self-similarity immediately demands that the number of side streams of order $`\mu 1`$ must be statistically independent of $`\mu `$. This number is of course $`T_1`$. Therefore, the typical length of an order $`\mu `$ absorbing stream must be
$$l_\mu ^{(\mathrm{s})}=T_1/\rho _{\mu 1}.$$
(5)
Using equation (4) to replace $`\rho _{\mu 1}`$ in the above equation, we find
$$T_1/\rho _{\mu 1}=R_\rho T_1/\rho _{\mu 2}.$$
(6)
Thus, $`T_2=R_\rho T_1`$ and, in general $`T_k=(R_\rho )^{k1}T_1`$. This is Tokunaga’s law and we therefore have
$$R_\rho R_T.$$
(7)
Equation (4) and equation (5) also give
$$l_\mu ^{(\mathrm{s})}=T_1/\rho _{\mu 1}=R_\rho T_1/\rho _{\mu 2}=R_\rho l_{\mu 1}^{(\mathrm{s})}.$$
(8)
On comparison with equation (3), we see that the above is our Hortonian law of stream segment lengths and that
$$R_\rho R_{l^{(\mathrm{s})}}.$$
(9)
As $`R_{l^{(\mathrm{s})}}`$ is the basic length-scale ratio in the problem, we rewrite equation (4), our Hortonian law of drainage density, as
$$\rho _{\nu +1}/\rho _\nu =1/R_{l^{(\mathrm{s})}}.$$
(10)
The above statement becomes our definition of the self-similarity of drainage density.
## IV Basin allometry
Given that we have suggested the need for only a single relevant length ratio, we must remark here on basin allometry. Allometry refers to the relative growth or scaling of a shape’s dimensions and was originally introduced in the context of biology huxley36 . A growth or change being allometric usually implies it is not self-similar. A longstanding issue in the study of river networks has been whether or not basins are allometric rodriguez-iturbe97 ; dodds2000pa ; hack57 .
Consider two basins described by $`(L_1,W_1)`$ and $`(L_2,W_2)`$ within the same system where $`L_i`$ is a characteristic longitudinal basin length and $`W_i`$ a characteristic width. The basins being allometric means that $`(W_1/W_2)=(L_1/L_2)^H`$ where $`H<1`$. Thus, two length ratios are needed to describe the allometry of basins. If we consider basins defined by stream ordering then we have the Horton-like ratios $`R_L`$ and $`R_W=R_L^H`$. Now, when rescaling an entire basin, streams roughly aligned with a basin’s length will rescale with the factor $`R_L`$ and those perpendicular to the basin’s axis will rescale differently with $`R_W`$. This creates a conundrum: how can basins be allometric ($`R_LR_W`$) and yet individual streams be self-similar ($`R_L=R_W`$) as implied by Horton’s laws?
We contend the answer is that allometry must be restricted to directed networks and that self-similarity of basins must hold for non-directed networks. This is in agreement with Colaiori et al. colaiori97 who also distinguish between self-similar and allometric river basins although we stress here the qualification of directedness. Directed networks have a global direction of flow in which the direction of each individual stream flow has a positive component. A basic example is the random model of Scheidegger scheidegger67 which we describe below. For a directed network, $`R_L=R_l`$, and the rescaling of basin sizes matches up with the rescaling of stream lengths regardless of how the basin’s width rescales since all streams are on average aligned with the global direction of flow. Hence, our premise that streams rescale in a self-similar way is general enough to deal with systems whose basins rescale in an allometric fashion.
In considering the allometry of basins, we must also address the additional possibility that individual stream lengths may scale non-trivially with basin length. In this case, the main stream length $`l`$ would vary with the longitudinal basin length $`L`$ as $`lL^d`$. This is typically a weak dependence with $`1.0<d<1.15`$ maritan96a ; tarboton90 . Note that Horton’s laws still apply in this case. The exponent $`d`$ plays a part in determining whether or not a basin scales allometrically. The exponent $`H`$ introduced in the discussion of basin allometry can be found in terms of Horton’s ratios (or equivalently Tokunaga’s parameters) and $`d`$ as $`H=d\mathrm{ln}R_n/\mathrm{ln}R_l1`$ dodds99pa .
Thus, for a directed network $`d=1`$ and $`H1`$ (e.g., Scheidegger scheidegger67 ) whereas for undirected, self-similar networks $`H=1`$ and $`d1`$ (e.g., random undirected networks manna92 ; manna96 ). River networks are in practice often neither fully directed or undirected. Scaling laws observed in such cases will show deviations from pure scaling that may well be gradual and difficult to detect dodds2000ua .
## V Tokunaga distributions
The laws of Tokunaga and Horton relate averages of quantities. In the remainder of this paper, we investigate the underlying distributions from which these averages are made. We are able to find general scaling forms of a number of distributions and in many cases also identify the basic form of the relevant scaling function.
To aid and motivate our investigations, we examine, as we have done in both dodds2000ua and dodds2000ub , a simple model of directed random networks that was first introduced by Scheidegger scheidegger67 Since we make much of use this model in the present work, we provide a self-contained discussion. Consider the triangular lattice of sites oriented as in Figure 2. At each site of the lattice a stream flow direction is randomly chosen between the two possible diagonal directions shown. It is therefore trivial to generate the model on a large scale, allowing for a thorough investigation of its river network statistics. The small, tilted box with a dashed boundary represents the area drained by the enclosed site. As with many discrete-space models, the details of the underlying lattice are unimportant. On a square lattice, the model’s streams would have three choices of flow, two diagonals and straight down the page. However, the choice of a triangular lattice does simplify implementation and calculation of statistics. For example, only one tributary can exist at each site along a stream and stream paths and basin boundaries are precisely those of the usual discrete-space random walk feller68I .
Since random walks are well understood, the exponents of many river network scaling laws are exactly known for the Scheidegger model takayasu88 ; takayasu89a ; takayasu91 ; huber91 and analogies may also be drawn with the Abelian sandpile model dhar99 . For example, a basin’s boundaries being random walks means that a basin of length $`L`$ will typically have a width $`WL^{1/2}`$ which gives $`H=1/2`$. Since the network is directed, stream length is the same as basin length, $`l=L`$, so we trivially have $`d=1`$. Basin area $`a`$ is estimated by $`WLL^{3/2}l^{3/2}`$ so $`la^{2/3}`$ giving Hack’s law with an exponent of $`2/3`$ hack57 .
Nevertheless, the Tokunaga parameters and the Horton ratios are not known analytically. Estimates from previous work dodds99pa find $`T_11.35`$, $`R_l=R_T3.00`$ and $`R_n5.20`$. Data for the present analysis was obtained on $`L=10^4`$ by $`W=3\times 10^3`$ lattices with periodic boundaries. Given the self-averaging present in any single instance on these networks, ensembles of 10 were deemed sufficient.
We first examine the distributions of Tokunaga ratios $`T_{\mu ,\nu }`$ and observe a strong link to the underlying distribution of $`l_\mu ^{(\mathrm{s})}`$. Both are well described by exponential distributions. To understand this link, we next consider the distances between neighboring side streams of like order. This provides a measure of fluctuations in drainage density and again, exponential distributions appear. We are then in a position to develop theory for the joint probability distribution between the Tokunaga ratios and stream segment lengths and, as a result, the distribution for the quantity $`T_{\mu ,\nu }/l_\mu ^{(\mathrm{s})}`$ and its inverse. In the limit of large $`\mu `$, the $`T_{\mu ,\nu }`$ are effectively proportional to $`l_\mu ^{(\mathrm{s})}`$ and all fluctuations of the former exactly follow those of the latter.
All investigations are initially carried out for the Scheidegger model where we may generate statistics of ever-improving quality. We find the same forms for all distributions for the Mississippi data (and for other river networks not presented here) and provide some pertinent examples. Perhaps the most significant benefit of the simple Scheidegger model is its ability to provide clean distributions whose form we can then search for in real data.
Figure 3 shows the distribution of the number of order $`\nu =2`$ side streams entering an order $`\mu =6`$ absorbing stream for the Scheidegger model. At first, it may seem surprising that this is not a single-peaked distribution centered around $`T_{\mu ,\nu }`$ dying off for small and large values of $`T_{\mu ,\nu }`$.
The distribution of $`T_{\mu ,\nu }`$ in Figure 3 is clearly well described by an exponential distribution. This can also be seen upon inspection of Figures 4(a) and 4(b). Figure 4(a) shows normalized distributions of $`T_{\mu ,\nu }`$ for $`\nu =2`$ and varying absorbing stream order $`\mu =4`$, 5 and 6. These distributions (plus the one for absorbing stream order $`\mu =7`$) are rescaled and presented in Figure 4(b). The single form thus obtained suggests a scaling form of the $`T_{\mu ,\nu }`$ distribution is given by
$$P(T_{\mu ,\nu })=(R_{l^{(\mathrm{s})}})^\mu F\left[T_{\mu ,\nu }(R_{l^{(\mathrm{s})}})^\mu \right].$$
(11)
where $`F`$ is an exponential scaling function. However, this only accounts for variations in $`\mu `$, the order of the absorbing stream.
Figures 5(a) and 5(b) show that a similar rescaling of the distributions may be effected when $`\nu `$ is varied. In this case, the data is for the Mississippi. The rescaling is now by $`R_{l^{(\mathrm{s})}}`$ rather than $`R_{l^{(\mathrm{s})}}^1`$ and equation (11) is improved to give
$$P(T_{\mu ,\nu })=(R_{l^{(\mathrm{s})}})^{\mu \nu 1}P_T\left[T_{\mu ,\nu }/(R_{l^{(\mathrm{s})}})^{\mu \nu 1}\right].$$
(12)
The function $`P_T`$ is a normalized exponential distribution independent of $`\mu `$ and $`\nu `$,
$$P_T(z)=\frac{1}{\xi _t}e^{z/\xi _t},$$
(13)
where $`\xi _t`$ is the characteristic number of side streams of one order lower than the absorbing stream, i.e., $`\xi _t=T_1`$. For the Mississippi, we observe $`\xi _t1.1`$ whereas for the Scheidegger model, $`\xi _t1.35`$. As expected, the Tokunaga distribution is dependent only on $`k=\mu \nu `$ so we can write
$$P(T_k)=(R_{l^{(\mathrm{s})}})^{k1}P_T\left[T_{\mu ,\nu }/(R_{l^{(\mathrm{s})}})^{k1}\right].$$
(14)
with $`P_T`$ as above.
## VI Distributions of stream segment lengths and randomness
As we have suggested, the distributions of the Tokunaga ratios depend strongly on the distributions of stream segment lengths. Figure 6 is the indicates why this is so. The form of the underlying distribution is itself exponential. We have already examined this fact extensively in dodds2000ub and here we develop its relationship with the Tokunaga distributions.
Figures 6(a) and 6(b) show that the distributions of $`l_\mu ^{(\mathrm{s})}`$ can be rescaled in the same way as the Tokunaga distributions. Thus, we write the distribution for stream segment lengths as dodds2000ub
$$P(l_\mu ^{(\mathrm{s})})=(R_{l^{(\mathrm{s})}})^{\mu +1}P_{l^{(\mathrm{s})}}\left[l_\mu ^{(\mathrm{s})}/(R_{l^{(\mathrm{s})}})^{\mu +1}\right].$$
(15)
As for $`P_T`$, the function $`P_{l^{(\mathrm{s})}}`$ is a normalized exponential distribution
$$P_{l^{(\mathrm{s})}}(z)=\frac{1}{\xi _{l^{(\mathrm{s})}}}e^{z/\xi _{l^{(\mathrm{s})}}},$$
(16)
where, in a strictly self-similar network, $`\xi _{l^{(\mathrm{s})}}`$ is the characteristic length of first-order stream segments, i.e., $`\xi _{l^{(\mathrm{s})}}=l_1^{(\mathrm{s})}`$. (Note that in dodds2000ub we use $`\xi `$ for $`\xi _{l^{(\mathrm{s})}}`$ for ease of notation). We qualify this by requiring the network to be exactly self-similar because in most models and all real networks this is certainly not the case. As should be expected, there are deviations from scaling for the largest and smallest orders. Therefore, $`\xi _{l^{(\mathrm{s})}}`$ is the characteristic size of a first-order stream as determined by scaling down the average lengths of those higher order streams that are in the self-similar structure of the network. It is thus in general different from $`l_1^{(\mathrm{s})}`$.
We therefore see that the distributions of $`T_{\mu ,\nu }`$ and $`l_\mu ^{(\mathrm{s})}`$ are both exponential in form. Variations in $`l_\mu ^{(\mathrm{s})}`$ largely govern the possible values of the $`T_{\mu ,\nu }`$. However, $`T_{\mu ,\nu }`$ is still only proportional to $`l_\mu ^{(\mathrm{s})}`$ on average and later on we will explore the joint distribution from which these individual exponentials arise.
The connection between the characteristic number $`\xi _t`$ and the length-scale $`\xi _{l^{(\mathrm{s})}}`$ follows from equations (3), (5), and (10):
$$\xi _t=\rho _1R_{l^{(\mathrm{s})}}\xi _{l^{(\mathrm{s})}}.$$
(17)
This presumes exact scaling of drainage densities and in the case where this is not so, $`\rho _1`$ would be chosen so that $`(R_{l^{(\mathrm{s})}})^{\nu 1}\rho _1`$ most closely approximates the higher order $`\rho _\nu `$.
We come to an important interpretation of the exponential distribution as a composition of independent probabilities. Consider the example of stream segment lengths. We write $`\stackrel{~}{p}_\mu `$ as the probability that a stream segment of order $`\mu `$ meets with (and thereby terminates at) a stream of order at least $`\mu `$. For simplicity, we assume only one side stream or none may join a stream at any site. We also take the lattice spacing $`\alpha `$ to be unity so that stream lengths are integers and therefore equate with the number of links between sites along a stream. For $`\alpha 1`$, derivations similar to below will apply with $`l_\mu ^{(\mathrm{s})}`$ replaced by $`[l_\mu ^{(\mathrm{s})}/\alpha ]`$, where $`[]`$ denotes rounding to the nearest integer. Note that extra complications arise when the distances between neighboring sites are not uniform.
Consider a single instance of an order $`\mu `$ stream segment. The probability of this segment having a length $`l_\mu ^{(\mathrm{s})}`$ is given by
$$P(l_\mu ^{(\mathrm{s})})=\stackrel{~}{p}_\mu (1\stackrel{~}{p}_\mu )^{l_\mu ^{(\mathrm{s})}}.$$
(18)
where $`\stackrel{~}{p}_\mu `$ is the probability that an order an order $`\mu `$ stream segment terminates on meeting a stream of equal or higher order. We can re-express the above equation as
$$P(l_\mu ^{(\mathrm{s})})\stackrel{~}{p}_\mu \mathrm{exp}\{l_\mu ^{(\mathrm{s})}\mathrm{ln}(1\stackrel{~}{p}_\mu )^1\},$$
(19)
and upon inspection of equations (15) and (16) we make the identification
$$(R_{l^{(\mathrm{s})}})^{\mu 1}\xi _{l^{(\mathrm{s})}}=[\mathrm{ln}(1\stackrel{~}{p}_\mu )]^1,$$
(20)
which has the inversion
$$\stackrel{~}{p}_\mu =1e^{1/(R_{l^{(\mathrm{s})}})^{\mu 1}\xi _{l^{(\mathrm{s})}}}.$$
(21)
For $`\mu `$ sufficiently large such that $`\stackrel{~}{p}_\mu 1`$, we have the simplification
$$\stackrel{~}{p}_\mu 1/(R_{l^{(\mathrm{s})}})^{\mu 1}\xi _{l^{(\mathrm{s})}}.$$
(22)
We see that the probabilities satisfy the Horton-like scaling law
$$\stackrel{~}{p}_\mu /\stackrel{~}{p}_{\mu 1}=1/R_{l^{(\mathrm{s})}}.$$
(23)
Thus, we begin to see the element of randomness in our expanded description of network architecture. The termination of a stream segment by meeting a larger branch is effectively a spatially random process.
## VII Generalized drainage density
Having observed the similarity of the distributions of $`T_{\mu ,\nu }`$ and $`l_\mu ^{(\mathrm{s})}`$, we proceed to examine the exact nature of the relationship beteween the two. To do so, we introduce three new measures of stream length. These are $`l_{\mu ,\nu }^{(\mathrm{s},\mathrm{b})}`$, the distance from the beginning of an order $`\mu `$ absorbing stream to the first order $`\nu `$ side stream; $`l_{\mu ,\nu }^{(\mathrm{s},\mathrm{i})}`$, the distance between any two adjacent internal order $`\nu `$ side streams along an order $`\mu `$ absorbing stream; and $`l_{\mu ,\nu }^{\mathrm{s},\mathrm{e}}`$, the distance from the last order $`\nu `$ side stream to the end of an order $`\mu `$ absorbing stream. By analysis of these inter-tributary lengths, we will be able to discern the distribution of side stream location along absorbing streams. This leads directly to a more general picture of drainage density which we fully expand upon in the following section.
Figure 7 compares normalized distributions of $`l_{\mu ,\nu }^{(\mathrm{s},\mathrm{b})}`$, $`l_{\mu ,\nu }^{(\mathrm{s},\mathrm{i})}`$ and $`l_{\mu ,\nu }^{\mathrm{s},\mathrm{e}}`$ for the Scheidegger model. The data is for the distance between order $`\nu =3`$ side streams entering order $`\mu =6`$ absorbing streams. Once again, the distributions are well approximated by exponential distributions. Moreover, they are indistinguishable. This indicates, at least for the Scheidegger model, that drainage density is independent of relative position of tributaries along an absorbing stream.
We now consider the effect on the distribution of internal inter-tributary distances $`l_{\mu ,\nu }^{(\mathrm{s},\mathrm{i})}`$ following from variations in $`\mu `$, the order of the absorbing stream. Figure 8(a) provides a comparison of $`l_{\mu ,\nu }^{(\mathrm{s},\mathrm{i})}`$ distributions for $`\nu =2`$ and $`\mu =3`$ through $`\mu =8`$. As $`\mu `$ increases, the distributions tend towards a limiting function. With increasing $`\mu `$ we are, on average, sampling absorbing streams of greater length and the full range of $`l_{\mu ,\nu }^{(\mathrm{s},\mathrm{i})}`$ becomes accordingly more accessible. This approach to a fixed distribution is reflected in the means of the distributions in Figure 8(a). Shown in Figure 8(b), the means $`l_{\mu ,\nu }^{(\mathrm{s},\mathrm{i})}`$ for $`\nu =2`$ approach a value of around $`10.1`$. The corresponding density of second-order streams for the Scheidegger model is thus $`\rho _2=1/l_{\mu ,\nu =2}^{(\mathrm{s},\mathrm{i})}0.01`$. Higher drainage densities follow from equation (10). However, since deviations occur for small $`\nu `$, there will also be an approach to uniform scaling to consider with drainage density.
## VIII Joint variation of Tokunaga ratios and stream segment length
We have so far observed that the individual distributions of the $`l_\mu ^{(\mathrm{s})}`$ and $`T_{\mu ,\nu }`$ are exponential and that they are related via the side-stream density $`\rho _nu`$. However, this is not an exact relationship. For example, given a collection of stream segments with a fixed length $`l_\mu ^{(\mathrm{s})}`$ we expect to find fluctuations in the corresponding Tokunaga ratios $`T_{\mu ,\nu }`$.
To investigate this further we now consider the joint variation of $`T_{\mu ,\nu }`$ with $`l_\mu ^{(\mathrm{s})}`$ from a number of perspectives. After discussing the full joint probability distribution $`P(T_{\mu ,\nu },l_\mu ^{(\mathrm{s})})`$ we then focus on the quotient $`v=T_{\mu ,\nu }/l_\mu ^{(\mathrm{s})}`$ and its reciprocal $`w=l_\mu ^{(\mathrm{s})}/T_{\mu ,\nu }`$. The latter two quantities are measures of drainage density and inter-tributary length for an individual absorbing stream.
### VIII.1 The joint probability distribution
We build the joint distribution of $`P(T_{\mu ,\nu },l_\mu ^{(\mathrm{s})})`$ from our conception that stream segments are randomly distributed throughout a basin. In equation (18), we have the probability of a stream segment terminating after $`l_\mu ^{(\mathrm{s})}`$ steps. We need to incorporate into this form the probability that the stream segment also has $`T_{\mu ,\nu }`$ order $`\nu `$ side streams. Since we assume placement of these side streams to be random, we modify equation (18) to find
$$P(l_\mu ^{(\mathrm{s})},T_{\mu ,\nu })=\stackrel{~}{p}_\mu \left(\genfrac{}{}{0pt}{}{l_\mu ^{(\mathrm{s})}1}{T_{\mu ,\nu }}\right)p_\nu ^{T_{\mu ,\nu }}(1p_\nu \stackrel{~}{p}_\mu )^{l_\mu ^{(\mathrm{s})}T_{\mu ,\nu }1},$$
(24)
where $`\left(\genfrac{}{}{0pt}{}{n}{k}\right)=n!/k!(nk)!`$ is the binomial coefficient and $`p_\nu `$ is the probability of absorbing an order $`\nu `$ side stream. The extra $`p_\nu `$ appears in the last factor $`(1p_\nu \stackrel{~}{p}_\mu )`$ because this term is the probability that at a particular site the stream segment neither terminates nor absorbs an order $`\nu `$ side stream. Also, it is simple to verify that the sum over $`l_\mu ^{(\mathrm{s})}`$ and $`T_{\mu ,\nu }`$ of the probability in equation (19) returns unity.
While equation (19) does precisely describe the joint distribution $`P(l_\mu ^{(\mathrm{s})},T_{\mu ,\nu })`$, it is somewhat cumbersome to work with. We therefore find an analogous form defined for continuous rather than discrete variables. We simplify our notation by writing $`p=p_\nu `$, $`q=(1p_\nu \stackrel{~}{p}_\mu )`$ and $`\stackrel{~}{p}=\stackrel{~}{p}_\mu `$. We also replace $`(l_\mu ^{(\mathrm{s})},T_{\mu ,\nu })`$ by $`(x,y)`$ where now $`x,y`$. Note that $`0yx1`$ since the number of side streams cannot be greater than the number of sites within a stream segment.
Equation (19) becomes
$$P(x,y)=N\stackrel{~}{p}\frac{\mathrm{\Gamma }(x)}{\mathrm{\Gamma }(y+1)\mathrm{\Gamma }(xy)}(p)^y(q)^{xy1},$$
(25)
where we have used $`\mathrm{\Gamma }(z+1)=z!`$ to generalize the binomial coefficient. We have included the normalization $`N`$ to account for the fact that we have moved to continuous variables and the resulting probability may not be cleanly normalized. Also we must allow that $`N=N(p,\stackrel{~}{p})`$ and we will be able to identify this form more fully later on. Using Stirling’s approximation bender78 , that $`\mathrm{\Gamma }(z+1)\sqrt{2\pi }z^{z+1/2}e^z`$, we then have
$`P(x,y)`$ $`=`$ $`N\stackrel{~}{p}p^yq^{xy1}{\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle \frac{(x1)^{x3/2}}{y^{y+1/2}(xy1)^{xy1/2}}}`$ (26)
$`=`$ $`N{\displaystyle \frac{\stackrel{~}{p}}{\sqrt{2\pi }q}}p^yq^{xy}(x1)^{1/2}`$
$`\times \left({\displaystyle \frac{y}{x1}}\right)^{y1/2}\left(1{\displaystyle \frac{y}{x1}}\right)^{x+y+1/2}`$
$``$ $`N^{}x^{1/2}\left[F(y/x)\right]^x`$
where we have absorbed $`N`$ and all terms involving only $`p`$ and $`\stackrel{~}{p}`$ into the prefactor $`N^{}=N^{}(p,\stackrel{~}{p})=N\stackrel{~}{p}/(\sqrt{2\pi }q)`$. We have also assumed $`x`$ is large such that $`x1x`$ and $`11/x0`$.
The function $`F(v)=F(v;p,q)`$ identified above has the form
$$F(v)=\left(\frac{1v}{q}\right)^{(1v)}\left(\frac{v}{p}\right)^v.$$
(27)
where $`0<v<1`$ (here and later, the variable $`v`$ will refer to $`y/x`$). Note that for fixed $`x`$, the conditional probability $`P(y|x)`$ is proportional to $`[F(y/x)]^x`$. Figure 9(a) shows $`[F(v)]^x`$ for a range of powers $`x`$. The basic function has a single peak situated near $`v=p`$. For increasing $`x`$ which corresponds to increasing $`l_\mu ^{(\mathrm{s})}`$, the peak becomes sharper approaching (when normalized) a delta function, i.e., $`lim_x\mathrm{}[F(v)]^x=\delta (vp)`$.
Figure 9(b) provides a comparison between data for the Scheidegger model and the analytic form of $`P(l_\mu ^{(\mathrm{s})},T_{\mu ,\nu })`$. For this example, $`\mu =6`$ and $`\nu =2`$ which corresponds to $`p0.10`$, $`q0.90`$ and $`\stackrel{~}{p}0.001`$ (using the results of the previous section). The smooth curve shown is the conditional probability $`P(y|X)`$ for the example value of $`X=l_\mu ^{(\mathrm{s})}340`$ following from equation (26). From simulations, we obtain a discretized approximation to $`P(l_\mu ^{(\mathrm{s})},T_{\mu ,\nu })`$. For each fixed $`x=l^{(\mathrm{s})}`$ in the range $`165l_\mu ^{(\mathrm{s})}345`$, we rescale the data using the following derived from equation (26),
$$P(X,y)=N^{}X^{1/2}\left(N_{}^{}{}_{}{}^{1}x^{1/2}P(x,y)\right)^{X/x}.$$
(28)
All rescaled data is then combined, binned and plotted as circles in Figure 9(b), showing excellent agreement with the theoretical curve.
### VIII.2 Distributions of side branches per unit stream length
Having obtained the general form of $`P(l_\mu ^{(\mathrm{s})},T_{\mu ,\nu })`$, we now delve further into its properties by investigating the distributions of the ratio $`v=T_{\mu ,\nu }/l_\mu ^{(\mathrm{s})}`$ and its reciprocal $`w`$.
The quantity $`T_{\mu ,\nu }/l_\mu ^{(\mathrm{s})}`$ is the number of side streams per length of a given absorbing stream and when averaged over an ensemble of absorbing streams gives
$$T_{\mu ,\nu }/l_\mu ^{(\mathrm{s})}=\rho _\nu .$$
(29)
Accordingly, the reciprocal $`l_\mu ^{(\mathrm{s})}/T_{\mu ,\nu }`$ is the average separation of side streams of order $`\nu `$.
First, we derive $`P(T_{\mu ,\nu }/l^{(\mathrm{s})})`$ from $`P(l_\mu ^{(\mathrm{s})},T_{\mu ,\nu })`$. We then consider some intuitive rescalings which will allow us to deduce the form of the normalization $`N(p,q)`$.
We rewrite equation (26) as
$$P(x,y)=N^{}x^{1/2}\mathrm{exp}\left\{x\mathrm{ln}\left[F(y/x)\right]\right\}.$$
(30)
We transform $`(x,y)`$ to the modified polar coordinate system described by $`(u,v)`$ with the relations
$$u^2=x^2+y^2\text{and}v=y/x.$$
(31)
The inverse relations are $`x=u/(1+v^2)`$ and $`y=uv/(1+v^2)`$ and we also have $`\text{d}x\text{d}y=x\text{d}u\text{d}v`$. Equation (32) leads to
$$P(u,v)=N^{}\left(\frac{u}{1+v^2}\right)^{1/2}\mathrm{exp}\left\{\frac{u}{1+v^2}\mathrm{ln}\left[F(v)\right]\right\}.$$
(32)
To find $`P(v)`$ we integrate out over the radial dimension $`u`$:
$`P(v)`$ $`=`$ $`{\displaystyle _{u=0}^{\mathrm{}}}\text{d}uP(u,v)`$ (33)
$`=`$ $`N^{}{\displaystyle _{u=0}^{\mathrm{}}}\text{d}u\left({\displaystyle \frac{u}{1+v^2}}\right)^{1/2}\mathrm{exp}\left\{{\displaystyle \frac{u}{1+v^2}}\mathrm{ln}\left[F(v)\right]\right\}`$
$`=`$ $`N^{}(1+v^2)(\mathrm{ln}[F(v)])^{3/2}{\displaystyle _{z=0}^{\mathrm{}}}\text{d}zz^{1/2}e^z`$
$`=`$ $`N^{\prime \prime }{\displaystyle \frac{1+v^2}{(\mathrm{ln}[F(v)])^{3/2}}}.`$
Here, $`N^{\prime \prime }=N^{}\mathrm{\Gamma }(3/2)=N^{}\sqrt{\pi }/2`$ and we have used the substitution $`z=u/(1+v^2)\mathrm{ln}[F(v)]`$.
The distribution for $`w=l_\mu ^{(\mathrm{s})}/T_{\mu ,\nu }=1/v`$ follows simply from equation (33) and we find
$$P(w)=N^{\prime \prime }\frac{1+w^2}{w^4(\mathrm{ln}[F(1/w)])^{3/2}}.$$
(34)
Figures 10(a) and 10(b) compare the predicted forms of $`P(v)`$ and $`P(w)`$ with data from the Scheidegger model. In both cases, the data is for order $`\nu =2`$ side streams being absorbed by streams of order $`\mu =6`$. Note that both distributions show an initially exponential-like decay away from a central peak. Moreover, the agreement is excellent, offering further support to the notion that the spatial distribution of stream segments is random.
Finally, we quantify how changes in the orders $`\mu `$ and $`\nu `$ affect the width of the distributions by considering some natural rescalings. Figure 11(a) shows binned, normalized distributions of $`T_{\mu ,\nu }/l_\mu ^{(\mathrm{s})}`$ for the Scheidegger model. Here, the side stream order is $`\nu =2`$ and the absorbing stream orders range over $`\mu =5`$ to $`\mu =8`$. All distributions are centered around $`\rho _20.10`$.
Because the average length of $`l_\mu ^{(\mathrm{s})}`$ increases by a factor $`R_{l^{(\mathrm{s})}}`$ with $`\mu `$, the typical number of side streams increases by the same factor. Since we can decompose $`l_\mu ^{(\mathrm{s})}`$ as
$$l_\mu ^{(\mathrm{s})}=l_{\mu ,\nu }^{(\mathrm{s},\mathrm{b})}+l_{\mu ,\nu }^{(\mathrm{s},\mathrm{i})}+\mathrm{}+l_{\mu ,\nu }^{(\mathrm{s},\mathrm{i})}+l_{\mu ,\nu }^{\mathrm{s},\mathrm{e}},$$
(35)
where there are $`T_{\mu ,\nu }1`$ instances of $`l_{\mu ,\nu }^{(\mathrm{s},\mathrm{i})}`$, $`l_\mu ^{(\mathrm{s})}`$ becomes better and better approximated by $`(T_{\mu ,\nu }+1)l_{\mu ,\nu }^{(\mathrm{s},\mathrm{i})}`$.
Hence, the distribution of $`T_{\mu ,\nu }/l_\mu ^{(\mathrm{s})}`$ peaks up around $`\rho _2`$ as $`\mu `$ increases, the typical width reducing by a factor of $`1/\sqrt{R_{l^{(\mathrm{s})}}}`$ for every step in $`\mu `$. Using this observation, Figure 11(b) shows a rescaling of the same distributions shown in Figure 11(a). The form of this rescaling is
$$P(T_{\mu ,\nu }/l_\nu ^{(\mathrm{s})})=(R_{l^{(\mathrm{s})}})^{\mu /2}G_1\left([T_{\mu ,\nu }/l_\nu ^{(\mathrm{s})}\rho _2](R_{l^{(\mathrm{s})}})^{\mu /2}\right)$$
(36)
where the function is similar to the form of $`P(v)`$ given in equation (33). The mean drainage density of $`\rho _2`$ has been subtracted to center the distribution.
We are able to generalize this scaling form of the distribution further by taking into account side stream order. Figures 12(a) and 12(b) respectively show the unrescaled and rescaled distributions of $`T_{\mu ,\nu }/l_\mu ^{(\mathrm{s})}`$ with $`\nu `$ allowed to vary. This particular example taken from the Scheidegger model is for $`\mu =6`$ and the range $`\nu =1`$ to $`\nu =5`$. Since $`\nu `$ is now changing, the centers are situated at the separate values of the $`\rho _\nu `$. Also, the typical number of side streams changes with order $`\nu `$ so the widths of the distributions dilate as for the varying $`\mu `$ case by a factor $`\sqrt{R_{l^{(\mathrm{s})}}}`$. Notice that the rescaling works well for $`\nu =2,\mathrm{},5`$ but not $`\nu =1`$. As we have noted, deviations from scaling from small orders are to be expected. In this case, we are led to write down
$$P(T_{\mu ,\nu }/l_\nu ^{(\mathrm{s})})(R_{l^{(\mathrm{s})}})^{\nu /2}G_2\left([T_{\mu ,\nu }/l_\nu ^{(\mathrm{s})}\rho _\nu ](R_{l^{(\mathrm{s})}})^{\nu /2}\right)$$
(37)
where, again, $`G_2(z)`$ is similar in form to $`P(v)`$.
We find the same rescalings apply for the Mississippi data. For example, Figure 13(a) shows unrescaled distributions of $`T_{\mu ,\nu }/l_\nu ^{(\mathrm{s})}`$ for varying $`\nu `$. Figure 13(b) then shows reasonable agreement with the form of equation (37). In this case, the Scheidegger model clearly affords valuable guidance in our investigations of real river networks. The ratio $`R_{l^{(\mathrm{s})}}=2.40`$ was calculated from an analysis of $`l_\omega ^{(\mathrm{s})}`$ and $`l_\omega `$. The density $`\rho _20.0004`$ was estimated directly from the distributions of $`T_{\mu ,\nu }/l_\nu ^{(\mathrm{s})}`$ and means that approximately four second-order streams appear every ten kilometers.
Combining equations (VIII.2) and (VIII.2), we obtain the complete scaling form
$`P(T_{\mu ,\nu }/l_\nu ^{(\mathrm{s})})`$
$`=`$ $`(R_{l^{(\mathrm{s})}})^{(\mu \nu 1)/2}G\left([T_{\mu ,\nu }/l_\nu ^{(\mathrm{s})}\rho _\nu ](R_{l^{(\mathrm{s})}})^{(\mu \nu 1)/2}\right).`$
As per $`G_1`$ and $`G_2`$, the function $`G`$ is similar in form to $`P(v)`$.
The above scaling form makes intuitive sense but is not obviously obtained from an inspection of (33). We therefore examine $`P(v)`$ by determining the position and magnitude of its maximum. Rather than solve $`P^{}(v)=0`$ directly, we find an approximate solution by considering the argument of the denominator, $`\mathrm{ln}F(v)`$, with $`F(v)`$ given in equation (27). Since the numerator of $`P(v)`$ is $`1+v^2`$ and the maximum occurs for small $`v`$ this is a justifiable step. Setting $`\text{d}F/\text{d}v=0`$, we thus have
$$\mathrm{ln}\frac{1v}{q}+\mathrm{ln}vp=0,$$
(39)
which gives $`v_m=p/(q+p)=p/(1\stackrel{~}{p})`$. Note that for $`\stackrel{~}{p}1`$, we have $`v_mp`$.
Substituting $`v=v_m=p/(1\stackrel{~}{p}`$ into equation (33), we find
$$P(v_m)N^{\prime \prime }\stackrel{~}{p}^{3/2}=N\stackrel{~}{p}^{1/2}2^{3/2}$$
(40)
presuming $`p^21`$ and $`q1`$. Returning to the scaling form of equation (VIII.2), we see that the $`\stackrel{~}{p}^{1/2}`$ factor in equation (40) accounts for the factors of $`(R_{l^{(\mathrm{s})}})^{\mu /2}`$ since $`\stackrel{~}{p}=\stackrel{~}{p_\mu }`$ scales from level to level by the ratio $`R_{l^{(\mathrm{s})}}`$. We therefore find the other factor $`(R_{l^{(\mathrm{s})}})^{\nu /2}`$ of equation (VIII.2) gives $`N=cp^{1/2}`$ where $`c`$ is a constant. Since $`p=p_\nu `$, it is the only factor that can provide this variation. We thus have found the variation with stream order of the normalization $`N`$ and have fully characterized, $`P(x,y)`$, the continuum approximation of $`P(l_\mu ^{(\mathrm{s})},T_{\mu ,\nu })`$.
## IX Concluding remarks
We have extensively investigated river network architecture as viewed in planform. We identify the self-similarity of a form of drainage density as the essence of the average connectivity and structure of networks. From previous work in dodds99pa , we then understand this to be a base from which all river network scaling laws may obtained.
We have extended the description of tributary structure provided by Tokunaga’s law to find that side stream numbers are distributed exponentially. This in turn is seen to follow from the fact that the length of stream segments are themselves exponentially distributed. We interpret this to be consequence of randomness in the spatial distribution of stream segments. Furthermore, the presence of exponential distributions indicate fluctuations in variables are significant being on the order of mean values. For the example of stream segment lengths, we thus identify $`\xi _{l^{(\mathrm{s})}}`$, a single parameter needed to describe all moments. This is simply related to $`\xi _t`$, which describes the distributions of Tokunaga ratios. The exponential distribution becomes the null hypothesis for the distributions of these variables to be used in the examination of real river networks.
We are able to discern the finer details of the connection between stream segment length and tributary numbers. Analysis of the placement of side streams along a stream segment again reveals exponential distributions. We are then able to postulate a joint probability distribution for stream segment lengths and the Tokunaga ratios. The functional form obtained agrees well with both model and real network data. By further considering distributions of the number of side streams per unit length of individual stream segments, we are able to capture how variations in the separation of side streams are averaged out along higher-order absorbing streams.
By expanding our knowledge of the underlying distributions through empiricism, modeling and theory, we obtain a more detailed picture of network structure with which to compare real and theoretical networks. We have also further shown that the simple random network model of Scheidegger has an impressive ability to produce statistics whose form may then be observed in nature. Indeed, the only distinction between the two is the exact value of the scaling exponents and ratios involved since all distributions match up in functional form.
We end with a brief comment on the work of Cui et al. cui99 who have recently also proposed a stochastic generalization of Tokunaga’s law. They postulate that the underlying distribution for the $`T_{\mu ,\nu }`$ is a negative binomial distribution. One parameter additional to $`T_1`$ and $`R_T`$, $`\alpha `$, was introduced to reflect “regional variability,” i.e., statistical fluctuations in network structure. This is in the same spirit as our identification of a single parameter $`\xi _t`$. However, our work disagrees on the nature of the underlying distribution of $`T_{\mu ,\nu }`$. We have consistently observed exponential distributions for $`T_{\mu ,\nu }`$ in both model and real networks.
In closing, by finding randomness in the spatial distribution of stream segments, we have arrived at the most basic description of river network architecture. Understanding the origin of the exact values of quantities such as drainage density remains an open problem.
## Acknowledgements
The authors would like to thank J.S. Weitz for useful discussions. This work was supported in part by NSF grant EAR-9706220 and the Department of Energy grant DE FG02-99ER 15004.
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# Discovery of 424 ms pulsations from the radio-quiet neutron star in the PKS 1209–52 supernova remnant
## 1 Introduction
Radio-quiet compact central objects (CCOs) of supernova remnants (SNRs) have emerged recently as a separate class of X-ray sources (see, e.g., Kaspi 2000, for a review). The nature of these sources, which are expected to be either neutron stars (NSs) or black holes formed in supernova explosions, remains enigmatic. They are characterized by soft, apparently thermal, X-ray spectra and a lack of visible pulsar activity and optical counterparts. One of the most important properties, which potentially allows one to elucidate the nature of these sources, is periodicity of their radiation. Some of these sources show periods in the range of 6–12 s, much longer than those of typical radio pulsars. These object form a subclass of anomalous X-ray pulsars (e.g., Gotthelf & Vasisht 2000), which have been interpreted as magnetars (NSs with superstrong magnetic fields — see Thompson 2000). A putative period of 75 ms was proposed by Pavlov, Zavlin, & Trümper (1999) for the Puppis A CCO; however, a low significance of that detection requires it to be confirmed by other observations. An unusually long period of 6 hours has been reported for the RCW 103 CCO (Garmire et al. 2000a), the origin of this periodicity is unclear. No periodicity has been found for the recently discovered CCO of the Cas A SNR (Pavlov et al. 2000; Chakrabarty et al. 2000).
In this paper we report detection of a period of 1E 1207.4–5209, the CCO of PKS 1209–52, a barrel-shaped radio, X-ray, and optical SNR (also known as G296.5+10.0). From the analysis of radio and optical observations of this SNR, Roger et al. (1988) estimated its age $`7000`$ yr, with an uncertainty of a factor of 3. A recent estimate of the distance to PKS 1209–52, $`d=2.1_{0.8}^{+1.8}`$ kpc, is given by Giacani et al. (2000). Estimates of the interstellar hydrogen column density from the radio, optical, and UV data yield $`n_{\mathrm{H},21}n_H/(10^{21}\mathrm{cm}^2)1.0`$–1.8 (see Kellet et al. 1987; Roger et al. 1988; Giacani et al. 2000), consistent with a distance $`d_2d/(2\mathrm{kpc})1`$.
The point source 1E 1207.4–5209 was discovered with the Einstein observatory (Helfand & Becker 1984), $`6^{}`$ off-center the $`81^{}`$ diameter SNR. Mereghetti, Bignami & Caraveo (1996) and Vasisht et al. (1997) showed that the ROSAT and ASCA spectra of 1E 1207.4–5209 can be interpreted as blackbody emission of $`T3`$ MK from an area with radius $`R1.5d_2`$ km and suggested that this radiation comes from hot spots on NS surface. The spots may be heated either by dissipative heating in the NS interior or by the bombardment of polar caps by relativistic particles from the NS magnetosphere if 1E 1207.4–5209 is an active pulsar. However, the former hypothesis can hardly explain the small sizes of the hot spots, even with allowance for large anisotropy of thermal conductivity of the magnetized NS crust. The latter heating mechanism is also in doubt because of the absence of radio and $`\gamma `$-ray emission from 1E 1207.4–5209.
From observations at 4.8 GHz, Mereghetti et al. (1996) found an upper limit of $`0.1`$ mJy on the radio flux from 1E 1207.4–5209. They also set a deep limit of $`V>25`$ for an optical counterpart in the Einstein HRI error circle, which supports the hypothesis that 1E 1207.4–5209 is indeed an isolated NS.
Zavlin, Pavlov & Trümper (1998) reanalyzed the ROSAT and ASCA data, fitting the observed spectra with NS atmosphere models. They have shown that the hydrogen atmosphere fits yield more realistic parameters of the NS and the intervening hydrogen column than the traditional blackbody fit. In particular, for a NS of mass $`1.4M_{}`$ and radius 10 km, they obtained a NS surface temperature $`T_{\mathrm{eff}}=(1.4`$$`1.9)`$ MK and distance $`d=1.6`$–3.3 kpc, versus $`T=(4.2`$$`4.6)`$ MK and implausibly large $`d=11`$–13 kpc for the blackbody fit, at a 90% confidence level. The hydrogen column density inferred from the atmosphere fits, $`n_{\mathrm{H},21}=0.7`$–2.2, agrees fairly well with independent estimates obtained from UV observations of nearby stars, radio data, and X-ray spectrum of the shell of the supernova remnant, whereas the blackbody and power-law fits give considerably lower and greater values, $`n_{\mathrm{H},21}=0.2`$–0.4 and $`n_{\mathrm{H},21}=5.2`$–7.0, respectively. The NS surface temperature inferred from the atmosphere fits is consistent with standard NS cooling models.
All the previous observations failed to detect pulsations of X-ray radiation from 1E 1207.4–5209. An upper limit of 18% (at a 95% confidence level) on flux modulation was set by Mereghetti et al. (1996) from the analysis of the ROSAT PSPC data. Due to its higher sensitivity and the possibility of continuous observations, the Chandra X-ray Observatory is much more capable to search for periodicity of X-ray sources. We have employed this capability to search for the period of 1E 1207.4–5209. The observation and data reduction are described in §2, the timing analysis of the data is presented in §3, and some implementations are briefly discussed in §4.
## 2 Observation and data reduction
1E 1207.4–5209 was observed with Chandra on 2000 January 6-7 with the spectroscopic array of the Advanced CCD Imaging Spectrometer (ACIS — see Garmire et al. 2000b) in the Continuous Clocking (CC) mode. This mode provides the highest time resolution of 2.85 ms available with ACIS by means of sacrificing spatial resolution in one dimension. The source was imaged on the back-illuminated chip S3. The total duration of the observation was 32.6 ks. Time history of detected events reveals that there were three relatively short time intervals with very strongly, by an order of magnitude, increased background (background “flares”) distributed over the whole detector. To mitigate the contamination of the source by the flares, we excluded these intervals from further analysis, which resulted in the effective exposure of 29,283 s.
The event arrival times were not properly corrected for the satellite wobbling (dither) by the standard pipe-line processing. From our analysis (Sanwal et al. 2000) of Chandra data on PSR 1055–52, observed in the same mode, we have been aware that the lack of this correction results in false periodicities (side peaks in the power spectrum) of a pulsar. Therefore, we corrected the event arrival times $`t`$ in the event file making use of the formula (Glenn Allen, private communication):
$$t_{\mathrm{corr}}=t+c_1\left[(\alpha \alpha _\mathrm{m})\mathrm{sin}\zeta \mathrm{cos}\delta _\mathrm{m}(\delta \delta _\mathrm{m})\mathrm{cos}\zeta \right],$$
(1)
where $`c_1=c_2\delta t`$, $`c_2=3600^{\prime \prime }/0\stackrel{}{\mathrm{.}}4919=7318.5605`$ is the scaling coefficient transforming degrees to the detector pixels, $`\delta t=2.85`$ ms is the integration time in the CC mode, $`\zeta `$ is the roll angle, and $`\alpha \alpha _m`$ and $`\delta \delta _m`$ (in degrees) are the deviations of right ascension and declination from their median values (calculated with the use of the aspect solution file). We have checked that this correction removes the artificial periodicities in the data on PSR 1055–52.
The 1D image of 1E 1207.4–5209 in “sky pixels” (Fig. 1) is consistent with the ACIS PSF. To obtain the source count rate, we extracted 23,337 source+background counts from a 1D segment of 8 pixel length centered at the source position (10 pixels contain 23,602 counts). The background was taken from similar segments adjacent to the 1D source aperture. Subtracting the background, we find the source countrate of $`0.76\pm 0.01`$ s<sup>-1</sup>. A preliminary analysis of the source spectrum allows us to estimate the observed source energy flux $`f_x=2.2\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, in the 0.1–5.0 keV range. The corresponding (unabsorbed) luminosity is $`L_x1.3\times 10^{33}d_2^2`$ erg s<sup>-1</sup>, and the bolometric luminosity (corrected for the gravitational redshift at the surface of NS with $`M=1.4M_{}`$, $`R=10`$ km) is $`L_x5\times 10^{33}d_2^2`$ erg s<sup>-1</sup>, in agreement with the estimates obtained by Zavlin et al. (1998).
## 3 Timing analysis
Previous $`ROSAT`$ and $`ASCA`$ observations of 1E 1207.4–5209 have not revealed pulsations of X-ray radiation from this object. This may be explained by small amount of counts collected in those observations. Moreover, those data were spread over very long time intervals, that always heavily complicates searching for weak pulsations.
One of the main advantages of the Chandra observation is that a large number of the source counts (about 10 times the number of counts previously detected) has been acquired in a short time span. The spectral analysis shows that the background exceeds the source radiation at energies above 3 keV. Therefore, for the timing analysis we chose 22,535 counts with energies below 3 keV from a 6-pixel segment centered at the 1E 1207.4–5209 position. Of these counts, about 98% were estimated to belong to the source. To search for pulsations, we used the well-known $`Z_1^2`$ (Rayleigh) test (Buccheri et al. 1983), which is expected to be optimal to search for smooth pulsations. We ran the test in the 0.01–100 Hz frequency range with a step $`\delta f=3`$ $`\mu `$Hz. The oversampling by a factor of 10 (compared to the expected widths of $`1/T`$ of the $`Z_1^2`$ peaks, where $`T=32.6`$ ks is the observational span) was chosen to resolve separate $`Z_1^2`$ peaks; it guarantees that the peak corresponding to a periodic signal would not be missed. The number of statistically independent trials in the chosen frequency range can be estimated as $`𝒩=100\mathrm{Hz}\times T=3.26\times 10^6`$. The test yields only one high peak of $`Z_{1,\mathrm{max}}^2=65.0`$ at the frequency $`f_{}=2.357769`$ Hz (see Fig. 2). Since the variable $`Z_1^2`$ has a probability density function equal to that of a $`\chi ^2`$ with two degrees of freedom, the probability to obtain a noise peak of a given height in one trial is $`\mathrm{exp}(Z_1^2/2)`$. This means that the probability to obtain by chance a peak of $`Z_1^2=65.0`$ in $`𝒩`$ independent trials is $`\rho =𝒩\mathrm{exp}(Z_{1,\mathrm{max}}^2/2)=2.5\times 10^8`$, that corresponds to a detection of pulsations at a confidence level of $`C=(1\rho )\times 100\%=99.9999975\%`$, or $`5.5\sigma `$. The next heighest peak is $`Z_1^2=32.0`$. The probability to obtain such a peak by chance in $`𝒩`$ trials is 35%.
To be sure that the pulsations found at the frequency $`f_{}`$ are not associated with instrumental effects, we ran the $`Z_1^2`$ test at $`f=f_{}`$ on the same amounts of counts extracted from a few background regions and obtained maximum $`Z_1^2`$ equal to 5. We also varied the size of the segment for source counts extraction and found that the initial choice of 6-pixel length was optimal to produce the maximum $`Z_1^2`$ value. We found that contributions from higher harmonics to the value of $`Z^2`$ are insignificant: $`Z_2^2=65.1`$ and $`Z_3^2=65.2`$.
To evaluate the pulsation frequency more precisely and find its uncertainty, we employed the method suggested by Gregory & Loredo (1996; GL hereafter), based on the Bayesian formalism. The method uses the phase-averaged epoch-folding algorithm to calculate a frequency-dependent odds ratio, $`O_m(f)`$, which specifies how the data favor a periodic model of a given frequency $`f`$ with $`m`$ phase bins over the unpulsed model (see Gregory & Loredo 1992 for details of odds ratio computations). To weaken the dependence on number of bins, GL suggested to use the odds ratio $`O_{\mathrm{per}}(f)=_{m=2}^{m_{\mathrm{max}}}O_m(f)`$, with a characteristic number of $`m_{\mathrm{max}}=10`$–15. The distribution of probability for a signal to be periodic, with frequency $`f`$ in a chosen frequency range $`(f_1,f_2)`$, can be written as
$$p(f)=\frac{O_{\mathrm{per}}(f)}{f}\left[(1+O_{\mathrm{per}}^{})\mathrm{ln}\frac{f_2}{f_1}\right]^1,$$
(2)
where
$$O_{\mathrm{per}}^{}=\left(\mathrm{ln}\frac{f_2}{f_1}\right)^1_{f_1}^{f_2}\frac{\mathrm{d}f}{f}O_{\mathrm{per}}(f),$$
(3)
and the probability for the signal to be periodic in $`(f_1,f_2)`$ is $`O_{\mathrm{per}}^{}/(1+O_{\mathrm{per}}^{})`$. Using PSR 0540–69 as an example, GL have demonstrated that this method allows one to determine the pulsation frequency with much higher accuracy than it can be done with more traditional methods. The most probable frequency $`f_0`$ is given by the maximum of $`O_{\mathrm{per}}(f)`$ in $`(f_1,f_2)`$. The uncertainty $`\delta f`$ at a given confidence level $`C=\stackrel{~}{\rho }\times 100\%`$ can be calculated from the following equation
$$_{f_0\delta f/2}^{f_0+\delta f/2}p(f)df=\stackrel{~}{\rho }.$$
(4)
We implemented the GL method, taking $`m_{\mathrm{max}}=12`$ and $`f_2f_{}=f_{}f_1=200`$ $`\mu `$Hz. The frequency dependence of odds ratio is shown in Figure 2. The maximum value, $`O_{\mathrm{per}}^{\mathrm{max}}=6.2\times 10^8`$, is at $`f=f_0=2.3577717`$ Hz ($`f_0f_{}3\mu \mathrm{Hz}0.1/T`$). The uncertainty of $`f_0`$ at 68%, 90%, and 95% confidence levels is $`\delta f=1.3`$, 2.5, and 3.5 $`\mu `$Hz, respectively. This estimate is practically independent of choice of frequency range if $`\mathrm{\Gamma }f_2f_1f_{}`$, where $`\mathrm{\Gamma }3`$ $`\mu `$Hz is a characteristic width of the peak of the odds ratio (see Fig. 2). Thus, we finally derive the frequency and the period of the detected pulsations
$$f_0=2.3577717\pm 1.3\times 10^6\mathrm{Hz},$$
(5)
$$P_0=0.42412927\pm 2.3\times 10^7\mathrm{s},$$
(6)
at the epoch of 51549.630051303 MJD (TDB).
The light curve extracted at $`f=f_0`$ (Fig. 3) reveals one broad pulse per period with intrinsic source pulsed fraction of $`f_p=9\pm 2\%`$. Assuming that the detected signal is sinusoidal, we obtain an estimate on the pulsed fraction $`f_p=(2Z_{1,\mathrm{max}}^2/N)^{1/2}=7.6\%`$ (where $`N`$ is the number of counts), in a fair agreement with the value calculated from the extracted light curve. Figure 3 indicates that the shape of the light curve may vary slightly with photon energy.
## 4 Discussion
The detection of the period, $`P424`$ ms, proves that 1E 1207.4–5209, the central compact object of the PKS 1209–52 SNR, is the neutron star.
Let us assume that 1E 1207.4–5209 is an active pulsar of an age $`\tau `$ within a range of 2–20 kyr (estimated limits of the SNR age). Then, for a braking index of 2.5, we should expect $`\dot{f}f/(1.5\tau )(2.5`$$`25)\times 10^{12}`$ Hz s<sup>-1</sup>, $`\dot{P}=(4.5`$$`45)\times 10^{13}`$ s s<sup>-1</sup>. A crude estimate for the pulsar magnetic field (at the magnetic pole) would be $`B6.4\times 10^{19}(P\dot{P})^{1/2}(3`$$`9)\times 10^{13}`$ G, close to the critical value $`B_{\mathrm{cr}}=4.4\times 10^{13}`$ G, above which nonlinear QED effects can affect the processes in the pulsar magnetosphere. The pulsar spin-down luminosity (rotation energy loss) would be $`\dot{E}=(2.3`$$`23)\times 10^{35}I_{45}`$ erg s<sup>-1</sup>. This $`\dot{E}`$ may be high enough to power a compact synchrotron nebula — for instance, Gaensler et al. (1998) found a radio nebula with a radius of $`0.3(d/7\mathrm{kpc})`$ pc around a 100-kyr old PSR B0906–49, whose spin-down luminosity is $`4.9\times 10^{35}`$ erg s<sup>-1</sup>. A nebula of similar size would have an angular radius of $`30^{\prime \prime }`$ at a distance of 2 kpc. The 1D image of 1E 1207.4–5209 puts an upper limit of $`3^{\prime \prime }`$$`4^{\prime \prime }`$ on the radius of a compact nebula around the pulsar.
If 1E 1207.4–5209 is an ordinary radio pulsar whose radio-quiet nature is due to an unfavorable orientation of the pulsar beam, the pulsar could be detected in deep radio observations at low frequencies, where radio beams are broader. On the other hand, the lack of manifestations of pulsar activity may indicate that 1E 1207.4–5209 is not an active pulsar, e.g., because a very high magnetic field can inhibit the cascade processes in the pulsar’s acceleration zone (Baring & Harding 1997). In this case, the most natural explanation of the observed pulsations would be anisotropy of temperature distribution caused by anisotropic heat conduction in a superstrong magnetic field (Greenstein & Hartke 1983). This hypothesis can be verified by a phase-dependent spectral analysis which will be done on the same data when the ACIS response is known with better precision. Particularly useful for elucidating the nature of the NS would be measuring of the period derivative, which requires at least one more observation. For instance, a similar ACIS observation in CC mode, taken a year later, would allow us to determine the frequency with an accuracy of $`2\times 10^6`$ Hz and detect a frequency derivative if it exceeds $`2\times 10^{13}`$ Hz s<sup>-1</sup>, which is substantially smaller than the above estimates.
Our thanks are due to Glenn Allen, who provided the algorithm to correct the event times for dither. We are grateful to Gordon Garmire, John Nousek, Leisa Townsley and George Chartas for the useful advice on the analysis of ACIS data. This work was partly supported by SAO grant GO0-1012X.
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# LWS spectroscopy of the luminous Blue Compact Galaxy Haro 11Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA.
## 1 Introduction
Blue compact galaxies (BCGs) are characterized by their high star formation rate and low chemical abundances. The high ratio between blue luminosity and HI mass, L<sub>B</sub>/$``$<sub>HI</sub>, combined with a high relative HI mass fraction, indicates a high star formation efficiency and a high gas consumption rate. The apparent conflict between intense star formation and low chemical abundances can only be explained if the gas containing the recently produced elements is expelled from the galaxy, if fresh gas is accreted onto the galaxy or if the galaxy is young (e.g. Searle & Sargent searle (1972); Makarova & Karachentsev makarova (1998); Papaderos et al. papaderos (1998)). We still have rather vague ideas about what is triggering the intense global burst. Our previous observations (Östlin et al. 1999a , 1999b ) indicate that gas dynamical instabilities in connection with mergers could be the key mechanism. Under such circumstances the heating processes of the ISM could become a complicated mixture of shocks caused by infalling clouds and winds from massive young stars and on the other hand a strong radiation field from the starburst.
While optical and radio observations give us important information about the conditions in the hot ionized and the cold neutral gas in star forming galaxies we often lack information about the the ISM in the transition zone, the photodissociation regions (PDRs). The gas in PDRs contains neutral atoms and ions of low ionization potential ($`\chi <`$ 13.6 eV), e.g. C<sup>+</sup>. With the different possible heating mechanisms involved one would expect that the geometry and spatial extension of the ISM and the PDRs would be quite different in galaxies like Haro 11 and those with quiet star formation activity or with active nuclei, resulting in radically different spectral signatures.
Spectral diagnostics are valuable tools in dissecting the balance between different heating and cooling processes and the physical conditions in the gas. The first observations of \[OI\]63$`\mu `$ in M82 was presented by Watson et al. (1972b ). Crawford et al. (crawford (1985)) for the first time showed that important information about the PDR physics in external galaxies can be obtained from the collisionally excited \[CII\]$`\lambda `$158$`\mu `$. Later, several studies of star forming regions in the Milky Way as well as other starburst galaxies and ultraluminous infrared galaxies (’ULIGs’, possibly hiding active nuclei) have demonstrated the power of more detailed far-IR spectroscopy (e.g. Colbert et al. colbert (1999)). An advantage with observations in the IR region is that we diminish the influence of dust extinction, allowing a more direct comparison between ’naked’ starbursts in low mass galaxies and hidden starbursts in ULIGs.
The fact that \[CII\]$`\lambda `$158$`\mu `$ emission constitutes the peak of the flux distribution f<sub>ν</sub> in starforming galaxies has led to a discussion about the possibility of using \[CII\] to detect young galaxies at high redshifts (e.g. Petrosian et al. petrosian (1969) and the detailed discussion by Stark stark (1997)). Despite some emerging pessimism about this possibility (Gerin & Phillips gerin (1998)) we will argue that this still is a viable idea.
As part of our multifrequency study of three of the most luminous BCGs in the southern hemisphere this report will focus on observations using the Long Wavelength Spectrometer (LWS) (Clegg et al. clegg (1996)) on the Infrared Space Observatory (ISO) (Kessler et al. kessler (1996)) of the luminous BCG Haro 11. Only a few emission lines in the LWS window are detectable in most nearby BCGs in a reasonably long integration time. In our case we chose to obtain clear detection of a few of these lines, with a sufficiently high S/N to allow comparisons with models and other galaxies. These lines are the \[OI\]63$`\mu `$, \[OIII\]88$`\mu `$ and \[CII\]158$`\mu `$ lines, which constitute important coolants in the far-IR regions. Among these lines, \[CII\] is of particular importance for reasons just mentioned and discussed below. In addition to these three lines, we obtained an upper limit on the flux in the \[OI\]145$`\mu `$ line in one of the adjoining detectors.
## 2 Observations and reductions
The observations were carried out during the routine observation phase of ISO in the low-resolution LWS02 observation mode with a resolution of 0.6 $`\mu `$. We used a scan width of 7 resolution elements, 8 samples per element and fast scanning. Each spectrum was scanned 6-111 times with a nominal 0.50 s integration time giving in total an effective integration time of 2820 s.
The data presented in this paper have been taken from pipeline 6.1. All the glitches in each individual scan have been removed by using the Interactive Analysis Software Package. About 15% of the total scans were rejected due to heavy cosmic rays contamination. In addition to the requested data the LWS spectrometer provided data on all the other detectors. These data have been also analyzed and an upper limit on the fine structure line of \[OI\]$`\lambda `$145$`\mu `$ has been derived. The resulting data were resampled and coadded. The data are shown in figure 1.
Exept for the LWS data we also present some additional previously unpublished data in Table 1. The CO observations were carried out during June 9-12 1988 with the SEST telescope equipped with cooled Schottky receivers at 115 GHz (Jörsäter & Bergvall steven (1999)). Six blue compact galaxies were observed and none was detected.
Optical spectroscopy was obtained with ESOs 1.5-m telescope in 1983, using the Image Dissector Scanner and an aperture of 4”x4” arcseconds<sup>2</sup> (Bergvall & Olofsson bergvall1 (1986)). Additional data were obtained with the IDS at the ESO 3.6-m telescope in 1986. Standard reduction techniques were used. More information is found in Bergvall & Östtlin (1999a ).
## 3 Results
### 3.1 General properties of the target galaxy
Haro 11 (ESO 338-IG04) is a galaxy with extreme starburst properties, summarized in Table 1. In the following we will assume a Hubble constant of H<sub>0</sub> = 65 km s<sup>-1</sup>Mpc<sup>-1</sup>. Most of the data are from our previous observations in the UV-radio region (Bergvall & Olofsson bergvall1 (1986); Östlin et al 1999a , 1999b ; Bergvall & Östlin 1999a and other unpublished data). Our observations show three conspicious condensations in the centre embedded in a regular H$`\alpha `$ halo. HST images (Malkan et al. malkan (1998); Östlin et al 1999b ) reveal a merger morphology. Faint whisps and indications of shell structures characteristic of galactic disk progenitors are seen in the outer regions. Our studies of the kinematics of the ionized gas indicate that the starburst is triggered by coalescing clouds involved in the merger process. The H$`\alpha `$ velocity field is requires a mass of the main body of $`>`$10<sup>9</sup> $`_{}`$ (Östlin et al. 1999a ; 1999b ).
The oxygen abundance in the central region was calculated from the optical spectra, using the standard method. We assumed a homogeneous single metallicity model with two temperature zones. The electron temperature of the O<sup>++</sup> zone was derived from the \[OIII\]$`\lambda `$4363/(\[OIII\]$`\lambda \lambda `$(4959+5007) ratio and that of the O<sup>+</sup> zone from the semi-empirical relationship derived by Vila-Costas and Edmunds (vilacostas (1993)). The density was derived from the \[SII\]$`\lambda `$6717/6731 line ratio and was found to be low. With this approach we obtain a metallicity close to 10% solar. Such a value is unusually low for a galaxy of this luminosity and deviates strongly (Bergvall et al. bergvall5 (1998)) from the metallicity-luminosity relationship of normal dwarf galaxies (Skillman et al. skillman (1989)). Part of the optical emission may have a shock origin and we made an estimate how large the contribution from shocks may be by mixing predicted line spectra from shocked regions (Raymond raymond (1979)) with those of HII regions from Stasińska (stasinska1 (1982)). The best fits generally gave a contribution from the shock component to H$`\beta `$ of $`<`$ 5%. This is in most cases not enough to change the result of the calculations of the metallicities drastically.
A surprising fact is that no HI has been detected in the galaxy (Bergvall et al. 1999b ). With an upper limit of $``$$`{}_{HI}{}^{}=`$10<sup>8</sup> $`_{}`$ and consequently $``$<sub>HI</sub>/L$`{}_{B,total}{}^{}<`$0.01, it seems to be remarkably devoid of neutral hydrogen for its star formation rate. One possible explanation is that the neutral gas is ionized by the central starburst. Another possibility is that much of the gas is in molecular form. The ISO observations discussed here, which relate to the transition zone between the ionized and the cold molecular medium, will help to derive a coherent picture of the different phases of the gas.
### 3.2 The dust properties
From the 60 and 100 $`\mu `$ IRAS fluxes and assuming a FIR grain opacity of $`\kappa `$ = 2.5 10$`{}_{}{}^{3}\lambda _{}^{1}`$ cm<sup>2</sup>g<sup>-1</sup> (Hildebrand hildebrand (1983)), where $`\lambda `$ is the wavelength in microns, we calculate a dust temperature of T<sub>dust</sub> = 49K. We can then obtain a rough estimate of the mass of the warm dust from (Devereux & Young devereux (1990))
$$_{dust}=4.58f_{100}r^2(e^{144/T_{dust}}1)_{}$$
f<sub>100</sub> is in Jy and r is the distance in Mpc. We obtain $`_{dust}`$ 4 10<sup>6</sup> $`_{}`$.
It is interesting to compare gas-to-dust mass ratio with that of other gas rich galaxies. In Sect. 3.4.3 we find that the total hydrogen gas mass, i.e. HI, HII and H<sub>2</sub>, is approximately 10<sup>9</sup> $`_{}`$. We thus obtain $`_{gas}`$/$`_{dust}`$ 200. In a study of normal spiral galaxies, Devereux & Young (devereux (1990)), taking into account only the inner part of the HI disk containing the warm dust, found that $`_{gas}`$/$`_{dust}`$ 1100. There is thus a difference a factor of 5-6 difference between Haro 11 and normal spirals. If the dust content correlates with metallicity, one would expect that $`_{gas}`$/$`_{dust}`$ would increase with decreasing metallicity and indeed this is the case for dIrrs as was shown by Lisenfeld & Ferrara (lisenfeld (1998)). If we apply this metallicity correction in the comparison, it would lead to an even larger difference between Haro 11 and the spirals in the $`_{gas}`$/$`_{dust}`$ ratio. The high value for spiral galaxies has been shown (e.g. Chini et al. chini (1995)) to be due to the fact that a large amount of the dust is of the cold, ’cirrus’ type, radiating at longer wavelengths and thus not contributing to the mass estimates based on the IRAS data. The fact that $`_{gas}`$/$`_{dust}`$ of Haro 11 is so low (in fact close to the value found for our own Galaxy, including the total dust mass, i.e. $`_{gas}`$/$`_{dust,total}`$ 100-150) indicates that the galaxy contains very little dust at low temperature.
### 3.3 The LWS spectra
Figure 1 shows the spectral profiles of the three detected lines and the region around \[OI\]145$`\mu `$, for which we derived an upper limit. Table 2 summarizes the spectral data. From our optical spectra we have derived an extinction coefficient in the brightest region of the galaxy of A$`{}_{V}{}^{}`$ 0.8 mag. Thus the corrections for galactic extinction in the infrared are considerably less than the observational errors in the observed emission lines (Adams et al. adams (1988)).
### 3.4 Properties of the PDR gas
One of the purposes of this study is to investigate how much the PDR gas contributes to the total mass in this type of objects, and in particular Haro 11 due to its remarkably low HI content. It is known that the relative volume of the PDR that occupies the weakly ionized (only low ionization stages) gas may be quite substantial. Regions having a total extinction of A$`{}_{V}{}^{}`$10 are partly photodissociated (see e.g. Crawford et al. crawford (1985) or Tielens & Hollenbach tielens (1985) for a general discussion about the PDR physics). This is the situation for most of the neutral gas in the Galaxy. Estimates also show that $``$ 40% of the total gas mass in starburst nuclei are PDRs, i.e. the photodissociated gas is the dominant phase in these extreme environments (Stacey et al. stacey (1991)). The most widely used model for calculating the energy balance and obtain predictions about emissivities of the gas and dust was presented by Tielens & Hollenbach (tielens (1985)). Later upgrades and applications have been discussed by Hollenbach et al. (hollenbach (1991)), Wolfire et al. (wolfire (1990)) and Kaufman et al. (kaufman (1999)).
The Kaufman et al. standard model is based on an oxygen abundance of 40% solar, close to the local Galactic value and a factor of 4-5 higher than in Haro11. From our derived carbon abundance in ESO 338-IG04 (Tol 1924-416), a galaxy quite similar to Haro 11, we may assume that N(C)/N(O) $``$ 0.2 (Bergvall bergvall (1985)). Similar values have been obtained for other metal poor galaxies (Garnett et al. garnett (1995)). Thus, oxygen and carbon are underabundant with factors 4-8 with respect to the standard setup in Kaufmans et al. models (their models have N(C)/N(O) $``$ 0.5). The predicted temperature of the surface gas layers obtained from their basic diagnostic diagrams is therefore probably somewhat overestimated while still the line intensities are reasonably correct (Wolfire et al. wolfire (1990)). However, Kaufman et al. also discuss a case of lower metallicity, about 4% solar and present diagrams showing how \[CII\], FIR and CO luminosities relate to density, n, incident far UV radiation intensity, G<sub>0</sub>, and the two metallicities $``$ 4 and 40% solar. In the following we use these models to estimate some fundamental parameters of the PDR regions.
#### 3.4.1 Optical thickness
To enable a comparison of the line intensities with the models of Kaufman et al., we first estimate the contribution from HII regions to the line intensities and also the optical depth of the lines. The optical emission lines (Bergvall & Olofsson bergvall1 (1986) and newer data) have been used together with recent photoionization models (Stasińska stasinska (1990), and updates) to estimate the contribution to the far-IR lines from the HII regions. We used our observed line intensities in \[OI\], \[OII\], \[OIII\], \[NII\], \[SII\], HeI and HeII, corrected for extinction, underlying absorption and shock contribution to select the best fitting HII region models from Stasińskas library. As a result we find that the measured intensity of the \[OIII\]$`\lambda `$88$`\mu `$ line is fully consistent with photoionization by a cluster of stars with T<sub>eff</sub> of the order of 35000-40000 K. For the other 3 lines observed with the LWS, the contribution from HII regions to the measured intensity is strongest in \[CII\] but is still $``$ 20%. Thus, when discussing the \[CII\] line flux and flux ratios, we can directly compare to the predictions from Kaufman et al..
In a study of two star forming complexes in the Large Magellanic Cloud, Israel et al. (israel (1996)) concluded that these regions were optically thin at a minimum column density of 3 10<sup>21</sup> cm<sup>-2</sup>. Since these cloud aggregates probably have similar properties as Haro 11 we can conclude that the star forming regions in Haro 11 on a local scale probably are optically thin in 158$`\mu `$. The optical extinction derived from the H$`\alpha `$/H$`\beta `$-ratio in Haro 11 is only A<sub>V</sub> = 0.8 so we don’t expect that the \[CII\] line is optically thick even on a global scale. But the derived extinction depends on the spatial distribution of the HII regions so to assure that we don’t have conditions that allow hidden sources in the central area we will estimate the optical depth in \[CII\], $`\tau _{CII}`$. This calculation will also be used to make estimates of optical depths of the other galaxies involved in the discussion.
In the high-temperature (T$``$92K, the excitation potential for \[CII\]), high-density limit (n$``$4 10<sup>3</sup> cm<sup>-3</sup>, the limit for collisional deexcitation by H and H<sub>2</sub>), $`\tau _{[CII]}`$ can be expressed as (Crawford et al. crawford (1985))
$$\tau _{[CII]}=\mathrm{2.3\; 10}^{16}\frac{N_{C^+}}{T_{C+}\sigma _v}$$
where N<sub>C+</sub> is the column density of C<sup>+</sup> in cm<sup>-2</sup>, T<sub>C+</sub> is the temperature of the gas in K and $`\sigma _v`$ is the velocity dispersion in km s<sup>-1</sup>. If we as a first rough approximation assume that the global distribution of the gas is close to a homogeneous exponential disk we can derive the central column density in atoms cm<sup>-2</sup>
$$N_{gas}=\mathrm{1.3\; 10}^{14}\frac{_{gas}}{2\pi h_{gas}^2}$$
where $``$<sub>gas</sub> is the mass of the gas in solar masses and h<sub>gas</sub> is the scale length in kpc. As a first approximation we will assume h$`_{C^+}`$=h. From our H$`\alpha `$ images we derive h$`{}_{H\alpha }{}^{}`$ 2” = 1 kpc. As an estimate of the gas mass we will use our upper observational limit of $`_{HI}`$, i.e. around 10<sup>8</sup> $`_{}`$. With this we obtain a central column density in HI of N$`{}_{HI}{}^{}`$ 2 10<sup>21</sup> atoms cm<sup>-2</sup>. Later we will find it consistent to assume that most of the HI actually is associated with the PDR regions. Part of the C<sup>+</sup> region in the PDR also contains molecular gas. In sect. 3.4.2 we will show that this may increase the estimated column density with a factor of 2.
The C<sup>+</sup> column density can now be derived from
$$N_{C^+}\mathrm{1.3\; 10}^{14}X(C^+)\frac{_{HI}}{2\pi h_{H\alpha }^2}$$
where X(C<sup>+</sup>) is the relative abundance of carbon ions to hydrogen atoms. With these approximation we can write the optical depth as
$$\tau _{[CII]}0.03X(C^+)\frac{_{HI}}{T_{C+}\sigma _v2\pi h_{H\alpha }^2}$$
Since most of the carbon in the \[CII\] emitting region is in ionized form we will assume X(C<sup>+</sup>)$``$X(C<sub>total</sub>)$`\frac{N(C)}{N(O)}`$ X(O)$``$ 1.2 10<sup>-5</sup>. In starburst galaxies T<sub>C+</sub> is typically a few hundred K. We will later iteratively derive T<sub>C+</sub> and obtain a value below 400K. From our Fabry-Perot observations in H$`\alpha `$ of the central region we estimate $`\sigma _v`$ 60 km s<sup>-1</sup>. We then obtain $`\tau _{[CII]}`$ 0.001. In order to derive a more realistic estimate, we have to take into account that the gas has a lumpy, fractal distribution. But it is sufficient to note that, assuming a volume filling factor $`\alpha `$ of a few per mille (Crawford et al. crawford (1985)), the surface density would increase with $`\alpha ^{2/3}`$, i.e. a few tens. Thus it still would be safe to assume that the medium is thin in \[CII\]. The optically thick limit at a chemical abundance similar to Haro 11 would be at a mean column density of N$`{}_{HI}{}^{}`$ 10<sup>23</sup> atoms cm<sup>-2</sup> and with solar abundances N$`{}_{HI}{}^{}`$ 10<sup>22</sup> atoms cm<sup>-2</sup>, still with the assumption that the filling factor is a few per mille.
#### 3.4.2 The mass of the PDR gas
From the diagnostic diagrams by Kaufman et al. we can obtain the density and intensity of the incident far-UV flux G<sub>0</sub> in the PDR region. We use two diagrams for this purpose, one utilizing the \[OI\]63$`\mu `$/\[CII\]158$`\mu `$ intensity ratio and one utilizing the (\[OI\]63$`\mu `$+\[CII\]158$`\mu `$)/FIR ratio. As mentioned above, before we use the predictions we need to correct for the difference in metallicity between Haro 11 and the models, a factor of 3-4. We can estimate what effect this may have by comparing the predicted fluxes of these lines in an HII region at the different metallicities. The column densities and thus the volume emissivities of the low ionization lines are lower in HII regions than in a typical PDR region but the flux ratios should vary in approximately the same way. From Stasińskas models we find that the difference in flux ratios, assuming a constant N(C)/N(O), is about 2 %. N(C)/N(O) probably differ with a factor of about 2 however, and so we should make an appropriate correction for this which means a factor of 1.5-2 in the \[OI\]63$`\mu `$/\[CII\]158$`\mu `$ ratio. Kaufman et al. show in later diagrams that the metallicity has a minor effect on the \[CII\]158$`\mu `$/FIR ratio in the region of the parameter space where we are residing.
We now feel confident to compare the slightly corrected line intensities with the predictions from the PDR models. Using the \[OI\]63$`\mu `$, \[CII\]158$`\mu `$ line fluxes and the FIR flux, Kaufmans et al. diagnostic diagrams show that G<sub>0</sub> is 1-2 10<sup>3</sup> in units of the Galactic far-UV field (1.6 10<sup>-6</sup> W m<sup>-2</sup>). The obtained value for the density is n $``$ 2 10<sup>3</sup> cm<sup>-3</sup> and thus G<sub>0</sub>/n $``$ 1-2 cm<sup>3</sup>. These are typical values for normal starforming regions and galaxies (Tielens & Hollenbach, tielens (1985); Wolfire et al. wolfire (1990); Carral et al. carral (1994); Fischer et al. fischer (1996)). The surface temperature is T $``$ 400K. The data are summarized in Table 3. With these data at hand we can obtain the predicted \[OI\]145$`\mu `$/\[OI\]63$`\mu `$ ratio which is approximately 0.06. This is within the errors of the measured line ratio and thus supports the reliability of the model.
From the measured \[CII\] flux we may now calculate the mass of the warm atomic gas according to (e.g. Wolfire et al. wolfire (1990))
$$_{PDR}=\mathrm{5.96\; 10}^{15}\frac{r^2f_{[CII]}m_H}{\mathrm{\Lambda }(CII)X(C^+)}_{}$$
where r is the distance to the galaxy in Mpc, f<sub>\[CII\]</sub> is the \[CII\] flux in W m<sup>-2</sup>, m<sub>H</sub> is the mass of the hydrogen atom in kg, and $`\mathrm{\Lambda }`$(CII) is the cooling rate per atom in the 158$`\mu `$ line in W. In the high-temperature and high-density limit $`\mathrm{\Lambda }`$(CII) will be nearly independent of the temperature (Wolfire et al. wolfire (1990)). $`\mathrm{\Lambda }`$(CII) can therefore be rather safely estimated and we obtain $`\mathrm{\Lambda }`$(CII) $``$ 1.3 10<sup>-26</sup> W atom<sup>-1</sup>. In the surface layers of the PDR carbon is mainly in the low ionization monoatomic state so X(C<sup>+</sup>)$``$ X(C). We then derive $``$$`{}_{PDR}{}^{}`$ 2$`{}_{1}{}^{}{}_{}{}^{+2}`$ 10<sup>8</sup> $`_{}`$.
#### 3.4.3 The mass fractions of cold, warm and ionized gas
The value we have derived for the mass of the PDR:s is larger than or comparable to our upper limit of the HI mass, i.e. 10<sup>8</sup>$`_{}`$, indicating that most of the HI gas is located in the PDR regions if these are dominated by the atomic gas component. The total molecular gas mass in normal star forming regions is normally 5-10 times larger than in the PDR regions. Data from other starburst galaxies however indicate that considerably lower values are more likely (e.g. Wolfire et al. wolfire (1990); Stacey et al. stacey (1991)). This is probably a metallicity effect. In galaxies with low metallicities the UV radiation can penetrate deeper into the molecular clouds, increasing the PDR zone relative to the CO core. An upper limit is therefore probably around 10<sup>9</sup> $`_{}`$ in our case. This mass estimate may be compared to the upper limit of the H<sub>2</sub> mass, derived from the CO (1-0) observations which is $``$ 10<sup>8</sup> $`_{}`$. The calculation of this last value is however rather uncertain due to the uncertainty in the conversion factor between the CO flux and $``$$`_{H_2}`$ (Maloney & Black maloney (1988), Taylor et al taylor (1998)). As just mentioned, the low metallicity also allows a larger proportion of H<sub>2</sub> to be hidden in the CO dissociated zone. Therefore, in addition to the problem caused by the uncertainty in the conversion factor between CO and H<sub>2</sub>, the CO fluxes are poor indicators of the H<sub>2</sub> masses in low metallicity galaxies.
Using the available photometry of the halo and the burst, our spectral evolutionary models predict a total stellar mass of 1.5 10<sup>10</sup> $`_{}`$ , while the observed rotation results in a mass of $`\mathrm{2\; 10}^9`$ $`_{}`$ , assuming dynamical equilibrium (Östlin et al. 1999a ; 1999b ). The discrepancy between the photometric and dynamical mass estimates can be resolved if the galaxy is dominated by velocity dispersion, which is reasonable if the width of the H$`\alpha `$ line reflects potential motions, or if the galaxy is not in dynamical equilibrium.
From the global H$`\alpha `$ luminosity and the mean density of the gas, derived from our spectra, the mass of the ionized gas is estimated to be $``$ 10<sup>8</sup> \- 10<sup>9</sup> $`_{}`$ (Östlin et al. 1999a ,1999b ). It is a remarkable fact that the HI mass is so low and also lower than or comparable to the estimated mass of ionized or molecular gas. A possible explanation may be that HI to a large extent has become molecular and/or ionized in the merging process. If this is correct, then the gas/total mass estimates may have been severely underestimated in many dwarf starburst cases. This will complicate the efforts to understand dwarf galaxy evolution in general. A summary of the mass estimates is found in Table 4.
### 3.5 Thermal balance and the \[CII\]/FIR ratio
The important coolant \[CII\]$`\lambda `$158$`\mu `$ originates from the surface layers of PDRs where the C<sup>+</sup> ions are excited by photoelectrons ejected from grains heated by the UV radiation (Watson 1972a ). PAHs are believed to be the dominating agent of this process. Stacey et al. (stacey (1991)) studied a mixed sample of spiral galaxies, starburst galaxies, giant molecular clouds and galactic star forming regions and noticed that while there is a linear relation between the energy density of the UV field and the cooling rate of the dust, the efficiency with which the \[CII\] line is excited seems to decrease with the intensity of the UV field. One should have in mind however that this is not a homogeneous sample. Malhotra et al. (malhotra (1997)) investigated the same relationship between the \[CII\] cooling and dust cooling for different types of galaxies, using the \[CII\]/FIR luminosity ratio. An anticorrelation was found between the \[CII\]/FIR ratio and dust temperature, as measured by the IRAS f<sub>60</sub>/f<sub>100</sub> ratio. After a discussion of several possible explanations, Malhotra et al. found it most probable that this is a consequence of a decreasing efficiency of ejection of photoelectrons from the grains, reducing the \[CII\] intensity. Theoretically it can be understood if the ratio of UV flux to gas density, and thus the dust temperature, becomes high because the dust grains would be more positively charged, forming a stronger Coloumb barrier against photoejection.
One should be aware, both as regards the results from Malhotra et al. and in the following discussion, that the sample for which useful data are available is quite heterogeneous. The galaxy types involved range from normal galaxies to Seyferts for which the relative importance of different heating sources (stars and AGNs) and different chemical abundances are not satisfactorily controlled. Moreover, in the first approximation, two dust components, warm and cold, having different spatial distribution, can be separated. This increases the complexity in the interpretation of the global data and the 60/100 index as well as the L<sub>FIR</sub>/L<sub>B</sub> index cannot be regarded as clearcut tracers of the star formation activity. As a consequence we expect that the scatter in the discovered trends will remain quite large.
We will now compare the data for Haro 11 with a data set similar to that used by Malhotra et al. but including more detected galaxies at relatively high f<sub>60</sub>/f<sub>100</sub>. These data were obtained from Stacey et al. (stacey (1991)), Luhman et al. (luhman (1998)), Malhotra et al. (malhotra (1997)) and Lord at al. (lord (1996)).
For the most nearby galaxies, having velocities $``$ 500 kms<sup>-1</sup>, we have scanned the literature to find the most recent distance determinations derived from photometric distance indicators. This concerns the galaxies M82, NGC 3109 (Kennicutt et al. kennicutt2 (1998)), NGC 4414 (Turner et al. turner (1998)), NGC 4736 (Garman & Young garman (1986)), NGC 5128 (Tonry et al. tonry (1995)), NGC 5194 (Sandage sandage (1987)), NGC 6946 (Freedman et al. freedman (1994)) and IC 342 (Krismer et al. krismer (1995)). In order to take proper account of the mid-IR contribution to the FIR luminosities, the FIR luminosities were calculated from FIR = 5.35 10<sup>5</sup> r<sup>2</sup>(12.66 f<sub>12</sub>+5.00 f<sub>25</sub>+2.55 f<sub>60</sub>+1.01 f<sub>100</sub>) L, where r is the distance to the galaxies in Mpc and f<sub>12</sub>, f<sub>25</sub>, f<sub>60</sub> and f<sub>100</sub> are the apparent flux densities in Jy (Belfort et al. belfort (1987)), instead of the normal approximation, e.g. FIR = 3.75 10<sup>5</sup> r<sup>2</sup>(2.58 f<sub>60</sub>+1.0 f<sub>100</sub>) L (Lonsdale & Helou lonsdale (1985)). This choice makes a small but noticeable change in the distribution of the data in the diagrams.
Fig 2 displays the distribution of Haro 11 and the comparison sample in the IRAS f<sub>25</sub>/f<sub>60</sub> vs. f<sub>25</sub>/f<sub>100</sub> diagram, showing the extreme position of Haro 11. Fig. 3 shows the FIR/B-60/100 diagram for the sample galaxies. In the diagram are also indicated the loci of normal spiral galaxies and blue compact galaxies (Dultzin et al. deborah1 (1988)). As we see, Haro 11 also in this case lies at the extreme of the BCG distribution. The extremely strong 25$`\mu `$ emission confirms the lack of cold dust and indicates the presence of a strong nuclear starburst (Dultzin-Hacyan et al. deborah1 (1988); Hawarden et al hawarden (1986); see also Boulanger et al. boulanger (1994); Taniguchi et al. taniguchi (1990)), destroying the smallest dust grains (e.g. Desert desert (1990)).
Fig. 4 and 5 correspond to Malhotras et al. diagram 1. To extract the most important global parameters involved in the relation reflected in fig. 4, we made a cross correlation test between the \[CII\]/FIR flux ratio, the IRAS data (i.e. the ratios between the 12, 25, 60 and 100$`\mu `$ fluxes), the FIR/B flux ratios and the FIR flux. We found strong correlations ($`6\sigma `$) between only three of these parameters: f<sub>60</sub>/f<sub>100</sub>, FIR/B and the \[CII\]/FIR ratio.
In figure 4 the trend found by Malhotra et al. is confirmed and appears even stronger. But it is also evident that the position of Haro 11 deviates from the envelope defined by the other galaxies. The possible explanations for the deviant position of Haro 11 are either that there is an excess in the \[CII\] over FIR luminosity or that the f<sub>60</sub>/f<sub>100</sub> temperature is high, or both. Let us look at these two options one at a time. A consequence of lower metallicity is that the UV photons penetrate deeper into the cloud, thus increasing the volume of the PDR region at the expense of the colder molecular core. Strong observational support of this in terms of a high \[CII\]/CO flux ratio have been reported by e.g. Mochizuki et el. (mochizuki (1994)), Poglitsch et al. (poglitsch (1995)), Madden et al. (madden1 (1997)), Smith & Madden (smith (1997)) and Israel et al. (israel (1996)). Haro 11 behaves in the same way with its extremely high \[CII\]/CO ratio ($`>`$ 4 10<sup>5</sup>). The interesting question is how much the \[CII\]/FIR ratio will be influenced by varying metallicity under these conditions. The models by Kaufman et al. show that \[CII\]/FIR is almost unaffected by metallicity variations in cases where G<sub>0</sub>/n have normal values. This is also in agreement with observations. In 30 Dor and the cases discussed by Israel et al. however, the G<sub>0</sub>/n ratio is low ($``$ 0.1) which results in a diluted radiation field, an increased heating efficiency and brighter \[CII\] emission. In Haro 11 the G<sub>0</sub>/n ratio is normal (G<sub>0</sub>/n=1 cm<sup>3</sup>), as seen from table 3. So, both observations and modelling argue against an increase in \[CII\]/FIR in Haro 11 due to low metallicity. There is also another reason. In fig. 5 Haro 11 adheres to the other galaxies in the diagram. Therefore as a first guess it is reasonable to assume that \[CII\]/FIR is normal while f<sub>60</sub>/f<sub>100</sub> is deviating.
We concluded in section 3.2 that Haro 11 contains very little warm or cold dust. Haro 11 is a merging system of galaxies and it is possible that the merging galaxies have lost their cold dust components, e.g. extended disks, in the merger process. Is it possible that the lack of a ’normal’ cold dust component could explain the high f<sub>60</sub>/f<sub>100</sub> ratio? Would the contribution from such a component be sufficient to move the position of Haro 11 into the region of the other galaxies? It is an important question to ask since it separates global phenomena from local ones. We have made a rough check of this possibility by selecting the galaxies in the comparison sample having apparent sizes sufficiently large to allow us to obtain IRAS data both from the very centre and also the integrated IRAS fluxes of the entire galaxy. In this way we can have an idea about how much the extended component influences the 60/100 ratio. It turns out that in the majority of the cases the IRAS temperature actually increases when the extended component is included. Moreover, for the hottest cases, the hot component is always significantly stronger than the warm/cold component, if existing. Our impression is therefore that the IRAS temperature of the warm dust of the most active regions in Haro 11 and the galaxies below Haro 11 in fig. 4 are similar.
Although we expect a large scatter in fig. 4 we should examine if there may be other parameters involved ruling the balance between \[CII\] and FIR emission and explaining part of the scatter in the diagram and the deviant position of Haro 11. Since both the spectrum of the radiation field and the spatial distribution of gas and dust in general differs between low-mass starbursts and galaxies with AGNs, we could suspect that we would also find systematic differences in the \[CII\]/FIR ratios between these types of galaxies. When comparing galaxies classified in the optical region as Sy1, Sy2, liners and HII galaxies we find that the HII galaxies appear to have a broad distribution in f<sub>60</sub>/f<sub>100</sub> while the Sy2 and liners tend to concentrate towards high IRAS temperatures and high FIR/B. Similarly, in a statistical investigation of the IR/optical properties of BCGs as related to Seyfert and Liner galaxies (Dultzin-Hacyan et al. deborah1 (1988), deborah2 (1990)) it was shown that the f<sub>25</sub>/f<sub>100</sub> is a powerful index that can be used to discriminate between pure starbursts and cases where AGNs are important for the radiation field.
### 3.6 What rules the \[CII\]/FIR ratio?
Several possible physical mechanisms have been proposed to explain the trends in \[CII\]/FIR seen in fig. 4 and 5 (e.g. Malhotra et al. malhotra (1997); Luhman et al. luhman (1998)). As we already mentioned, Malhotra et al. argued that the reduced efficiency in photoejection due to a hard radiation field could be a plausible explanation. The data for Haro 11 however, acts as a counterexample of this. Both optical spectroscopy and the ’hot’ IRAS indices indicate a hard radiation field. Still, the \[CII\]/FIR is close to the values of normal, more passive galaxies. An alternative explanation, that have been discussed by Malhotra et al. and others is that \[CII\] actually is optically thick in many of the cases where the \[CII\]/FIR ratio is low. The general opinion seems to be that the \[CII\] line is optically thin. These conclusions are based on estimates of the optical depth in a few well known cases, e.g. M82 (Crawford crawford (1985)). But M82 is behaving well in the diagrams (see fig. 5) and the central column density of the gas is not abnormal. From the discussion above we see that, depending on the velocity field and the relative amount of molecular gas in the C$`+`$ zone, the medium could become optically thick in \[CII\] at a column density of N $``$ 10<sup>22-23</sup> cm<sup>-2</sup>. Such high numbers can be found in the central regions of some luminous starburst galaxies. Gerin & Phillips (gerin (1998)) elaborated on the low \[CII\]/FIR ratio in Arp 220 and it is interesting to compare this galaxy with M82 and Haro 11 (see fig. 5). They came to the conclusion that with an estimated column density of N $``$ 10<sup>24</sup> cm<sup>-2</sup>, the dust opacity and emission would be sufficiently high to explain at least part of the low value of \[CII\]/FIR. But in fact there may be three sources responsible for reducing the ratio: dust opacity, dust emission and self absorption. What is maybe most important is how to explain why the \[CII\]/L<sub>Bol</sub> ratio decreases with FIR/B. Therefore the combined effect of dust absorption and \[CII\] selfabsorption should be investigated.
Fig. 6 shows the same data as in fig. 5, divided into different galaxy types, Sy1, Sy2, HII and Liners. As we can see, there is a predominance of Sy2 galaxies and Liners at high FIR/B. Most of these galaxies are ULIGs. There are strong indications that Liners are galaxies with massive nuclear starbursts and likewise some of the Sy2 galaxies. Many of these galaxies are probably remnants of massive merges. Even though the luminosity of some of them are not extreme, high column densities could possibly be obtained also in less massive nuclear starbursts, provided they are sufficiently compact. Such very compact nuclear starbursts are preferentially found in the type of galaxy we find at high FIR/B in the diagram (Heisler & Vader heisler (1986)). In a simplified form fig. 6 shows the effect of increasing optical depth on the sample. Where the gas is optically thick in \[CII\] the ratio between the \[CII\]/B flux ratio should be approximately constant. Where it is thin, \[CII\]/FIR should be constant.
## 4 \[CII\] as a probe of high redshift starbursts
As we can see from figure 4 and 5, the \[CII\] line in starburst galaxies has a maximum luminosity of a few per mille of the bolometric luminosity and constitutes the peak in the spectral distribution from radio to X-ray. It has therefore been proposed (e.g. Petrosian et al. petrosian (1969); Loeb loeb (1993); Stark stark (1997)) that it can be used as a probe in searches for high redshift starburst galaxies. At redshifts of z$``$ 5-6, where current models of galaxy formation predict the first massive galaxies to form and where optical-near IR observations start to reach the confusion limit, the \[CII\] line is shifted into the submillimeter (submm) region and could be detected with a suitable submm telescope at a cold dry site like the South Pole (Stark stark (1997)). The most interesting targets would probably be massive starburst mergers, characterized by a high 60/100 ratio and a high FIR/B ratio. The tendencies seen in the study of Malhotra et al. (malhotra (1997)) however, have created some pessimism (e.g. Gerin & Phillips gerin (1998)) about the feasibility, due to the rapid decrease in the relative strength of the \[CII\] line with IRAS temperature. The properties of Haro 11 however, indicate that the conditions may be more promising than anticipated since the \[CII\]/L<sub>Bol</sub> of this galaxy is close to the value obtained for the normal galaxies. If this is due to the low metallicity of the galaxy or not remains to be explored but if so it may apply on distant metal-poor galaxies. The detection probabilities derived by Stark (stark (1997)) predict that it would be quite possible to detect a galaxy with a mass equal to that of a present day L galaxy at z=4-5 in a few hours, using a 10m submm telescope at the South Pole. At the formation epoch a massive galaxy would be several magnitudes brighter than today, i.e. $``$ 10<sup>12</sup> L. Thus we would need to merge about 10 Haro 11 type galaxies to produce this luminosity. A beam size of 20”-30” for a 10m telescope corresponds to about 100 kpc at the redshift we discuss. Thus it would easily encompass a compact merging group of dwarf galaxies and make a detection feasible. Moreover, with the advent of planned large sub-mm/mm arrays (e.g. ALMA), the prospects for using the \[CII\]158 line as a probe of high redshift starburst would be even more interesting.
## 5 Conclusions
We have reported about ISO LWS observations of the luminous metal poor BCG Haro 11 in the three prominent emission lines \[OI\]63$`\mu `$, \[OIII\]88$`\mu `$ and \[CII\]1588$`\mu `$. The galaxy is involved in an intense global starburst probably due to a merger of dwarf galaxies. The ISO data and previous observations in the optical-radio region in combination with model predictions are used to derive information about the physical properties of gas and dust in the galaxy.
Haro 11 is one of the hottest IRAS galaxies observed and we find no trace of cold dust. Most of the neutral hydrogen is confined to the photodissociation regions. The main part of the gas however, appears to be in ionized or molecular state. This possibility is generally neglected when discussing the gas content in starburst dwarfs, leading to serious underestimates of the total gas mass.
We also reinvestigate the claimed correlation between \[CII\] emission, far-IR luminosity and the IRAS 60/100 temperature for a number of starforming galaxies of different types, active and non-active. We confirm and tighten the correlation and note that Haro 11 deviates in the sense that its \[CII\]/L<sub>FIR</sub> is higher than expected. We argue that this may indicate that the previously preferred explanation of the relationship, a decreasing efficiency with IRAS temperature in the ejection of photoelectrons from UV-illuminated grains may not be entirely correct. An alternative, or complementary explanation seems to be that the optical depth increases with increasing IRAS temperature. This would be expected if the IRAS temperature correlates with compactness and/or mass of the central starburst region. The metallicity may also play an important role, in which case the somewhat pessimistic opinions about using \[CII\]158$`\mu `$ to detect high redshift massive starburst galaxies may not be valid. The whole situation will probably be clarified once a more representative sample of galaxies, including both normal dwarfs and massive galaxies becomes available.
###### Acknowledgements.
We gratefully acknowledge partial support from the Swedish Natural Science Research Council and the Swedish Space Board. JM and JC are supported by Spanish CICYT grant ESP98-1351. JM has been also supported by DIGCYT grant PB93-0139. G. Östlin acknowledges support from The Swedish Foundation for International Cooperation in Research and Higher Education (STINT). We thank our referee for useful suggestions for improvements of the manuscript. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration.
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# Center Vortices at Strong Couplings and All Couplings 11footnote 1Talk presented at Confinement 2000, Osaka, Japan, March 7-10 2000.
## 1 Why Center Vortices?
We begin with a simple question: If confinement is defined by the Wilson area-law criterion, then what charge is actually confined in an SU(N) gauge theory? The first answer that comes to mind is that confined charge is just the SU(N) color charge. But this answer can’t be quite right, at least according to the Wilson criterion, because not all color charges are confined in this sense. For example, due to color screening, there exists no asymptotic linear potential between heavy charges in the adjoint representation. A second possible answer, motivated by dual-superconductor models, is that abelian electric charge (identified by abelian projection $`SU(N)U(1)^{N1}`$) is the charge that is confined. But this doesn’t work either, since not all electric charges are confined. In, e.g., SU(2) lattice gauge theory, there is no asymptotic linear potential between charges of $`q=\pm 2`$ multiples of the elementary charge. Finally, consider N-ality, i.e. the charge associated with the $`Z_N`$ subgroup of SU(N). It is well known that in SU(N) gauge theory, only color charges with non-zero N-ality are confined, and the asymptotic string tension depends only on the N-ality of the color charge representation. Thus we conclude that it is really $`Z_N`$ charge which is confined in non-abelian gauge theories.
This simple fact is interesting, because the type of charge confined is an indication of the type of field configuration which does the confining. Of course, the N-ality dependence of the QCD string tension is no mystery; it is simply due to color screening of higher charge representations by gluons. We have an intuitive picture of one or more gluons bound to a static charge. On the other hand, Wilson loops can also be interpreted as a probe of vacuum fluctuations in the absence of external sources (think of evaluating spacelike loops in the Hamiltonian formulation), and it is generally assumed that loop observables are disordered by certain large-scale topological configurations. If that is the case, then such configurations must have the very non-trivial property that N-ality $`=0`$ loops are somehow *not* disordered, and that the induced string tension in general depends only on N-ality. The only known gluonic field configurations with this property are the center vortices.
The center group is also singled out by the deconfinement phase transition, which involves the breaking of a global $`Z_N`$ symmetry, and certain features found in the deconfined phase are elegantly explained in terms of the vortex picture. Note that only global symmetries can actually break spontaneously. In the absence of gauge fixing, the VEV of a Higgs field in any gauge theory, dual or otherwise, is zero in any phase. This is due to the Elitzur theorem, and also to the status of local gauge symmetry as a genuine redundancy in field variables.
Let us consider the non-confinement of abelian electric charge in more detail (cf. Ambjørn et al. for an extended discussion). Fixing to maximal abelian gauge, we write the link variables in the usual form $`U_\mu =W_\mu A_\mu `$ where $`A_\mu `$ is the abelian (diagonal) link variable. Imagine integrating out the $`W`$ and ghost fields, to obtain an effective abelian action
$$\mathrm{exp}\left[S_{eff}[A]\right]=DW_\mu D(\text{ghosts})e^{S+S_{gf}}$$
(1)
The reduction to $`U(1)^{N1}`$ degrees of freedom is not particularly significant in itself; this procedure could be carried out for any subgroup of SU(N), including the $`Z_N`$ center. Moreover, such effective actions are likely to be extremely complicated and non-local. The reduction only becomes physically interesting if $`S_{eff}[A]`$ takes on some particularly simple form at large scales, e.g. if $`S_{eff}[A]`$ describes a dual superconductor and/or monopole Coulomb gas of some kind. There exist, in fact, some very concrete proposals along these lines. One feature which is always found in such proposals is that *all* multiples of abelian electric charge are confined by the dual Meissner effect. If that is really so, then an immediate consequence is that abelian Polyakov lines corresponding to any integer multiple of the electric charge ought to vanish in the confined phase. In SU(2) lattice gauge theory, this means that
$$P_q=\text{Tr}[(A_0A_0\mathrm{}A_0)^q]_{S_{eff}}=0\text{all }q$$
(2)
In addition, if $`P_{Mq}`$ denotes the “monopole dominance” approximation to $`P_q`$, and if $`S_{eff}[A]`$ is well described by a monopole Coulomb gas, we would also expect
$$P_{Mq}P_q$$
(3)
The fact is, however, that while both of the above relations hold at $`q=1`$ ($`Z_2`$ charged), neither relation holds at $`q=2`$ ($`Z_2`$ neutral). The relevant $`q=2`$ Polyakov lines in the confined phase (for $`T=4`$ lattice spacings in the time direction), are shown in fig. 1 below. From the data showing $`P_20`$, we conclude that there is no confinement for $`q=2`$. From $`|P_2||P_{M2}|`$, it seems there is no monopole dominance either. And from the fact that $`P_2<0`$ we find that positivity is broken as well. Actually, positivity can be restored by going to spacelike maximal abelian gauge, but then $`90^{}`$ rotation symmetry is lost, and $`q=2`$ string-breaking still occurs.
Non-confinement of $`q=2`$ charge, and the breakdown of monopole dominance, pose severe difficulties for dual-superconductor/dual abelian Higgs/ monopole Coulomb gas models based on the abelian projection. All of these pictures predict confinement of all $`q`$; therefore none of them is a good description of $`S_{eff}[A]`$.
A similar objection can be raised to the monopole confinement picture in the D=3 Georgi-Glashow model. Although the monopole Coulomb gas picture, developed by Polyakov, is certainly valid for some intermediate range of distances, this picture must break down asymptotically. The reason is that in a monopole Coulomb gas we find string tensions $`\sigma _qq`$ between objects with $`q`$ units of U(1) charge. But the Georgi-Glashow model has W-bosons with $`q=2`$. Taking charge screening by these fields into account, we must get eventually
$$\sigma _q=\{\begin{array}{cc}\sigma _1& \text{odd}q\hfill \\ 0& \text{even}q\hfill \end{array}$$
(4)
which contradicts the Coulomb gas picture. There is a moral here: In a confining theory, *massive charged fields are relevant to far-infrared vacuum structure,* and cannot be ignored.
These comments apply also to the Seiberg-Witten model. The low-energy effective action, derived in this model, again explicitly neglects the massive W-particles, and therefore misses the screening effects due to those particles. For this reason, the Seiberg-Witten effective action (which assumes locality) is not really the same thing as a Wilsonian effective action obtained from integrating out massive charged fields, and does not adequately describe physics beyond the double-charge screening scale.
In pure SU(2) gauge theory in a physical abelian gauge, non-confinement of $`q=`$ even charge can likewise be deduced from the inevitable electric charge screening by off-diagonal gluons (the W-fields of eq. (1)). The effect is no mystery, but the consequences are important. While non-confinement ($`\sigma _q=0`$) for $`q=`$ even must be a property of the true effective action $`S_{eff}[A]`$ for the abelian field, a very different $`q`$-dependence ($`\sigma _qq`$) is found in dual superconductor and monopole gas models. This indicates that the latter are either incorrect or, at best, incomplete in some way. In contrast, the correct $`q`$-dependence of the abelian string tension is quite natural in the framework of the vortex theory (cf. Ambjørn et al. ).
## 2 Center Vortices at Strong Couplings
In strong-coupling lattice gauge theory in $`D>2`$ dimensions, we have both confinement for N-ality $`0`$ charges, and color screening for N-ality $`=0`$ charges. These facts suggest the existence of a vortex mechanism. On the other hand, there is a bit of folklore about strong coupling, namely, that confinement in $`D>2`$ dimensions is just due to plaquette disorder, as in $`D=2`$ dimensions. If so, vortices (and any other topological objects), have nothing to do with confinement at strong coupling. This folklore, however, is misleading.
Consider SU(2) lattice gauge theory at strong-coupling, and denote by $`U(C)`$ the product of link variables around loop $`C`$. Let the minimal area of a planar loop be decomposed into a set of smaller areas, bounded by loops $`\{C_i\}`$. We ask: Do the $`\{U(C_i)\}`$ fluctuate (nearly) independently, for large areas and small $`\beta `$? The test is whether
$$<\underset{i}{}F[U(C_i)]>\stackrel{\mathrm{?}}{=}\underset{i}{}<F[U(C_i)]>$$
(5)
for any class function
$$F[g]=\underset{j0}{}f_j\chi _j[g]$$
(6)
In fact, in $`D=2`$ dimensions, it is easy to show that this equality is satisfied exactly. However, for dimensions $`D>2`$, evaluating the left- and right-hand sides of (5) we find for the exponential falloff on each side
$$e^{4\sigma P(C)}\underset{i}{}\frac{1}{3}f_1\underset{i}{}f_1e^{4\sigma P(C_i)}$$
(7)
where the inequality holds for perimeters $`P(C)_iP(C_i)`$. The conclusion is that the holonomies $`U(C_i)`$ do *not* fluctuate independently, even at strong-coupling, for $`D>2`$. Where, then, does the area-law falloff come from?
The question is resolved by extracting a center element from the holonomies
$$z[U(C)]=\text{signTr}[U(C)]Z_2$$
(8)
and asking if the center elements fluctuate independently; i.e
$$<\underset{i}{}z[U(C_i)]>\stackrel{\mathrm{?}}{=}\underset{i}{}<z[U(C_i)]>$$
(9)
In fact, it is easy to show that they do:
$$e^{\sigma A(C)}\underset{i}{}\frac{3}{4\pi }=\underset{i}{}\frac{3}{4\pi }e^{\sigma A(C_i)}$$
(10)
Thus, confining disorder is center disorder, at least at strong couplings. Confining configurations must disorder the center elements $`z`$, but not the coset elements, of SU(2) holonomies $`U(C_i)`$. Again, the only configurations known to have this property are center vortices.
If center vortices are, in fact, the confining configurations of strong-coupling lattice gauge theory, then it would be interesting if this fact could be demonstrated analytically. A reasonable conjecture is that if the Wilsonian effective action could be computed at a scale beyond the vortex thickness (4-5 lattice spacings at strong couplings), then at this scale “thin” center vortices will be stable saddlepoints of the action.
Suppose we define an effective long-range action $`S_{eff}`$ by, e.g.
$$\mathrm{exp}\left[S_{eff}[V]\right]=DU\underset{l^{}}{}\delta [V_l^{}^{}(UU..U)_l^{}I]e^{S_W[U]}$$
(11)
where the V-lattice spacing is $`L`$ U-lattice spacings. In D=2 dimensions
$`\mathrm{exp}\left[S_{eff}[V]\right]`$ (12)
$`=`$ $`𝒩\mathrm{exp}\left[{\displaystyle \underset{P^{}}{}}\mathrm{log}\left(1+{\displaystyle \underset{j=\frac{1}{2},1,\frac{3}{2}}{}}(2j+1)\left({\displaystyle \frac{I_{2j+1}(\beta )}{I_1(\beta )}}\right)^{L^2}\chi _j[V(P^{})]\right)\right]`$
But this must be wrong for $`D>2`$ dimensions, because it leads to a perimeter-law falloff
$$\chi _1[V(C)]\mathrm{exp}[\mu 𝒫(C)]$$
(13)
with an $`L`$dependent “gluelump” mass
$$\mu =4L\mathrm{log}\left(\frac{\beta }{4}\right)\text{(wrong)}$$
(14)
The correct coefficient is
$$\mu =4\mathrm{log}\left(\frac{\beta }{4}\right)$$
(15)
In fact, $`S_{eff}`$ in $`D>2`$ dimensions is non-local. It contains loops in all $`j=`$ integer representations, with perimeter-law weightings, derived from diagrams in which plaquettes on the U-lattice form a “tube” around a contour $`C`$ on the V-lattice. These diagrams lead to non-local contributions to $`S_{eff}`$ such as
$$S_{eff}[V]\left(\frac{\beta }{4}\right)^{4(𝒫(C)4)}\chi _1[V(C)]$$
(16)
We would like to derive a local effective action which would produce at least the *leading* contribution to any Wilson loop on the V-lattice. To achieve this, we integrate over all links on the U-lattice except on 2-cubes surrounding V-lattice sites, as shown in fig. 2.<sup>2</sup><sup>2</sup>2We work in D=3 dimensions. The extension to D=4 should be straightforward. This defines an effective action $`\stackrel{~}{S}_L`$
$`Z`$ $`=`$ $`{\displaystyle DV\underset{l2cubes}{}d\stackrel{~}{U}_l}`$
$`\left\{{\displaystyle }{\displaystyle \underset{l^{\prime \prime }2cubes}{}}dU_{l^{\prime \prime }}{\displaystyle \underset{l^{}}{}}\delta [V_l^{}^{}(UU..U)_l^{}I]e^{S_W[U]}\right\}`$
$`=`$ $`{\displaystyle DV\underset{l2cubes}{}d\stackrel{~}{U}_l\mathrm{exp}\left[\stackrel{~}{S}_L[V,\stackrel{~}{U}]\right]}`$
Introduce group-valued plaquette variables in the 2-cubes $`h,g`$, where $`h`$ variables run around plaquettes on the surface of the 2-cube, and $`g`$ variables run around plaquettes in the interior. Both sets of contours begin and end at the center of the 2-cube. After a number of manipulations, which include changing variables from $`\stackrel{~}{U}`$ to $`h,g`$ and integrating over the $`g`$ variables, we obtain
$`Z`$ $``$ $`{\displaystyle }DVDh{\displaystyle \underset{2cubesK}{}}\{1+2\left({\displaystyle \frac{\beta }{4}}\right)^3{\displaystyle \underset{cK}{}}\chi _{\frac{1}{2}}[(hhh)_c]`$ (17)
$`+2\left({\displaystyle \frac{\beta }{4}}\right)^4{\displaystyle \underset{\stackrel{adjacent}{c_1c_2K}}{}}\chi _{\frac{1}{2}}[(hhh)_{c_1}(hhh)_{c_2}]+\mathrm{}\}`$
$`\times `$ $`\mathrm{exp}[{\displaystyle \frac{\beta }{2}}{\displaystyle }\text{Tr}[h]+2\left({\displaystyle \frac{\beta }{4}}\right)^{4(L2)}{\displaystyle \underset{l^{}}{}}f_l^{}^{ijkl}\text{Tr}[h_{ij}^{}V_l^{}h_{kl}^{}V_l^{}^{}]`$
$`+2\left({\displaystyle \frac{\beta }{4}}\right)^{L^2}{\displaystyle \underset{P^{}}{}}\text{Tr}[VVV^{}V^{}]]`$
where $`f_l^{}^{ijkl}=1`$ if two plaquettes on neighboring 2-cubes can be joined by a cylinder of plaquettes on the U-lattice adjacent to V-link $`l^{}`$; $`f_l^{}^{ijkl}=0`$ otherwise. This resembles an adjoint-Higgs theory, with an SU(2) gauge field $`V_\mu `$ coupled to 24 “matter” fields $`h`$ in the adjoint representation. Note that for large $`L`$, the “Higgs” potential term is much larger than the “kinetic” (V-link) and pure-gauge (V-plaquette) terms, so the $`h`$-fields fluctuate almost independent of $`V_\mu `$. We then do a unitary gauge-fix of the $`h`$-fields (which leaves a remnant $`Z_2`$ symmetry), and integrate out the remaining $`h`$ d.o.f. to obtain
$`S_{eff}[V]`$ $``$ $`S_{link}[V,h_h]+S_{plaq}[V]`$ (18)
$`=`$ $`2\left({\displaystyle \frac{\beta }{4}}\right)^{4(L2)}{\displaystyle \underset{l^{}}{}}f_l^{}^{ijkl}\text{Tr}\left[h_{ij}^{}_hV_l^{}h_{kl}^{}_hV_l^{}^{}\right]`$
$`+2\left({\displaystyle \frac{\beta }{4}}\right)^{L^2}{\displaystyle \underset{P^{}}{}}\text{Tr}[VVV^{}V^{}]`$
Now look for saddlepoints. We find that $`S_{link}`$ is maximized at
$`V_\mu (\stackrel{}{n})`$ $`=`$ $`Z_\mu (\stackrel{}{n})\times g(\stackrel{}{n})g^{}(\stackrel{}{n}+\widehat{\mu })`$
$`Z_\mu `$ $`=`$ $`\pm 1`$
where $`g(\stackrel{}{n})g^{}(\stackrel{}{n}+\mu )`$ is fixed by the particular unitary gauge choice, while $`S_{plaq}`$ is maximized if $`ZZZZ=+1`$. This is the unitary gauge ground state. Create a thin center vortex on this state by a discontinuous gauge transformation, e.g.
$`Z_y(\stackrel{}{n})`$ $`=`$ $`\{\begin{array}{cc}\hfill 1& n_12,n_2=1\hfill \\ \hfill +1& \text{otherwise}\hfill \end{array}`$ (21)
$`Z_x(\stackrel{}{n})`$ $`=`$ $`Z_z(\stackrel{}{n})=1`$
This configuration is stationary: $`S_{link}[V]`$ is still a maximum, and $`S_{plaq}`$ is extremal (max or min) on all plaquettes. Stability depends on the eigenvalues of
$$\frac{\delta ^2S_{eff}}{\delta V_\mu (n_1)\delta V_\nu (n_2)}=\frac{\delta ^2S_{link}}{\delta V_\mu (n_1)\delta V_\nu (n_2)}+\frac{\delta ^2S_{plaq}}{\delta V_\mu (n_1)\delta V_\nu (n_2)}$$
(22)
and we find
$`{\displaystyle \frac{\delta ^2S_{link}}{\delta V_\mu (n_1)\delta V_\nu (n_2)}}`$ $``$ $`\left({\displaystyle \frac{\beta }{4}}\right)^{4(L2)+12}`$
$`{\displaystyle \frac{\delta ^2S_{plaq}}{\delta V_\mu (n_1)\delta V_\nu (n_2)}}`$ $``$ $`\left({\displaystyle \frac{\beta }{4}}\right)^{L^2}`$ (23)
The crucial observation is that for $`\beta /41`$ and
$$4(L2)+12<L^2L5$$
(24)
the contribution of $`\delta ^2S_{plaq}/\delta V\delta V`$ to the stability matrix (and therefore to the eigenvalues of the stability matrix) is negligible compared to $`\delta ^2S_{link}/\delta V\delta V`$, which has only stable modes. This implies:
1. *Vortex Stability:* The thin vortex is a stable saddlepoint of the full effective action $`S_{eff}`$ at $`L5`$.
2. *Vortex Thickness:* A “thin” vortex on the V-lattice means thickness $`<L`$ on the U-lattice. This means that stable center vortices are $`45`$ lattice spacings thick. For the strong coupling Wilson action, this is the distance where the adjoint string breaks! The correspondence between the adjoint string-breaking length, and the thickness of center vortices, has been emphasized by our group in connection with Casimir scaling (see also Cornwall ).
3. *Percolation:* From $`S_{eff}`$, we see that center vortices in D=3 cost an action $`8(\beta /4)^{L^2}`$/unit length, while the entropy is O(1)/unit length. Since entropy $``$ action, this implies that vortices percolate through the lattice, and confine N-ality $`0`$ charge.
## 3 P-Vortices, Gauge Copies, and Lattice Size
The original calculations of center-projected Creutz ratios $`\chi _{cp}(I,I)`$, in direct maximal center gauge, used 3 gauge copies for gauge-fixing each lattice (picking the best of the three). Very recently Bornyakov et al. have claimed that $`\chi _{cp}(I,I)`$ varies with the number of gauge copies used, and disagrees, in the large copy number limit, with the unprojected string tension by as much as 30%.
In our opinion, the reported strong disagreement between projected and unprojected string tensions is due to finite-size effects. Lattice sizes used by Bornyakov et al. were $`12^4`$ at $`\beta =2.3,2.4`$, and $`16^4`$ at $`\beta =2.5`$ while our published results were obtained on $`16^4`$ lattices at $`\beta =2.3,2.4`$ and $`22^4`$ lattices at $`\beta =2.5`$ Projected lattices are more sensitive to finite-size effects than unprojected lattices, and this is probably due to the fact that center projection has difficulty finding vortices, when most of the lattice volume is taken up by the vortex cores. There are now good estimates for vortex thickness, coming from three sources: First, from the ratio of “vortex-limited” Wilson loops. Second, from the adjoint string-breaking distance, measured by de Forcrand and Philipsen. Third, the vortex thickness is found in a very interesting calculation, reported here by Terry Tomboulis, of the vortex free energy vs. lattice size. All three estimates are in rough agreement, and give a vortex thickness of a little over one fermi. This means that center vortices are $`1214`$ lattice spacings thick at $`\beta =2.5`$; a $`16^4`$ lattice may just be too small on this scale.
We have therefore repeated the calculation of Bornyakov et al. on a variety of lattice sizes. Some typical results at $`\beta =2.5`$ are shown in fig. 3, where the finite size dependence is clearly seen. When the lattice volume is large enough, increasing the number of copies does not seem to make any substantial difference to our previously reported results for projected Creutz ratios and vortex densities. The details will be presented in a separate publication.
Finally, some questions about propagating ghosts in center gauges were raised at this meeting in the summary talk by Schierholz. Center gauges, like Landau gauge, are not ghost-free, and this would be a real problem if the aim of center gauge-fixing were to eliminate all unphysical modes in the Lagrangian. But the issue has little relevance, in our opinion, to the actual purpose of center gauge fixing, which is used in conjunction with center projection as a vortex finder. The rationale underlying this procedure, and its empirical success in finding vortices on thermalized lattices, have been discussed at length elsewhere. The same speaker questions whether vortex physics will be found to be consistent with instanton physics. We see no evidence of a problem in this area; in fact there are some suggestive findings to the effect that removing vortices from a lattice configuration also removes the topological charge. In any case, the study of instanton physics in relation to center vortices has only just begun, and there is no reasonable basis, at this stage, for strong conclusions.
In this talk I have indicated how the center vortex picture of confinement can be derived at strong lattice couplings, and why this picture is attractive at any gauge coupling. The strong-coupling analysis shows that center vortices are stabilized by color-screening terms in the long-range effective action, and that screening terms dominate the action at the adjoint string-breaking scale. We expect that these qualitative features of the strongly coupled gauge theory are also found in the continuum $`\beta \mathrm{}`$ limit.
## Acknowledgments
I would like to thank Hideo Suganuma, Hiroshi Toki, and the other organizers of Confinement 2000 for inviting me to participate in this stimulating meeting. This work was supported by the US Dept. of Energy under Grant No. DE-FG03-92ER40711.
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# The Broad Band Spectrum of MXB 1728–34 Observed by BeppoSAX
## 1 Introduction
The recent broad band spectral analysis of neutron star Low Mass X-ray Binaries (hereafter NS) has shown that their spectra are often more complex than previously thought. In fact, similarly to accreting Black Holes (BH), these spectra could be described by a power law with high energy cutoff and a soft excess at 0.5–1 keV. The thermal Comptonization of these soft photons by hot electrons probably originates the power law. Moreover an emission line is usually present at $`6.4`$ keV interpreted as fluorescence from iron in low ionization states. The probable origin of this line is from the reprocessed emission from the accretion disk surface illuminated by the primary comptonized spectrum. In this case one would also expect the presence of a bump between 20 and 40 keV due to Compton reflection of the primary spectrum by the disk. Indeed this reflection bump has been observed in the spectra of some NSs (Barret et al. 1999; Piraino et al. 1999; Yoshida et al. 1993), usually with reflection amplitudes (i.e. the solid angle $`\mathrm{\Omega }/2\pi `$ subtended by the reflector as seen from the corona) lower than 0.3. In these cases a correlation has been observed between the photon index of the primary spectrum and the reflection amplitude of the reprocessed component (Zdziarski et al. 1999; Barret et al. 1999; Piraino et al. 1999), the same observed in Seyfert galaxies and galactic BHs. It was originally thought that the electron temperature in the scattering cloud is lower for the NSs than for the BHs (Tavani & Barret 1997; Zdziarski et al. 1998; Churazov et al. 1997). This is in agreement with the expectation that an extra cooling should be present in the NSs because of the soft photons emitted by the surface. This seems indeed true for those systems in which a high energy cutoff has been observed, although some LMXBs do not show any cutoff in their spectra up to $`100`$ keV (Barret et al. 1991; Harmon et al. 1996; Piraino et al. 1999).
This similarity with the spectra of the BHs suggests that the same emission mechanism and geometry operate in NSs. In this scenario the soft excess comes from the accretion disk. The comptonized component originates in a hot and optically thin cloud (corona), probably placed between the neutron star and the accretion disk. A fraction of this component is reprocessed by the surface of the accretion disk arising the reflection bump and the iron K<sub>α</sub> line. The weakness of these components suggests that the reflector subtends a small solid angle ($`\mathrm{\Omega }/2\pi <0.3`$) as seen from the corona. This is possible if the disk is truncated and its inner part is absent, or if the inner accretion disk is highly ionized (but see Done & $`\dot{\mathrm{Z}}`$ycki 1999). In this geometry, a decrease of the inner radius of the disk causes an increase of the solid angle subtended by the reflector and a steepness of the power law. This could be responsible of the observed correlation between the photon index of the power law component and the reflection amplitude.
An Advection Dominated Accretion Flow (ADAF, see Narayan & Yi 1995 for a description) has been proposed as the origin of the scattering corona in NSs as in BHs (see Barret et al. 1999 for a discussion on this argument). However there is an important difference between NSs and BHs, because the neutron stars have a solid surface and the advected energy should be released at the surface or in a boundary layer. This extra soft emission would enhance the Compton cooling and result in softer spectra. Barret et al. (1999) suggest that the boundary layer is optically thin and its emission is hard, even if a reprocessed component from the neutron star surface could also be present.
In this paper we concentrate on the broad band (0.1–100 keV) spectral analysis performed on a BeppoSAX observation of MXB 1728–34 (4U 1728–34). MXB 1728–34 is a Low Mass X-ray Binary (LMXB), belonging to the class of atoll sources (Hasinger and van der Klis 1989). The optical counterpart has not been identified yet, due to the high optical extinction in the galactic center direction. The distance to this source is poorly known. It is probably between 4 and 14 kpc (Grindlay & Hertz 1981).
Observations with the Rossi X-Ray Timing Explorer (RXTE) have shown kilohertz quasi-periodic oscillations (kHz QPO) in the persistent emission of MXB 1728–34, with frequencies ranging from a few hundred Hz up to $`1200`$ Hz (Strohmayer et al. 1996). Usually two kHz peaks are simultaneously observed, with a nearly constant difference between their frequencies. If the frequency upper kHz QPO is interpreted as the Keplerian frequency at the inner rim of the accretion disk (e.g. Miller, Lamb, & Psaltis 1998, Stella & Vietri 1999), this implies that the disk can arrive very close to the neutron star in these systems. MXB 1728–34 shows frequent type-I X-ray bursts (Basinska et al. 1984). Recently bursts from this source were observed by RXTE (Strohmayer et al. 1997), and nearly coherent oscillations were discovered during the rising phase of the bursts, at a frequency of 363 Hz. The oscillation amplitude decreased as the X-ray burst flux increased. These oscillations likely result from spin modulation of a not uniform X-ray flux produced by the thermonuclear burning that causes the burst. The frequency separation between the two simultaneous kHz QPOs is similar to the frequency of the burst oscillations, suggesting a beat frequency mechanism as the origin of the kHz QPOs (Strohmayer et al. 1996, Miller et al. 1998). However precise measures of the peak separation between the kHz QPOs in this source have shown that it is always smaller than 363 Hz, and decreases significantly at higher inferred accretion rates (Méndez & van der Klis 1999).
The spectrum during the bursts is usually fitted by a blackbody. In the rising phase the blackbody temperature increases from $`1`$ to $`2.5`$ keV (Foster et al. 1986; Day & Tawara 1990), and, during the burst decay, decreases down to $`1`$ keV (Foster et al. 1986). There is some evidence of the presence of a hard tail during bursts (Foster et al. 1986). From the burst spectrum the radius of the emitting region can be derived and compared with the radius of the neutron star. Usually values well below 10 km are found (van Paradijs 1978, see also Basinska et al. 1984 and references therein), i.e. below the lowest radius allowed for a 1.4 $`M_{}`$ neutron star. Foster et al. (1986) showed that the measured color temperature of the blackbody can be higher than the effective temperature. Using appropriate correcting factors they derive a radius of the emitting region which is consistent with current equations of state, and a distance to the source of 4.3 kpc.
The persistent spectrum of MXB 1728–34 in the 1–20 keV energy range, observed by Einstein MPC, was fitted by a thermal bremsstrahlung with a temperature of $`18`$ keV (Grindlay and Hertz 1981). The spectrum in the same energy range was also obtained from SAS 3 data (Basinska et al. 1984). During the SAS 3 long observation ($`48`$ days) the source intensity (in the 1–20 keV energy range) was found to vary by a factor $`2`$. Accordingly the temperature associated to the spectrum was found to vary between 4 and 11 keV, positively correlated with the X-ray flux. The hard energy spectrum (30–200 keV), observed by SIGMA telescope (Claret et al. 1994), was also fitted by a thermal bremsstrahlung but with a higher temperature ($`38`$ keV). The 30–200 keV flux was $`7\times 10^{10}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$.
## 2 Observations
The Narrow Field Instruments (NFI) on board BeppoSAX satellite (Boella et al. 1997) observed MXB 1728–34 on 1998 August 23 and 24, for a total exposure time of 24 ks. The NFIs are four co-aligned instruments which cover more than three decades of energy, from 0.1 keV up to 200 keV, with good spectral resolution in the whole range. LECS (operating in the range 0.1–10 keV) and MECS (1–11 keV) have imaging capabilities with a Field of View (FOV) of $`20^{}`$ and $`30^{}`$ radius respectively. We selected the data for scientific analysis in circular regions of $`8^{}`$ and $`4^{}`$ radii for the LECS and MECS respectively centered on the source. The background subtraction was obtained using blank sky observations, in which we extracted the background spectra in regions of the FOV similar to those used for the source. MXB 1728–34 is near the Galactic plane and in this case simultaneous backgrounds should be used (Parmar et al. 1999). However MXB 1728–34 is an intense source and the background contributes to less than 1% and 0.3% of the total count rate in the LECS and MECS respectively. Therefore we preferred to use the standard method for background subtraction in this case. HPGSPC (7–60 keV) and PDS (13–200 keV) have no imaging capabilities, because the FOVs, of $`1^{}`$ FWHM, are delimited by collimators. The background subtraction for these instruments was obtained using the off-source data accumulated during the rocking of the collimators.
During the persistent emission the MECS light curve shows little variability of the intensity, the count rate ranging between 25 and 30 counts/s. We divided the MECS data in two energy bands, soft (1.8–4 keV) and hard (4–10 keV), and calculated the corresponding hardness ratio. This hardness ratio is constant implying no spectral variations during the BeppoSAX observation. The average observed flux of the source (in the 0.1–100 keV energy range) during the persistent emission is $`3.6\times 10^9`$ ergs cm<sup>-2</sup> s<sup>-1</sup>, which corresponds to an unabsorbed luminosity of $`1.2\times 10^{37}D_4^2`$ ergs/s, where $`D_4`$ is the distance to the source in units of 4 kpc.
During the BeppoSAX observation of MXB 1728–34, three type-I bursts are present. The duration of these bursts is between 10 and 15 s, considering the start and end times of each burst when the flux is 10% of the peak above the persistent emission level. The peak count rate (measured by the MECS) is around 500 counts/s for the first and the second burst, and $`300`$ counts/s for the third one. Figure 1 shows the MECS light curve (upper panel), in the energy band 1.8–10 keV and with the persistent level subtracted, of the first burst and the corresponding hardness ratio (4–10 keV/1.8–4 keV, lower panel). The hardness ratio seems to decrease from $`2.5`$ to $`1`$ during the burst decay. This means that the spectrum softens, as expected by the decay of the blackbody temperature during the burst. We obtain similar results for the other two bursts.
## 3 Spectral Analysis
The energy ranges used in the spectral analysis for each NFI are: 0.12–4 keV for the LECS, 1.8–10 keV for the MECS, 8–30 keV for the HPGSPC and 15–100 keV for the PDS. Different normalizations of the four NFIs are considered by including in the model normalizing factors, fixed to 1 for the MECS and kept free for the other instruments.
### 3.1 Spectral Analysis of the Burst Emission
We used the MECS and PDS data to produce the spectra during the bursts, considering the start and end times of each burst when the rate in the MECS is 10% of the peak above the persistent emission level. LECS and HPGSPC data were not used in this case because of the low statistics due to the short duration. With a preliminary fit we verified that the spectra of the three bursts presented the same shape. Therefore, to increase the statistics, we summed the three spectra obtained from the MECS. In the PDS, only the second and third bursts were observed and summed together. During the first burst, which falls at the beginning of a BeppoSAX orbit, the PDS was switched off. From the obtained burst spectrum we subtracted the spectrum of the persistent emission, and the result is shown in Figure 2 (upper panel). We fitted this spectrum with a blackbody modified by interstellar absorption from cold matter. Because of the lack of data at energies below 1.8 keV, we fixed the value of the equivalent hydrogen column N<sub>H</sub> at $`2.5\times 10^{22}`$ cm<sup>-2</sup>, in agreement with the value obtained from the fit of the BeppoSAX spectrum during the persistent emission (Table 2, model 4) and with previous results (Hoffman et al. 1979, Grindlay & Hertz 1981, Foster et al. 1986). In Table 1 we report the best fit values of the parameters and in Figure 2 (lower panel) the residuals (in units of $`\sigma `$) with respect to the best fit model. A hard excess between 20 and 30 keV is visible in these residuals with low statistical significance. This hard excess is still present when we restrict the MECS data to only the second and third bursts, that are simultaneously observed by MECS and PDS. Indeed, fitting only the PDS data to a blackbody we obtain a higher temperature, between 2.3 and 2.9 keV.
### 3.2 Spectral Analysis of the Persistent Emission
The broad band (0.1–100 keV) spectrum during the persistent emission was obtained excluding intervals of $`100`$ s around each burst, and is shown in Figure 3 (upper panel). A simple model (consisting of a single component, like a thermal bremsstrahlung with photoelectric absorption by cold matter) was not sufficient to fit the spectrum in the whole energy band.
We tried several models and we obtained a good representation of the spectrum using a soft emission that could be equivalently modeled by a multi-color disk blackbody (diskbb in XSPEC, Mitsuda et al. 1984) or a blackbody, a comptonized component fitted with the comptt model (described in Titarchuk, 1994), and two gaussian emission lines. The results of this fits are shown in Table 2 (model 1 and 2 respectively). The comptonized spectrum seems to be better fitted by the thComp model (see the appendix of Zdziarski et al. 1996 for a description of this model). The results of this fit are shown in Table 2, model 3 with the soft excess described by a multi-color disk blackbody, and model 4 with the soft excess described by a blackbody. The unfolded spectrum, together with the spectral components used to fit the data, are shown in Figure 4.
The soft emission could be described by a blackbody with a temperature of $`0.5`$ keV or a multi-color disk blackbody with an inner temperature $`0.8`$ keV. As it is evident from Table 2, the values of the parameters characterizing the other components do not vary significantly adopting a multi-color disk blackbody or a blackbody to fit the soft excess. For the comptonized spectrum, we obtain an electron temperature of $`711`$ keV and an optical depth of $`\tau 5`$ for a spherical scattering cloud, with a temperature of the seed photons $`1.5`$ keV. We used a broad gaussian to fit an iron K<sub>α</sub> line at $`6.7`$ keV, with equivalent width of 50–70 eV. We also used a narrow gaussian line to fit an excess in the residuals around 1.7 keV. The addition of a gaussian emission line at this energy is statistically significant at more than 99.99% confidence level. The residuals with respect to model 4 are shown in Figure 3 (lower panel). We also tried the presence of an iron edge at 7–8 keV, but the addition of this component to the model does not improve the fit.
The addition of a reflection component (pexriv in XSPEC, Magdziarz & Zdziarski 1995) does not improve significantly the fit. The upper limit on the reflection amplitude, corresponding to the solid angle (in units of $`2\pi `$) that the reflector subtends as viewed from the corona, is $`0.7`$, using an inclination angle $`\mathrm{cos}i=0.5`$. However the pexriv model does not contain the iron emission line that should be present in the reflection component. Therefore we tried a reflection model (that we shall call fe-refl) in which the iron line is self-consistently calculated for the given ionization state, temperature, spectral shape and abundances (this model is described in detail in $`\dot{\mathrm{Z}}`$ycki et al. 1998). Excluding the gaussian line at $`6.7`$ keV and adding this self-consistent reflection model (with the inclination angle $`\mathrm{cos}i`$ fixed to 0.5), we obtain a fit equivalent to that shown in Table 2 (model 4), with $`\chi ^2=663`$ for 614 degrees of freedom. The observed iron line requires, to be fitted, reflection from a ionized disk (with $`\xi =L_\mathrm{X}/n_\mathrm{e}r^2280`$) and a reflection amplitude $`\mathrm{\Omega }/2\pi =0.15\pm 0.08`$.
## 4 Discussion
We performed a spectral analysis of the bursts observed by BeppoSAX in the 1.8–70 keV energy range. We fitted the averaged spectrum during the bursts to a blackbody with a temperature $`2`$ keV. The total luminosity of this blackbody is $`3.6\times 10^{37}D_4^2`$ ergs/s, where $`D_4`$ is the distance to the source in units of 4 kpc. This corresponds to a radius of the emitting region of $`4.8D_4`$ km. Because the most probable distance to the source is not much different from 4 kpc (see Foster et al. 1986, and below), this radius is rather low to be identified with the neutron star radius. However the color temperature of the spectrum can be higher than the effective temperature by a factor 1.4 (Foster et al. 1986). Considering this correcting factor we find an effective radius of the emitting region of $`R_{\mathrm{eff}}9.4D_4`$ km, that is comparable with the radius of a neutron star.
The distance to the source is not well known and is estimated between 4 and 14 kpc (Grindlay & Hertz 1981). However the bursts of this source can be used to infer the distance, as shown by Basinska et al. (1984). They found that the maximum burst flux ($`F_{\mathrm{max}}`$) and the integrated burst flux ($`E_\mathrm{b}`$) are strongly correlated at low flux levels. For high flux levels $`F_{\mathrm{max}}`$ saturates at the value of $`7\times 10^8`$ ergs cm<sup>-2</sup> s<sup>-1</sup> (calculated in the 1–20 keV energy range). This is therefore a critical luminosity and might be interpreted as the Eddington limit luminosity. The bursts we see in MXB 1728–34 during the BeppoSAX observation are far from this saturation level. The nearest to saturation is the first burst, for which $`E_\mathrm{b}=3.3\times 10^7`$ ergs cm<sup>-2</sup> and $`F_{\mathrm{max}}=4.5\times 10^8`$ ergs cm<sup>-2</sup> s<sup>-1</sup> (corresponding to $`64\%`$ of the critical luminosity level), both calculated in the 1–20 keV energy range and with the persistent level subtracted as in Basinska et al. (1984). Using the fit parameters of Table 1, we can correct the maximum burst flux for absorption and extrapolate it in the 0.1–100 keV energy range. This corresponds to a total maximum luminosity of $`6.2\times 10^{36}D^2`$ ergs/s, where D is the distance in kpc. In the interpretation above, this should be equal to 64% of the Eddington luminosity, that is $`L_\mathrm{E}=2.5\times 10^{38}`$ ergs/s for a 1.4 $`M_{}`$ (van Paradijs & McClintock 1994), giving a distance to the source of $`5.1`$ kpc. This is similar to previous estimations based on the bursts (e.g. 4.3 kpc, Foster et al. 1986, and 4.2 kpc, van Paradijs 1978).
A rather complex model is needed to fit the persistent spectrum. The continuum consists of a soft emission that could be described by a blackbody or a multi-color disk blackbody plus a comptonized spectrum. This model is similar to that used for other type-I X-ray bursters, such as 1E 1724–3045 (Guainazzi et al. 1998, Barret et al. 1999) and KS 1731–260 (Barret et al. 1999). For these sources and MXB 1728–34, the electron temperature of the scattering region is between 3 and 28 keV and the optical depth between 3 and 10, with higher temperatures corresponding to lower optical depths, as expected. In the following we will describe the results we obtained for MXB 1728–34 in more detail.
The value of the equivalent absorption column $`N_\mathrm{H}`$ is 2.5–3 $`\times 10^{22}`$ cm<sup>-2</sup>, depending on the model used to fit the soft emission (blackbody or disk multi-temperature blackbody). If the distance of 5.1 kpc inferred from the bursts is correct, this implies a visual extinction of $`A_\mathrm{v}=7.8\pm 1.7`$ for MXB 1728–34 (calculated from Hakkila et al. 1997). Using the observed correlation between visual extinction and absorption column (Predehl & Schmitt 1995) we find $`N_\mathrm{H}=(1.4\pm 0.3)\times 10^{22}`$ cm<sup>-2</sup>, which is around half the value we find from the spectral fit. Large values of the absorption column, between 1.5 and 3.5 $`\times 10^{22}`$ cm<sup>-2</sup> and even larger, were also found in previous observations of MXB 1728–34 (see Basinska et al. 1994). Therefore there is an excess of attenuation of the soft emission that has to be caused by matter close to the X-ray source, e.g. by the outer cold regions of an accretion disk or corona.
Using a blackbody to describe the soft emission, we get a temperature of $`0.5`$ keV. The luminosity of this component is $`3\times 10^{36}D_4^2`$ ergs/s. This implies a radius of the spherical emitting region of $`15D_4`$ km. Indeed the distance to the source could be higher than the 4 kpc (Grindlay & Hertz 1981) and therefore the luminosity and the radius of the blackbody emission region could be higher. For instance, adopting a distance of 5.1 kpc, the calculated radius is $`20`$ km. Therefore the radius of the blackbody emitting region could be rather large to be identified with the neutron star surface. Using a multi-color disk blackbody to model the soft excess we obtain an inner disk temperature of $`0.8`$ keV, and a lower limit for the inner radius of the disk of $`R_{\mathrm{in}}\sqrt{\mathrm{cos}i}8D_4`$ km. This source does not show dips in the light curve and therefore the inclination of the normal to the plane of the disk with respect to the line of sight should be less than $`60^{}`$. Assuming an inclination angle $`i`$ of $`60^{}`$ we obtain $`11D_4`$ km for the inner radius of the accretion disk. Moreover the simple multi-color disk blackbody could be not appropriate to describe the emission of accretion disks in X-ray binaries, because the electron scattering will modify the spectrum (Shakura & Sunyaev 1973; White, Stella & Parmar 1988). In this case, the measured color temperature is related to the effective temperature of the inner disk $`T_{\mathrm{col}}=fT_{\mathrm{eff}}`$, where $`f`$ is the spectral hardening factor. The factor $`f`$ has been estimated by Shimura & Takahara (1995) to be $`1.7`$ for a luminosity $`10\%`$ of the Eddington limit, with a little dependence on the mass of the compact object and the radial position. Applying this correction to the values of $`T_{\mathrm{in}}`$ and $`R_{\mathrm{in}}\sqrt{\mathrm{cos}i}`$ reported in Table 2, we obtain an effective temperature of $`0.5`$ keV, which corresponds to an inner radius $`R_{\mathrm{eff}}\sqrt{\mathrm{cos}i}=f^2R_{\mathrm{in}}\sqrt{\mathrm{cos}i}20D_4`$ km.
In conclusion, we identify the soft component with the emission from the accretion disk, although we obtain a slightly better fit using the blackbody instead of the multi-color disk blackbody, especially when we describe the Comptonized spectrum with thComp (compare model 3 and model 4 in Table 2). In the hypothesis that the blackbody represents the emission from the inner accretion disk, we can calculate the inner radius of the disk from this component, in order to compare the result with the estimate obtained using the multi-color disk blackbody for the soft excess. The blackbody luminosity is then the total potential energy that has been released at $`R_{\mathrm{in}}`$:
$$L_{\mathrm{BB}}=\frac{GM\dot{M}}{2R_{\mathrm{in}}}$$
(1)
where we assumed a Keplerian, geometrically thin, optically thick accretion disk. We attribute the measured temperature of the blackbody to the maximum temperature in the disk, that is reached close to the inner radius at $`R=49/36`$ $`R_{\mathrm{in}}`$ (see e.g. Frank, King, & Raine 1985).
$$T_{\mathrm{max}}=0.488\left(\frac{3GM\dot{M}}{8\pi \sigma R_{\mathrm{in}}^3}\right)^{1/4}$$
(2)
where we assumed that the disk is truncated at $`R_{\mathrm{in}}`$ and that the zero-torque condition can be applied at $`R_{\mathrm{in}}`$. Using these two equations we can find the inner radius of the disk as a function of the measured temperature and luminosity of the blackbody:
$$R_{\mathrm{in}}=3.69T_{\mathrm{keV}}^2L_{37}^{1/2}\mathrm{km}$$
(3)
Considering the values reported in Table 2 (model 4), we obtain $`R_{\mathrm{in}}=6.3D_4`$ km for MXB 1728–34, that becomes $`18D_4`$ km considering a spectral hardening factor of $`1.7`$. This value is similar to the result obtained with the multi-color disk blackbody.
The second component of our model is a comptonized spectrum, produced in a hot ($`kT_\mathrm{e}711`$ keV) region of spherical shape and moderate optical depth ($`\tau 5`$) surrounding the neutron star. This comptonized spectrum requires soft seed photons at a temperature of $`1.41.5`$ keV, that is significantly higher than the temperature of the disk. The emission from the neutron star surface and/or from an optically thick or thin boundary layer (see the discussion in Barret et al. 1999 on the possible existence of an optically thin boundary layer) could then account for the seed photons subsequently comptonized in the surrounding corona. This could explain why we do not observe directly the emission from the neutron star and is in agreement with our identification of the blackbody as the emission from the accretion disk.
Adopting the geometry described above one would expect a Compton reflection by the accretion disk of the hard spectrum emitted by the corona. This would produce a bump in the spectrum, that can be described by the pexriv model. This reflection is not significantly detected in MXB 1728–34 (nor in the similar sources 1E 1724–3045, Guainazzi et al. 1998, Barret et al. 1999, and KS 1731–260, Barret et al. 1999). A possible reason could be the low temperature of the spectrum emitted by the corona. Because most of the hard photons have energies below 10 keV, they will be photoelectric absorbed by the matter in the disk, and re-emitted at the temperature of the disk, rather than Compton scattered. This explains why we obtained a rather high upper limit in the reflection amplitude ($`0.7`$) using pexriv. It is probable that the reflection component can be observed in a harder state of this source.
This argument could also be applied to KS 1731–260, whose spectrum is indeed soft, with an electron temperature of $`2.8`$ keV (Barret et al. 1999), but shows some difficulties when applyed to 1E 1724–3045, whose spectrum is harder with an electron temperature $`28`$ keV (Guainazzi et al. 1998, Barret et al. 1999). In this case Barret et al. (1999) suggest that the inclination of the source could be high and this could reduce the reflection component. We note that up to date the reflection component has been detected only in sources with very hard spectra, such as GS 1826–238 (with a cutoff energy $`90`$ keV, Barret et al. 1999), SLX 1735–269 (with a cutoff energy $`200`$ keV, Barret et al. 1999), 4U 0614+09 (with a cutoff energy $`>200`$ keV, Piraino et al. 1999) and 4U 1608–52 (with a cutoff energy at 300 keV, Yoshida et al. 1993).
The irradiation of the accretion disk by the X-ray primary spectrum should also produce an iron emission line at $`6.4`$ keV. Therefore we can better constrain the parameters of the reflection component with the fe-refl model, which includes a self-consistent calculation of this line. Using such a model, instead of the gaussian, to fit the iron line, we find that the reflector has to be ionized ($`\xi 280`$) and we derive an upper limit on the reflection amplitude of $`0.2`$, much stringent than the value found using pexriv.
Two emission lines, at 6.7 keV and 1.66 keV, are needed to fit the spectrum of MXB 1728–34. The line at $`6.7`$ keV is interpreted as emission from highly ionized iron. A broad emission line at 6.7 keV is compatible with being emitted in a strongly ionized region such as the corona or the inner part of the disk. The second emission line, at $`1.66`$ keV, is compatible with the radiative recombination emission from Mg XI. These X-ray emission lines at low energies likely arise in a photoionized corona (Liedahl et al. 1992, see White, Nagase & Parmar 1995 as a review).
## 5 Conclusions
We analysed data from a BeppoSAX observation of MXB 1728–34 performed in 1998 August between 23 and 24. Three type-I X-ray bursts were present during this observation. The spectrum during the bursts is fitted by a blackbody with a temperature $`2`$ keV. We calculated the corresponding effective radius of the blackbody emitting region, that is $`R_{\mathrm{eff}}9.4D_4`$ km, comparable with the neutron star radius. Assuming that the peak flux vs. fluence relationship of the bursts saturates at the Eddington limit luminosity, we find a distance to the source of 5.1 kpc.
The spectrum during the persistent emission is described by a blackbody with a temperature of $`0.5`$ keV, a comptonized spectrum, and two gaussian emission lines at 1.66 keV and 6.7 keV respectively. In the hypothesis that the blackbody is emitted by the accretion disk we estimate an inner disk radius of $`20`$ km. The comptonized spectrum is probably produced in a spherical, hot ($`10`$ keV) region of moderate optical depth ($`\tau 5`$) surrounding the neutron star. The soft seed photons for the Comptonization, with a temperature $`1.5`$ keV, probably come from the neutron star surface and/or boundary layer. The presence of a reflection component is not detected with high statistical significance. To constrain the reflection parameters we used a self-consistent model that also calculates the iron line produced in the reflection. This is important because MXB 1728–34 was in a soft state during the BeppoSAX observation and the low temperature of the primary emission significantly reduce the reflection bump between 20 and 40 keV. Using this self-consistent model to fit the iron line, we find that the reflector is ionized and the reflection amplitude is $`0.2`$. Alternatively the iron line could be emitted in a strongly ionized corona. Further observations are needed to clarify the origin of the iron line.
The authors want to thank P.T. $`\dot{\mathrm{Z}}`$ycki and C. Done for kindly supplying their reflection model fe-refl and the Comptonization model thComp, and for useful discussions. This work was supported by the Italian Space Agency (ASI), by the Ministero della Ricerca Scientifica e Tecnologica (MURST).
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# Griffiths-McCoy singularities in random quantum spin chains: Exact results through renormalization
## Abstract
The Ma-Dasgupta-Hu renormalization group (RG) scheme is used to study singular quantities in the Griffiths phase of random quantum spin chains. For the random transverse-field Ising spin chain we have extended Fisher’s analytical solution to the off-critical region and calculated the dynamical exponent exactly. Concerning other random chains we argue by scaling considerations that the RG method generally becomes asymptotically exact for large times, both at the critical point and in the whole Griffiths phase. This statement is checked via numerical calculations on the random Heisenberg and quantum Potts models by the density matrix renormalization group method.
In a random quantum system at zero temperature several physical quantities are singular not only at the critical point, but in a whole region, as well, which extends on both sides of the transition point. In this Griffiths phase a random quantum system is non-critical in the space direction (spatial correlations decay exponentially), whereas it is critical in the time direction and the corresponding behavior due to Griffiths-McCoy singularities is controlled by a line of semi-critical fixed points characterized by the dynamical exponent $`z(\delta )`$, which depends on the value of the quantum control-parameter, $`\delta `$. For example average autocorrelations decay in (imaginary) time, $`\tau `$, as $`\tau ^{z(\delta )}`$; for a small magnetic field, $`H0`$, the magnetization behaves as $`H^{1/z(\delta )}`$; the low temperature susceptibility and specific heat are singular as $`T^{1/z(\delta )1}`$ and $`T^{1/z(\delta )}`$, respectively. Let us mention that Griffiths-McCoy singularities are relevant in experiments on quantum spin glasses and provides a theoretical explanation about the non-Fermi liquid behavior in U and Ce intermetallics.
Among the theoretical methods developed to study random quantum systems the renormalization group (RG) scheme introduced by Ma, Dasgupta and Hu plays a special rôle. For a class of systems, the critical behavior of those is controlled by an infinite-randomness fixed point (IRFP), the RG method becomes asymptotically exact during iteration. For some one-dimensional models, random transverse-field Ising model (RTIM) and the random XXZ-model, Fisher has obtained analytical solution of the RG equations and in this way many new exact results and new physical insight about the critical behavior of these models have been gained. Subsequent analytical and numerical investigations of the models are in agreement with Fisher’s results. The RG-scheme has been numerically implemented in higher dimensions, as well, to study the critical behavior of the RTIM and reasonable agreement with the results of quantum Monte-Carlo simulations has been found. Considering the Griffiths-phase of random quantum spin chains here the RG-scheme has been rarely used, mainly due to the general belief that the method looses its asymptotically exact properties by leaving the vicinity of the scale invariant critical point.
Our aim in the present Letter is to clarify the applicability of the Ma-Dasgupta-Hu RG-method in the Griffiths phase of random quantum spin chains. We start with the RTIM, extend Fisher’s calculation to the Griffiths phase of the model and present the analytical solution of the RG-equations. Then, for general models, we analyze by scaling considerations the structure of the RG equations around the line of semi-critical fixed points and arrive to the conclusion that the RG method becomes asymptotically exact in the whole Griffiths region. This statement is then checked numerically on the random Heisenberg chain and the random quantum Potts model (RQPM) using the density matrix renormalization group (DMRG) method.
We start with the 1d RTIM which is defined by the Hamiltonian:
$$H_I=\underset{i}{}J_i\sigma _i^x\sigma _{i+1}^x\underset{i}{}h_i\sigma _i^z,$$
(1)
where the $`\sigma _i^{x,z}`$ are Pauli matrices at site $`i`$ and $`J_i`$ and $`h_i`$ are the couplings and the transverse fields, respectively, which are independent random variables. The quantum control-parameter of the model is defined as
$$\delta =\frac{[\mathrm{ln}h]_{\mathrm{av}}[\mathrm{ln}J]_{\mathrm{av}}}{\mathrm{var}[\mathrm{ln}h]+\mathrm{var}[\mathrm{ln}J]},$$
(2)
where $`\mathrm{var}[x]`$ is the variance of $`x`$ and we use $`[\mathrm{}]_{\mathrm{av}}`$ to denote averaging over quenched disorder. For $`\delta >0`$ ($`\delta <0`$) the system is in the paramagnetic (ferromagnetic) phase, so that the random quantum critical point is at $`\delta =0`$.
In the Ma-Dasgupta-Hu RG-method the strongest term in the Hamiltonian, coupling or field, of strength $`\mathrm{\Omega }`$ are successively decimated out and the neighboring fields or couplings are replaced by weaker ones, which are generated by a perturbation calculation. The basic RG-equations for coupling and field decimations are given by:
$$\stackrel{~}{h}=\frac{h_ih_{i+1}}{J_i\kappa },\stackrel{~}{J}=\frac{J_{i1}J_i}{h_i\kappa },$$
(3)
respectively, which are related through duality. Here, for the RTIM $`\kappa =1`$. Under renormalization we follow the probability distributions of the couplings, $`R(J,\mathrm{\Omega })`$, and that of the fields, $`P(h,\mathrm{\Omega })`$. When the energy-scale is lowered as $`\mathrm{\Omega }\mathrm{\Omega }d\mathrm{\Omega }`$ the distribution of the couplings is changed as:
$`{\displaystyle \frac{\mathrm{d}R(J,\mathrm{\Omega })}{\mathrm{d}\mathrm{\Omega }}}`$ $`=`$ $`R(J,\mathrm{\Omega })[P(\mathrm{\Omega },\mathrm{\Omega })R(\mathrm{\Omega },\mathrm{\Omega })]`$ (4)
$`P(`$ $`\mathrm{\Omega }`$ $`,\mathrm{\Omega }){\displaystyle }_{J\kappa }^\mathrm{\Omega }\mathrm{d}J^{}R(J^{},\mathrm{\Omega })R({\displaystyle \frac{J\mathrm{\Omega }\kappa }{J^{}}},\mathrm{\Omega }){\displaystyle \frac{\mathrm{\Omega }\kappa }{J^{}}},`$ (5)
where the first term in the r.h.s. is due to a balance between decimated couplings and normalization, whereas the second term accounts for the generated new couplings. A similar equation for the field distribution follows from Eq.(5) through duality
$`{\displaystyle \frac{\mathrm{d}P(h,\mathrm{\Omega })}{\mathrm{d}\mathrm{\Omega }}}`$ $`=`$ $`P(h,\mathrm{\Omega })[R(\mathrm{\Omega },\mathrm{\Omega })P(\mathrm{\Omega },\mathrm{\Omega })]`$ (6)
$`R(`$ $`\mathrm{\Omega }`$ $`,\mathrm{\Omega }){\displaystyle }_{h\kappa }^\mathrm{\Omega }\mathrm{d}h^{}P(h^{},\mathrm{\Omega })P({\displaystyle \frac{h\mathrm{\Omega }\kappa }{h^{}}},\mathrm{\Omega }){\displaystyle \frac{\mathrm{\Omega }\kappa }{h^{}}},`$ (7)
which amounts to interchange $`Jh`$ and $`RP`$.
For the RTIM, i.e. for $`\kappa =1`$, we found one class of solutions of the RG equations in the form:
$`R(J,\mathrm{\Omega })`$ $`=`$ $`R(\mathrm{\Omega },\mathrm{\Omega })\left(\mathrm{\Omega }/J\right)^{1R(\mathrm{\Omega },\mathrm{\Omega })\mathrm{\Omega }}`$ (8)
$`P(h,\mathrm{\Omega })`$ $`=`$ $`P(\mathrm{\Omega },\mathrm{\Omega })\left(\mathrm{\Omega }/h\right)^{1P(\mathrm{\Omega },\mathrm{\Omega })\mathrm{\Omega }},`$ (9)
where the distributions involve the parameters, $`\stackrel{~}{R}(\mathrm{\Omega })R(\mathrm{\Omega },\mathrm{\Omega })`$ and $`\stackrel{~}{P}(\mathrm{\Omega })P(\mathrm{\Omega },\mathrm{\Omega })`$, which satisfy the relation $`(\stackrel{~}{P}\stackrel{~}{R})\mathrm{\Omega }=1/z=const`$. Thus the solution is characterized by one parameter, $`z=z(\delta )`$, which depends on the quantum control-parameter, $`\delta `$: at the critical point, $`1/z(0)=0`$, whereas in the paramagnetic phase, $`\delta >0`$, $`1/z(\delta )>0`$ and monotonically increases with $`\delta `$.
In terms of the variables, $`y=\stackrel{~}{R}\mathrm{\Omega }+1/2z=\stackrel{~}{P}\mathrm{\Omega }1/2z`$ and $`x=\mathrm{ln}\mathrm{\Omega }`$ we obtain the differential equation:
$$\frac{\mathrm{d}y}{\mathrm{d}x}+y^2=\frac{1}{4z^2},$$
(10)
which has the solution, $`y=1/(xx_0)`$, $`x_0=\mathrm{const}`$, at the critical point with $`1/z=0`$. The distribution of $`\rho =R\mathrm{\Omega }`$ in terms of the variable $`\eta =(\mathrm{ln}\mathrm{\Omega }\mathrm{ln}J)/\mathrm{ln}\mathrm{\Omega }`$ is given by $`\rho (\eta )\mathrm{d}\eta =\mathrm{exp}(\eta )\mathrm{d}\eta `$, which corresponds to the fixed-point solution by Fisher. At this point we refer to Fisher’s analysis and conclude that the functions in Eqs(9) indeed represent the fixed-point distribution for all non-singular initial distributions.
In the Griffiths phase, $`\delta >0`$, the solution of Eq.(10) in terms of the original energy-scale variable, $`\mathrm{\Omega }`$, is given by
$$y=\frac{y_0/2z+1/4z^2\mathrm{tanh}\left[\mathrm{ln}(\mathrm{\Omega }_0/\mathrm{\Omega })/2z\right]}{1/2z+y_0\mathrm{tanh}\left[\mathrm{ln}(\mathrm{\Omega }_0/\mathrm{\Omega })/2z\right]},$$
(11)
where $`y=y_0`$ at a reference point, $`\mathrm{\Omega }=\mathrm{\Omega }_0`$. Approaching the line of semi-critical fixed points, i.e. for $`\mathrm{\Omega }/\mathrm{\Omega }_00`$, we have to leading order:
$$\stackrel{~}{R}\mathrm{\Omega }=\stackrel{~}{P}\mathrm{\Omega }1/z=\stackrel{~}{R}(\mathrm{\Omega }_0)/[\stackrel{~}{P}(\mathrm{\Omega }_0)z]\left(\mathrm{\Omega }/\mathrm{\Omega }_0\right)^{1/z}+\mathrm{}$$
(12)
thus $`\stackrel{~}{P}`$ and $`\stackrel{~}{R}`$ have different low energy asymptotics.
The physical relevance of $`1/z`$ can be obtained by studying the change of number of spins, $`n_\mathrm{\Omega }n_\mathrm{\Omega }\mathrm{d}n_\mathrm{\Omega }`$ connected with a change in the energy scale as $`\mathrm{\Omega }\mathrm{\Omega }\mathrm{d}\mathrm{\Omega }`$. This leads to the differential equation $`\mathrm{d}n_\mathrm{\Omega }/\mathrm{d}\mathrm{\Omega }=n_\mathrm{\Omega }\left[P(\mathrm{\Omega },\mathrm{\Omega })+R(\mathrm{\Omega },\mathrm{\Omega })\right]`$, the solution of which is given by:
$$n_\mathrm{\Omega }=\left\{\mathrm{cosh}\left[\mathrm{ln}(\mathrm{\Omega }_0/\mathrm{\Omega })/2z\right]+2zy_0\mathrm{sinh}\left[\mathrm{ln}(\mathrm{\Omega }_0/\mathrm{\Omega })/2z\right]\right\}^2,$$
(13)
which along the line of semi-critical fixed points has the asymptotic behavior $`n_\mathrm{\Omega }=const\mathrm{\Omega }^{1/z}`$, $`\mathrm{\Omega }0`$. Since the typical distance between remaining spins is $`L_\mathrm{\Omega }1/n_\mathrm{\Omega }\mathrm{\Omega }^{1/z}`$, we can identify $`z`$ as the dynamical exponent, which governs the relation between time- and length-scales as $`\tau L^z`$.
Next we show that $`z`$ is invariant along the RG trajectory and can be deduced from the original distributions. For this we consider the averages, $`[J^\mu ]_{\mathrm{av}}`$ and $`[h^\mu ]_{\mathrm{av}}`$, and using Eqs.(5) and (7) we calculate the derivative:
. $`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}\mathrm{\Omega }}}\left[\left(J/h\right)^\mu \right]_{\mathrm{av}}=\left(1\left[\left(J/h\right)^\mu \right]_{\mathrm{av}}\right)`$ (14)
$`\times `$ $`\left(P(\mathrm{\Omega },\mathrm{\Omega })\mathrm{\Omega }^\mu [J^\mu ]_{\mathrm{av}}+R(\mathrm{\Omega },\mathrm{\Omega })\mathrm{\Omega }^\mu [h^\mu ]_{\mathrm{av}}\right)`$ (15)
which is vanishing for $`\mu =\stackrel{~}{\mu }`$, if $`[(J/h)^{\stackrel{~}{\mu }}]_{\mathrm{av}}=1`$. Consequently $`\stackrel{~}{\mu }`$ stays invariant along the RG trajectory until the fixed point, where using the distribution in Eqs.(9) we obtain $`\stackrel{~}{\mu }=1/z`$. Thus the dynamical exponent for the RTIM is given by the solution of the equation:
$$\left[\left(J/h\right)^{1/z}\right]_{\mathrm{av}}=1,$$
(16)
which is then exact, since the RG-transformation becomes asymptotically exact as $`\mathrm{\Omega }0`$. This latter statement follows from the fact that the ratio of decimated bonds, $`\mathrm{\Delta }n_J`$, and decimated fields, $`\mathrm{\Delta }n_h`$, goes to zero as $`\mathrm{\Delta }n_J/\mathrm{\Delta }n_h=R(\mathrm{\Omega },\mathrm{\Omega })/P(\mathrm{\Omega },\mathrm{\Omega })\mathrm{\Omega }^{1/z}`$. Then the probability, $`Pr(\alpha )`$, that the value of a coupling, $`J`$, being neighbor to a decimated field is $`\mathrm{\Omega }<J<\alpha \mathrm{\Omega }`$ with $`0<\alpha <1`$ is given by $`Pr(\alpha )=1\alpha ^{\stackrel{~}{R}\mathrm{\Omega }}`$, which goes to zero for any non-zero $`\alpha `$, since $`\stackrel{~}{R}\mathrm{\Omega }\mathrm{\Omega }^{1/z}`$ at the fixed point. Consequently the decimations in Eq.(3) and the related RG equation in Eqs.(5) and (7) are indeed exact.
The relation in Eq.(16) can be recovered starting from the exact expression for the surface magnetization, $`m_s`$, on a chain with a fixed spin at site $`L+1`$ :
$$\frac{1}{m_s^2}=1+\underset{l=1}{\overset{L}{}}\underset{i=1}{\overset{l}{}}\left(\frac{h_i}{J_i}\right)^2.$$
(17)
This type of expression is the so-called Kesten-variable in the mathematical literature and the corresponding probability distribution is singular at $`m_s=0`$ in the thermodynamic limit for $`[\mathrm{ln}J]_{\mathrm{av}}>[\mathrm{ln}h]_{\mathrm{av}}`$. Then for the distribution of $`\mathrm{ln}m_s`$ we have the following singularity:
$$P(\mathrm{ln}m_s)m_s^{1/\stackrel{~}{z}},m_s0,$$
(18)
where $`\stackrel{~}{z}`$ is the solution of Eg.(16) in terms of the dual variables. The physical origine of the small $`m_s`$ tail of the distribution is due to such samples which have a weakly coupled domain (WCD), which effectively cuts the system inty two very weakly interacting parts and thus reduces the surface order enormously. In the dual system in the paramagnetic phase the dual object to a WCD is a strongly coupled domain (SCD) which results in a very small energy gap, $`ϵ`$. Thus in the tails of the distributions, $`m_s`$ and $`ϵ`$ are dual quantities. This remark indicates that the dynamical exponent found by the RG calculation in Eq.(16) is indeed exact.
Next, we consider general random quantum spin chains with a critical IRFP and analyze the structure of the RG equations close to the line of semi-critical fixed points, thus as $`\mathrm{\Omega }0`$. As for the RTIM, the decimation for fields and couplings is asymmetric and for $`\mathrm{\Omega }0`$ exclusively fields are decimated out, which are typically infinitly stronger, than the neighboring couplings. Therefore the RG decimation equations in Eq.(3) are asymptotically exact. The second point is to show that the dynamical exponent stays invariant along the RG trajectory, even though in the starting phase the RG equations are approximative. For this we consider the low energy tail of the distribution function of the first gap, $`P_1(\mathrm{ln}ϵ)`$, which involves in analogy to Eq.(18) the exponent $`z`$, and use the scaling result of Ref. This states that the probability distribution of the second, third, etc. gaps are related to $`P_1(\mathrm{ln}ϵ)`$ as $`P_2(\mathrm{ln}ϵ)P_1^2(\mathrm{ln}ϵ)`$, $`P_3(\mathrm{ln}ϵ)P_1^3(\mathrm{ln}ϵ)`$, etc, due to the fact that for a small second, third gap one needs two, three independent SCD-s and the corresponding probabilities are multiplied. In the RG decimation the SCD-s are only eliminated through coupling decimation, since their couplings are stronger than the average fields. If at some time a SCD with a small gap, $`ϵ`$, is eliminated then in the probability distribution, $`P_1(\mathrm{ln}ϵ)`$, one should consider the former second gap and use the corresponding conditional probability, $`P_1(\mathrm{ln}ϵ)P_2(\mathrm{ln}ϵ)/P_1(\mathrm{ln}ϵ)P_1(\mathrm{ln}ϵ)`$. Thus the small energy tail of the gap-distribution and consequently the dynamical exponent remains invariant under the renormalization procedure. The previously obtained exact results for the RTIM give strong support for the validity of these phenomenological considerations.
For a numerical demonstration of the validity of the above statement we considered two random quantum spin chains, the dimerized Heisenberg (XXX) chain and the q-state RQPM, both having a set of RG equations very similar to that of the RTIM in Eq.(3). For the dimerized XXX chain $`J`$ and $`h`$ in Eq.(3) are replaced by the Heisenberg couplings at odd and even positions, $`J_o`$ and $`J_e`$, respectively, and the parameter takes the value $`\kappa =2`$. The distance from the critical point is measured similarly to Eq.(2). For the q-state RQPM fields and couplings play analogous rôle as for the RTIM, the quantum control-parameter is given in Eq.(2), whereas $`\kappa `$ takes the value $`\kappa =q/2`$. We note that the RG equations for the XXX-chain and the $`q=4`$ state RQPM are identical.
At the critical point the RG-equations for $`1\kappa <\mathrm{}`$ has been solved by Senthil and Majumdar with the result that $`\kappa `$ is an irrelevant variable and the IRFP is the same as for the RTIM. In the Griffiths-phase we could not find a complete solution of the RG equations, in spite of the close similarity to that of the RTIM. We could, however, show that up to an accuracy of $`O(\mathrm{\Omega }^{1/z})`$ the solution is of the form of Eqs.(9) and thus there is infinite randomness along the line of fixed points. The $`z`$ exponent, however, does depend on the parameter $`\kappa `$, since the validity of the condition in Eq.(16) is limited to $`\kappa =1`$, thus in general $`z=z_\kappa (\delta )`$.
We have calculated the dynamical exponent by a numerical implementation of the RG scheme over 50000 samples of length $`L2^{14}`$. Starting with the uniform probability distribution: $`R_0(J)=\mathrm{\Theta }(1J)\mathrm{\Theta }(J)`$, $`P_0(h)=\mathrm{\Theta }(h_0h)\mathrm{\Theta }(h)/h_0`$ (and analogously for the XXX-chain) we got the estimates shown in Fig.1: $`1/z_\kappa `$ is a monotonously decreasing function of $`\kappa `$ and eventuelly it goes to zero in the whole Griffiths phase in the limit $`\kappa \mathrm{}`$.
The dynamical exponents of the XXX-chain and the RQPM are also calculated directly from the asymptotic behavior of the distribution of the surface magnetization, as given in Eq.(18). For the numerical calculations of the surface magnetizations we used the DMRG method for rather large finite chains with $`L64`$ and considered some 20000 samples. We found an overall agreement between the dynamical exponents calculated by the two methods. As a demonstration we show in Fig. 2 the distribution of $`m_s`$ for the XXX-chain, compared with that of the $`q=4`$ state RQPM, where for both models we are at the same distance from the transition point. As seen in Fig. 2 the asymtotic behavior of the two distributions is identical, as expected on the RG basis, since $`\kappa =2`$ for both models. Furthermore the dynamical exponents agree very well with those calculated by the RG method.
To conclude, we have shown in this paper that for random quantum spin chains the RG method of Ma, Dasgupta and Hu is asymptotically exact for large times, i.e. along the line of semi-critical fixed points. Consequently the dynamical exponent in the Griffiths phase calculated by the RG-method is exact. This result will hopefully stimulate new efforts to solve the RG equations for different systems analytically in 1d also outside the critical point, and in higher dimensions to perform numerical calculations and get more precise estimates for the Griffiths-McCoy singularities.
Acknowledgment: This work has been supported by the Hungarian National Research Fund under grant No OTKA TO23642, TO25139, F/7/026004, MO28418 and by the Ministery of Education under grant No. FKFP 0596/1999. We are indebted to E. Carlon for his participation in the early stages of the DMRG calculations. F. I. is grateful to H. Rieger and L. Turban for previous cooperation in the project and for useful discussions.
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# Orbifold Hodge numbers of the wreath product orbifolds
## 1. Introduction
In the study of orbifold string theory, Dixon, Harvey, Vafa and Witten introduced the notion of orbifold Euler number for a smooth manifold $`Y`$ acted on by a finite group $`G`$ and raised the question on the existence of a resolution of singularities of $`M/G`$ whose ordinary Euler number coincides with the orbifold Euler number. The orbifold Euler number has subsequently been interpreted as the Euler number for equivariant $`K`$-theory, cf. Atiyah-Segal . The notion of orbifold Euler number has been further refined to give rise to the notion of orbifold Hodge numbers and more generally the stringy Hodge numbers . The orbifold Hodge numbers of an orbifold are then conjectured to coincide with the ordinary Hodge numbers of a suitable resolution of the orbifold. For recent related development, see and the references therein.
A well-known series of examples with such a property is provided by the symmetric product of a surface which admits a resolution of singularities given by the Hilbert scheme of points. In this case the orbifold Euler number calculated by Hirzebruch and Höfer matches with the Euler number of the Hilbert scheme found earlier by Göttsche . It is further shown by Göttsche that the orbifold Hodge numbers matches with the Hodge numbers of the Hilbert scheme calculated by Göttsche and Soergel (also see Cheah ). The same method has been used by the second author , also for the calculation for higher dimensional complex manifolds.<sup>1</sup><sup>1</sup>1Seeing the math review (99c:14022) of but not the paper itself when he was writing , the second author got the wrong impression that Göttsche proved his result by establishing the strong McKay correspondence for symmetric products. This misunderstanding has been clarified when we actually looked into the paper during the preparation of the present paper.
The wreath product orbifolds, which are generalizations of the symmetric products, were shown by the first author (also see ) to have deep connections with Hilbert schemes of surfaces and vertex representations of infinite dimensional Lie algebras. More explicitly, if $`Y`$ is a smooth manifold acted upon by a finite group $`G`$, then there exists a natural action on the $`n`$-th Cartesian product $`Y^n`$ by the wreath product $`G_n`$ (which is the semidirect product of the symmetric group $`S_n`$ and the product group $`G^n`$). The orbifold Euler number for $`Y^n/G_n`$ has been explicitly calculated in loc. cit.. If in addition we assume that $`Y`$ is a quasi-projective surface and $`X`$ is a resolution of singularities of the orbifold $`Y/G`$, then the following commutative diagram
$$\begin{array}{ccc}X^{[n]}& & X^n/S_n\\ & & & & \\ Y^n/G_n& \stackrel{}{}& (Y/G)^n/S_n\end{array}$$
implies that the Hilbert scheme $`X^{[n]}`$ is a resolution of singularities of the orbifold $`X^n/G_n`$. It has been shown that if the ordinary Euler number of $`X`$ equals the orbifold Euler number of $`Y/G`$ then the ordinary Euler number of $`X^{[n]}`$ equals the orbifold Euler number of $`X^n/G_n`$ for all $`n`$. When $`G`$ is the trivial group and $`X`$ equals $`Y`$, one recovers the case of symmetric products.
The purpose of the present paper is to point out that the wreath product orbifolds also provide a large class of new higher dimensional examples which verify the orbifold Hodge number conjecture. More precisely, we show that if $`Y`$ is a quasi-projective surface and $`X`$ is a resolution of singularities of $`Y/G`$ such that the ordinary Hodge numbers of $`X`$ coincide with the orbifold Hodge numbers of the orbifold $`Y/G`$, then the orbifold Hodge numbers of the orbifold $`Y^n/G_n`$ coincide with the Hodge numbers of the Hilbert scheme $`X^{[n]}`$, which is a resolution of singularities. Our proof rely on the analysis of fixed-point set structures of the wreath product action on $`Y^n`$ (cf. ). As in our calculation of the orbifold Hodge numbers for $`Y^n/G_n`$ actually works for any complex $`G`$-manifold of even dimension.
In a very recent paper , Bryan, Donagi and Leung pointed out a series of examples (besides the well-known symmetric products) verifying the orbifold Hodge numbers conjecture. It turns out that their examples correspond to our special case when $`Y`$ is an abelian surface, $`G`$ is $`_2`$, and $`X`$ is the Kummer $`K3`$ surface. They remarked that most examples in literature are lower dimensional and no other higher dimensional examples known to them. Their remarks are largely responsible for us to decide to write up the results on the wreath product orbifolds which have been known to us for some time. We do not know any other higher dimensional examples which verify the orbifold Hodge number conjecture.
To conclude we also make two explicit conjectures on elliptic genera for wreath product orbifolds. These are motivated and in turn generalize the work of Dijkgraaf, Moore, Verlinde and Verlinde on elliptic genera for symmetric products.
The layout of this paper is as follows. In section 2, we recall the fixed-point set structures of the wreath product action, and the definition of the orbifold Hodge numbers, for both the compact and noncompact situations. In section 3, we prove our main results, Theorem 3.1 on the orbifold Hodge numbers of wreath product orbifolds and Theorem 3.2 on the verification of orbifold Hodge number conjecture. In section 4, we provide various examples illustrating our main results and in addition formulate two conjectures on elliptic genera.
## 2. preliminaries on the wreath product and orbifold Hodge numbers
In this section, we first review the definition of a wreath product $`G_n`$ associated to a finite group $`G`$, and the descriptions of conjugacy classes and centralizers for $`G_n`$, cf. e.g. . We also describe fixed-point sets for the action of $`G_n`$ on the $`n`$-th Cartesian product of a $`G`$-manifold, following . We then recall the definition of orbifold Hodge numbers, cf. .
### 2.1. The wreath product action on $`Y^n`$
Let $`G`$ be a finite group and denote by $`G_{}`$ the set of conjugacy classes of $`G`$. Let $`G^n=G\times \mathrm{}\times G`$ be the direct product of $`n`$ copies of $`G`$. Denote by $`[g]`$ the conjugacy class of $`gG`$. The symmetric group $`S_n`$ acts on $`G^n`$ by permuting the $`n`$ factors: $`s(g_1,\mathrm{},g_n)=(g_{s^1(1)},\mathrm{},g_{s^1(n)})`$. The wreath product $`G_n=GS_n`$ is defined to be the semidirect product $`G^nS_n`$ of $`G^n`$ and $`S_n`$, namely the multiplication on $`G_n`$ is given by $`(g,s)(h,t)=(g.s(h),st)`$, where $`g,hG^n,s,tS_n`$. Note when $`G`$ is the trivial one-element group the wreath product $`G_n`$ reduces to $`S_n`$, and when $`G`$ is $`_2`$ the wreath product $`G_n`$ is the hyperoctahedral group, the Weyl group of type $`C`$.
Given $`a=(g,s)G_n`$ where $`g=(g_1,\mathrm{},g_n)`$, we write $`sS_n`$ as a product of disjoint cycles: if $`z=(i_1,\mathrm{},i_r)`$ is one of them, the cycle-product $`g_{i_r}g_{i_{r1}}\mathrm{}g_{i_1}`$ of $`a`$ corresponding to the cycle $`z`$ is determined by $`g`$ and $`z`$ up to conjugacy. For each $`cG_{}`$ and each integer $`r1`$, let $`m_r(c)`$ be the number of $`r`$-cycles in $`s`$ whose cycle-product lies in $`c`$. Denote by $`\rho (c)`$ the partition having $`m_r(c)`$ parts equal to $`r`$ ($`r1`$) and denote by $`\rho =(\rho (c))_{cG_{}}`$ the corresponding partition-valued function on $`G_{}`$. Note that $`\rho :=_{cG_{}}|\rho (c)|=_{cG_{},r1}rm_r(c)=n`$, where $`|\rho (c)|`$ is the size of the partition $`\rho (c)`$. Thus we have defined a map from $`G_n`$ to $`𝒫_n(G_{})`$, the set of partition-valued function $`\rho =(\rho (c))_{cG_{}}`$ on $`G_{}`$ such that $`\rho =n`$. The function $`\rho `$ or the data $`\{m_r(c)\}_{r,c}`$ is called the type of $`a=(g,s)G_n`$. Denote $`𝒫(G_{})=_{n0}𝒫_n(G_{})`$. It is well known (cf. e.g. ) that two elements in $`G_n`$ are conjugate to each other if and only if they have the same type.
Let us describe the centralizer $`Z_{G_n}(a)`$ of $`aG_n`$, cf. . First we consider the typical case that $`a`$ has one $`n`$-cycle. As the centralizers of conjugate elements are conjugate subgroups, we may assume that $`a`$ is of the form $`a=((g,1,\mathrm{},1),\tau )`$, where $`\tau =(12\mathrm{}n)`$. Denote by $`Z_G^\mathrm{\Delta }(g)`$, or $`Z_G^{\mathrm{\Delta }_n}(g)`$ when it is necessary to specify $`n`$, the following diagonal subgroup of $`G^n`$ (and thus a subgroup of $`G_n`$):
$`Z_G^\mathrm{\Delta }(g)=\{((h,\mathrm{},h),1)G^nhZ_G(g)\}.`$
The centralizer $`Z_{G_n}(a)`$ of $`a`$ in $`G_n`$ is equal to the product $`Z_G^\mathrm{\Delta }(g)a`$, where $`a`$ is the cyclic subgroup of $`G_n`$ generated by $`a`$.
Take a generic element $`a=(g,s)G_n`$ of type $`\rho =(\rho (c))_{cG_{}}`$, where $`\rho (c)`$ has $`m_r(c)`$ $`r`$-cycles ($`r1`$). We may assume (by taking a conjugation if necessary) that the $`m_r(c)`$ $`r`$-cycles are of the form
$$g_{ur}(c)=((g,1,\mathrm{},1),(i_{u1},\mathrm{},i_{ur})),1um_r(c),gc.$$
Denote $`g_r(c)=((g,1,\mathrm{},1),(12\mathrm{}r)).`$ Throughout the paper, $`_{c,r}`$ is understood as the product $`_{cG_{},r1}`$. The centralizer $`Z_{G_n}(a)`$ of $`aG_n`$ is isomorphic to a direct product of the wreath products
$`{\displaystyle \underset{c,r}{}}\left(Z_{G_r}(g_r(c))S_{m_r(c)}\right).`$
Furthermore $`Z_{G_r}(g_r(c))`$ is isomorphic to $`Z_G^{\mathrm{\Delta }_r}(g)g_r(c)`$.
For a $`G`$-space $`Y`$, we define an action of $`G_n`$ on $`Y^n`$ as follows. Given $`a=((g_1,\mathrm{},g_n),s)`$, we let
(1) $`a.(x_1,\mathrm{},x_n)=(g_1x_{s^1(1)},\mathrm{},g_nx_{s^1(n)})`$
where $`x_1,\mathrm{},x_nY`$.
Next we recall the description of the fixed point set $`(Y^n)^a`$ for $`aG_n`$, cf. . Let us first look at the typical case $`a=((g,1,\mathrm{},1),\tau )G_n`$. Note that the centralizer group $`Z_G(g)`$ preserves the $`g`$-fixed point set $`X^g`$. The fixed point set is
$`(Y^n)^a=\{(x,\mathrm{},x)Y^nx=gx\}`$
which can be naturally identified with $`Y^g`$. The action of $`Z_{G_n}(a)`$ on $`(Y^n)^a`$ can be identified canonically with that of $`Z_G(g)`$ on $`Y^g`$ together with the trivial action of the cyclic group $`a`$. Thus $`(X^n)^a/Z_{G_n}(a)`$ can be identified with $`X^g/Z_G(g).`$
All $`Z_G(g)`$ are conjugate and all $`X^g`$ are homeomorphic to each other for different representatives $`g`$ in a fixed conjugacy class $`cG_{}`$. Also the orbit space $`X^g/Z_G(g)`$ can be identified with each other by conjugation for different representatives of $`g`$ in $`cG_{}`$. We agree to denote $`Z_G(g)`$ (resp. $`Y^g`$, $`Y^g/Z_G(g)`$) by $`Z_G(c)`$ (resp. $`Y^c`$, $`Y^c/Z_G(c)`$) by abuse of notations. Similar remarks apply to other situations below when the choice of representatives in a conjugacy class is irrelevant.
For an element $`aG_n`$ of type $`\{m_r(c)\}`$, the fixed-point set $`(Y^n)^a`$ can be naturally identified with $`_{c,r}(Y^c)^{m_r(c)}`$. Furthermore the orbit space $`(Y^n)^a/Z_{G_n}(a)`$ can be naturally identified with
(2) $`{\displaystyle \underset{c,r}{}}S^{m_r(c)}\left(Y^c/Z_G(c)\right)`$
where $`S^m(X)`$ denotes the $`m`$-th symmetric product $`X^m/S_m`$.
### 2.2. Definition of orbifold Hodge numbers
Let $`Y`$ be a compact complex manifold of complex dimension $`d`$ acted on by a finite group $`G`$ of automorphisms. For each conjugacy class $`c=[g]G_{}`$, let $`Y_1^g,\mathrm{},Y_{N_c}^g`$ be the connected components of the fixed-point set $`Y^g`$. Zaslow defined a shift number $`F_\alpha ^g`$ associated to each component $`Y_\alpha ^g`$ as follows. On the tangent space to each point in $`Y_\alpha ^g`$, $`g`$ acts as a diagonal matrix $`\mathrm{diag}(e^{2\pi \sqrt{1}\theta _1},\mathrm{},e^{2\pi \sqrt{1}\theta _d})`$, where $`0\theta _i<1`$. Then
$$F_\alpha ^g=\underset{j=1}{\overset{d}{}}\theta _j.$$
In general, $`F_\alpha ^g`$ is just a rational number. However, there are many occasions when it is an integer, e.g., when $`g`$ acts on the tangent space by a matrix in $`SL(n,)`$.
###### Remark 2.1.
In the case when $`Y`$ is a complex surface, the shift $`F_\alpha ^g`$ is an integer only if the component $`Y_\alpha ^g`$ is either an isolated point or two dimensional. Indeed a finite subgroup $`G`$ of $`GL(2,)`$ acting on $`^2`$ has integer shifts if and only if $`G`$ lies in $`SL(2,)`$. That is, the shift $`F_\alpha ^g=\theta _1+\theta _2`$ is an integer if and only if $`detg=e^{2\pi \sqrt{1}(\theta _1+\theta _2)}=1.`$
In the case all the shifts are integers, the orbifold Hodge numbers of the orbifold $`Y/G`$ are defined to be
(3) $`h^{s,t}(Y,G)={\displaystyle \underset{cG_{}}{}}{\displaystyle \underset{\alpha _c=1}{\overset{N_c}{}}}h^{sF_{\alpha _c}^c,tF_{\alpha _c}^c}(Y_{\alpha _c}^c/Z(c)).`$
The ordinary Dolbeault cohomology for an orbifold is given by (cf. Satake )
(4) $`H^,(Y/G)H^,(Y)^G.`$
Clearly the orbifold Hodge numbers can now be regarded as the dimensions of the corresponding orbifold cohomology groups (cf. )
(5) $`H^,(Y,G)={\displaystyle \underset{cG_{}}{}}{\displaystyle \underset{\alpha _c=1}{\overset{N_c}{}}}H^,(Y_{\alpha _c}^c/Z_c)\{F_{\alpha _c}^c\}.`$
Here and below we adopt the convention that if $`V=_{s,t}V^{s,t}`$ is a bigraded vector space, then $`V\{n\}`$ is the bigraded vector space with $`(V\{n\})^{s,t}=V^{sn,tn}.`$
It is convenient to form the generating function of bigraded spaces
$$H(Y,G;x,y)=\underset{s,t}{}H^{s,t}(Y,G)x^sy^t,$$
whose graded dimension is given by the orbifold Hodge polynomial
$$h(Y,G;x,y)=\underset{s,t}{}h^{s,t}(Y,G)x^sy^t.$$
Then we can rewrite the definition of orbifold cohomology groups as
$`H(Y,G;x,y)`$ $`=`$ $`{\displaystyle \underset{cG_{}}{}}{\displaystyle \underset{\alpha _c=1}{\overset{N_c}{}}}H(Y_{\alpha _c}^c/Z_c;x,y)\{F_{\alpha _c}^c\}`$
$`=`$ $`{\displaystyle \underset{cG_{}}{}}{\displaystyle \underset{\alpha _c=1}{\overset{N_c}{}}}H(Y_{\alpha _c}^c/Z_c;x,y)(xy)^{F_{\alpha _c}^c}.`$
For later use we define the orbifold virtual Hodge polynomial
$$e(Y,G;x,y)=\underset{s,t}{}(1)^{s+t}h^{s,t}(Y,G)x^sy^t.$$
We also define the usual virtual Hodge polynomial for the Hodge numbers $`h^{s,t}(Y)`$ associated to smooth $`Y`$ by letting $`e(Y;x,y)=_{s,t}(1)^{s+t}h^{s,t}(Y)x^sy^t.`$
### 2.3. The definition of orbifold virtual Hodge numbers
We now indicate how to extend the above definitions to the case of smooth quasi-projective varieties by using Deligne’s theory of mixed Hodge structures . Recall that a (pure) Hodge structure of weight $`m`$ on a complex vector space $`H`$ with a real structure is a direct sum decomposition:
$$H=\underset{s+t=m}{}H^{s,t},$$
such that $`\overline{H}^{s,t}=H^{t,s}`$ for all pairs $`(s,t)`$. A mixed Hodge structure (MHS) on $`H`$ consists of two filtrations
$$0\mathrm{}W_{m1}W_mW_{m+1}\mathrm{}H,$$
the ‘weight filtration’, and
$$H\mathrm{}F^{p1}F^pF^{p+1}\mathrm{}0,$$
the ‘Hodge filtration’, such that the filtration induced by the latter on $`Gr_m(W_{})=W_m/W_{m1}`$ defines a Hodge structure of weight $`m`$, for each $`m`$.
Define
$$I^{s,t}=F^sW_{s+t}\left[\overline{F^t}W_{s+t}+\underset{i2}{}\overline{F^{ti+1}}W_{s+ti}\right];$$
Then $`I^{s,t}W_{s+t}`$ maps isomorphically to the $`(s,t)`$ component in $`Gr_{s+t}(W_{})`$. One can show that
$`F^s(H)`$ $`={\displaystyle \underset{s^{}s}{}}{\displaystyle \underset{t}{}}I^{s^{},t}(H),`$ $`W_m(H)`$ $`={\displaystyle \underset{s+tm}{}}I^{s,t}(H).`$
It can be shown that $`\{I^{s,t}\}`$ is a splitting of $`H`$ characterized by the property that
$$I^{s,t}\overline{I^{t,s}}\left(mod\underset{s^{}<s,t^{}<t}{}I^{s^{},t^{}}\right)$$
(cf. ). We will refer to this splitting as the canonical splitting. Define
$$h^{s,t}(H)=dimI^{s,t}(H).$$
Let $`V=_{k0}V^k`$ be a graded vector space, with $`dimV^k<\mathrm{}`$ for all $`k`$. Assume that each $`V^k`$ is endowed with a MHS. We will refer to such a space as a graded vector space with MHS. The virtual Hodge numbers and the virtual Hodge polynomial of $`V`$ are defined by
$`e^{s,t}(V)={\displaystyle \underset{k0}{}}(1)^kh^{s,t}(V^k),`$
$`e_{x,y}(V)={\displaystyle \underset{s,t}{}}e^{s,t}(V)x^sy^t.`$
Alternatively, we have the splitting:
$$V=\underset{k0}{}\underset{s,t}{}I^{s,t}(V^k).$$
Consider the generating function
$$e_{x,y,z}(V)=\underset{k0}{}\underset{s,t}{}dimI^{s,t}(V^k)x^sy^tz^k.$$
Then $`e_{x,y}(V)=e_{x,y,1}(V)`$. We will use the following convention: for a graded vector space with MHS $`V=_{k0}V^k`$ and a positive integer $`n`$, $`V\{n\}`$ is the graded vector space with MHS such that for each $`k`$,
$`W_m(V^k\{n\})`$ $`=W_{m2n}(V^k\{n\}),`$ $`F^p(V^k\{n\})`$ $`=F^{pn}(V^k\{n\}).`$
It is straightforward to see that $`e^{s,t}(V\{n\})=e^{sn,tn}(V)`$, and so
$$e_{x,y}(V\{n\})=(xy)^ne_{x,y}(V).$$
Deligne has shown that for an arbitrary complex algebraic variety $`Y`$, the cohomology $`H^k(Y)`$ carries a MHS which coincides with the classical pure Hodge structure in the case of smooth projective varieties. Hence one can define the virtual Hodge number of $`Y`$
$$e^{s,t}(Y)=e^{s,t}(H^{}(Y))$$
and the virtual Hodge polynomial of $`Y`$
$$e(Y;x,y)=e_{x,y}(H^{}(Y)).$$
Assume that $`(Y,G)`$ is a pair consisting of a smooth quasi-projective variety $`Y`$ and a finite subgroup $`G`$ of automorphisms of $`Y`$. Then by functorial property, there is an induced action of $`G`$ on the MHS on $`H^{}(Y)`$ by automorphisms. By taking the invariant parts, we obtain a MHS on each $`H^k(Y/G)`$. One can also achieve this by taking a smooth compactification $`\overline{Y}`$ such that $`D=\overline{Y}Y`$ is a divisor with normal crossing singularities and such that the $`G`$-action extends to $`\overline{Y}`$. Then the MHS on $`H^{}(Y/G)`$ is obtained by using $`(\mathrm{\Omega }_Y^{}D)^G`$, the invariant part of the complex of differential forms with logarithmic poles. Using the above MHS on $`H^{}(Y/G)`$, we can now define $`e^{p,q}(Y/G)`$. Similar to the closed case (cf. (3)), we define the orbifold virtual Hodge number as follows:
$`e^{s,t}(Y,G)={\displaystyle \underset{cG_{}}{}}{\displaystyle \underset{\alpha _c=1}{\overset{N_c}{}}}e^{sF_{\alpha _c}^c,tF_{\alpha _c}^c}(Y_{\alpha _c}^c/Z(c)).`$
We also define the orbifold virtual Hodge polynomial:
$`e(Y,G;x,y)={\displaystyle \underset{s,t}{}}e^{s,t}(Y,G)x^sy^t.`$
It is clear that $`e(Y,G;x,y)`$ is the virtual Hodge polynomial of
$$H^{}(Y,G)=\underset{cG_{}}{}\underset{\alpha _c=1}{\overset{N}{}}H^{}(Y_{\alpha _c}^c/Z(c))\{F_{\alpha _c}^c\}$$
(cf (5)), where both sides are understood as graded vector spaces with MHS.
###### Remark 2.2.
One can replace $`H^{}(Y)`$ by the cohomology with compact support $`H_c^{}(Y)`$ in the above definitions.
## 3. The orbifold Hodge numbers of wreath product orbifolds
In this section, we calculate explicitly the ordinary and orbifold Hodge numbers of wreath product orbifolds $`Y^n/G_n`$ associated to an even-dimensional orbifold $`Y/G`$.
### 3.1. Two simple lemmas
Let $`V=_{s,t_+}V^{s,t}`$ be a bigraded complex vector space, such that $`dimV^{s,t}<\mathrm{}`$ for all $`s,t`$, where $`_+`$ is the set of non-negative integers. We introduce the generating function
$$h_{x,y}(V)=\underset{s,t_+}{}(dimV^{s,t})x^sy^t.$$
For example, when $`V`$ is the total Dolbeault cohomology group $`H(Y)`$, then $`h_{x,y}(V)`$ is its associated Hodge polynomial $`h(Y;x,y)`$. When $`V`$ is the total orbifold Dolbeault cohomology group $`H(Y,G)`$, then $`h_{x,y}(V)`$ is its associated orbifold Hodge polynomial $`h(Y,G;x,y)`$. It is actually more convenient to work with $`e_{x,y}(V)=h_{x,y}(V)`$. It is easy to see that
$`e_{x,y}(V_1V_2)=e_{x,y}(V_1)+e_{x,y}(V_2),`$
$`e_{x,y}(V_1V_2)=e_{x,y}(V_1)e_{x,y}(V_2).`$
The graded symmetric algebra of $`V`$ is by definition
$$S(V)=T(V)/I$$
where $`T(V)`$ is the tensor algebra of $`V`$, $`I`$ is the ideal generated by elements of the form
$$vw(1)^{(s+t)(p+q)}wv,vV^{s,t},wV^{p,q}.$$
The bigrading on $`V`$ induces a bigrading on $`T(V)`$ and also on $`S(V)`$, and hence $`e_{x,y}(S(V))`$ makes sense. Note that for bigraded vector spaces $`V_1`$ and $`V_2`$, we have $`S(V_1V_2)S(V_1)S(V_2)`$. Consequently,
(7) $`e_{x,y}(S(V_1V_2))=e_{x,y}(S(V_1))e_{x,y}(S(V_2)).`$
By introducing a formal variable $`q`$ to count the degree of symmetric power, we can write formally $`S(qV)=_{n0}S^n(V)q^n`$. By breaking $`V`$ into one-dimensional subspaces, one can easily prove the following.
###### Lemma 3.1.
For any bigraded vector space $`V=_{s,t0}V^{s,t}`$ with $`dimV^{s,t}<\mathrm{}`$ for all pairs $`(s,t)`$, we have
$$\underset{n0}{}e_{x,y}(S^n(V))q^n=\underset{s,t}{}\frac{1}{(1x^sy^tq)^{e^{s,t}(V)}},$$
where $`e^{s,t}(V)=(1)^{s+t}dimV^{s,t}`$.
For a formal power series $`_{r>0}V_rq^r`$, where each $`V_r`$ is a bigraded vector space of weight $`r`$ such that $`dimV_r^{s,t}<\mathrm{}`$, define
$$S(\underset{r>0}{}V_rq^r)=\underset{m0}{}\underset{_{j=1}^mjm_j=m}{}q^m\underset{j=1}{\overset{m}{}}S^{m_j}(V_j).$$
Formally we have
$$S(\underset{r>0}{}V_rq^r)=\underset{r>0}{}S(V_rq^r)$$
and
$$e_{x,y}(\underset{r>0}{}V_rq^r)=\underset{r>0}{}e_{x,y}(V_r)q^r.$$
Then the next lemma follows from Lemma 3.1.
###### Lemma 3.2.
For a sequence $`\{V_n\}`$ of bigraded vector spaces, we have the following formula:
$$e_{x,y}\left(S(\underset{n>0}{}V_nq^n)\right)=\underset{n>0}{}\underset{s,t}{}\frac{1}{(1x^sy^tq^n)^{e^{s,t}(V_n)}}.$$
###### Remark 3.1.
Using the canonical splitting, it is fairly straightforward to generalize Lemma 3.1 and Lemma 3.2 to the case of vector spaces with MHS.
### 3.2. The main theorems
Since $`G^n`$ is a normal subgroup of the wreath product $`G_n=G^nS_n`$, it is easy to see by (4) that
$$H(Y^n/G_n;x,y)H(Y^n;x,y)^{G^nS_n}S^n(H(Y;x,y)^G)S^n(H(Y/G;x,y)).$$
When $`Y`$ is a compact complex manifold, this is an isomorphism of bigraded vector spaces; when $`Y`$ is a quasi-projective smooth variety over $``$, this is an isomorphism of graded vector spaces with MHS. As a consequence of Lemma 3.1 and Remark 3.1, we obtain the following proposition.
###### Proposition 3.1.
If $`Y`$ is a compact complex manifold or a quasi-projective smooth variety, and $`G`$ is a finite subgroup of automorphisms on $`Y`$, then we have the following formula:
$`{\displaystyle \underset{n0}{}}e(Y^n/G_n;x,y)q^n={\displaystyle \underset{s,t}{}}{\displaystyle \frac{1}{(1x^sy^tq)^{e^{s,t}(Y/G)}}}.`$
The first main result of this paper is the following theorem.
###### Theorem 3.1.
Given a compact complex manifold or a smooth quasi-projective variety $`Y`$ of even complex dimension $`d`$, acted on by a finite group $`G`$ with integer shifts, we have the following formula for the orbifold Hodge numbers:
(8) $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e(Y^n,G_n;x,y)q^n={\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{s,t}{}}{\displaystyle \frac{1}{(1x^sy^tq^r(xy)^{(r1)d/2})^{e^{s,t}(Y,G)}}}.`$
###### Proof.
We first compute the shifts $`F^c`$ for the orbifold $`Y^n/G_n`$ associated to a conjugacy class $`c`$ in $`G_n`$. Consider the typical class containing
$$g\sigma _n=((g,1,\mathrm{},1),(12\mathrm{}n))$$
where $`gG`$. Recall from the previous section that a fixed point in $`Y^n`$ by the action of $`g\sigma _n`$ is of the form $`(x,\mathrm{},x)`$ where $`xY^g`$. Since the calculation can be done locally, we will assume that we take local coordinates $`(z_1,\mathrm{},z_d)`$ near a point $`xY^g`$ such that the action is given by
$$g(z_1,\mathrm{},z_d)=(e^{2\pi \sqrt{1}\theta _1}z_1,\mathrm{},e^{2\pi \sqrt{1}\theta _r}z_r,z_{r+1},\mathrm{},z_d).$$
Equivalently, $`g`$ is locally given by the diagonal matrix $`\mathrm{diag}(e^{2\pi \sqrt{1}\theta _1},\mathrm{},e^{2\pi \sqrt{1}\theta _d})`$ where $`\theta _{r+1}=\mathrm{}=\theta _d=0`$. Then on $`Y^n`$ near $`(x,\mathrm{},x)`$, $`g\sigma _n`$ is given by a block diagonal matrix with blocks of the form
$$\left(\begin{array}{ccccc}0& 1& 0& \mathrm{}& 0\\ 0& 0& 1& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& 0\\ 0& 0& 0& \mathrm{}& 1\\ e^{2\pi \sqrt{1}\theta _j}& 0& \mathrm{}& \mathrm{}& 0\end{array}\right)$$
The characteristic polynomial of this matrix is $`\lambda ^ne^{\sqrt{1}\theta _j}`$, hence it has eigenvalues
$$\lambda _{jk}=e^{2\pi \sqrt{1}(\theta _j+k)/n},k=0,\mathrm{},n1.$$
Notice that $`\lambda _{jk}=1`$ if and only $`\theta _j=k=0`$. So the shift for the component of $`(Y^n)^{g\sigma _n}`$ containing $`(x,\mathrm{},x)`$ is given by
$`F^{g\sigma _n}(x,\mathrm{},x)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{r}{}}}{\displaystyle \underset{k=0}{\overset{n1}{}}}{\displaystyle \frac{\theta _j+k}{n}}+(dr){\displaystyle \underset{k=1}{\overset{n1}{}}}{\displaystyle \frac{k}{n}}`$
$`=`$ $`{\displaystyle \underset{j=1}{\overset{r}{}}}\theta _j+(n1)d/2=F_{\alpha _c}^c+(n1)d/2.`$
Here we have assumed that $`xY^g`$ lies in the component $`Y_{\alpha _c}^g`$ $`(\alpha _c=1,\mathrm{},N_c)`$, and $`F_{\alpha _c}^c`$ is the shift for the component $`Y_{\alpha _c}^c/Z_G(c)`$.
For a general conjugacy class containing an element $`a`$ of type
$$\rho =\{m_r(c)\}_{r1,cG_{}},$$
where $`_{r,c}rm_r(c)=n`$, the description of the fixed-point set $`(Y^n)^a`$ given in (2) implies that the components for $`(Y^n)^a`$ can be listed as
$$(Y^n)_{\{m_{r,c}(\alpha _c)\}}^a=\underset{r,c}{}\underset{\alpha _c=1}{\overset{N_c}{}}S^{m_{r,c}(\alpha _c)}(Y_{\alpha _c}^c/Z_G(c)),$$
where $`(m_{r,c}(1),\mathrm{},m_{r,c}(N_c))`$ satisfies $`_{\alpha _c=1}^{N_c}m_{r,c}(\alpha _c)=m_r(c)`$. Then the shift for the component $`(Y^n)_{\{m_{r,c}(\alpha _c)\}}^a`$ is given by
(9) $`F_{\{m_{r,c}(\alpha _c)\}}={\displaystyle \underset{r,c}{}}{\displaystyle \underset{\alpha _c=1}{\overset{N_c}{}}}m_{r,c}(\alpha _c)\left(F_{\alpha _c}^c+(r1)d/2\right).`$
By using (2.2), (2), (9) and (7) we have
$`{\displaystyle \underset{n0}{}}H(Y^n,G_n;x,y)q^n`$
$`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\{m_r(c)\}𝒫_n(G_{})}{}}{\displaystyle \underset{r,c}{}}{\displaystyle \underset{\alpha _c=1}{\overset{N_c}{}}}H(S^{m_{r,c}(\alpha _c)}(Y_{\alpha _c}^c/Z_G(c));x,y)\{F_{\{m_{r,c}(\alpha _c)\}}\}q^n`$
$`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\{m_r(c)\}𝒫_n(G_{})}{}}`$
$`{\displaystyle \underset{r,c}{}}{\displaystyle \underset{\alpha _c=1}{\overset{N_c}{}}}S^{m_{r,c}(\alpha _c)}\left(H(Y_{\alpha _c}^c/Z_G(c);x,y)\{F_{\alpha _c}^c+(r1)d/2\}\right)q^n`$
$`=`$ $`{\displaystyle \underset{\{m_r(c)\}}{}}{\displaystyle \underset{r,c}{}}S^{m_{r,c}}\left({\displaystyle \underset{\alpha _c=1}{\overset{N_c}{}}}H(Y_{\alpha _c}^c/Z_G(c);x,y)\{F_{\alpha _c}^c+(r1)d/2\}q^r\right)`$
$`=`$ $`{\displaystyle \underset{\{m_r\}}{}}{\displaystyle \underset{r}{}}S^{m_r}\left({\displaystyle \underset{c}{}}{\displaystyle \underset{\alpha _c=1}{\overset{N_c}{}}}H(Y_{\alpha _c}^c/Z_G(c);x,y)\{F_{\alpha _c}^c+(r1)d/2\}q^r\right)`$
$`\text{where we let }m_r={\displaystyle \underset{c}{}}m_r(c)`$
$`=`$ $`{\displaystyle \underset{\{m_r\}}{}}{\displaystyle \underset{r1}{}}S^{m_r}\left(H(Y,G;x,y)\{(r1)d/2\}q^r\right)`$
$`=`$ $`S\left({\displaystyle \underset{r1}{}}H(Y,G;x,y)\{(r1)d/2\}q^r\right).`$
Namely we have proved that
$`{\displaystyle \underset{n0}{}}H(Y^n,G_n;x,y)q^n=S\left({\displaystyle \underset{r>0}{}}H(Y,G;x,y)(xy)^{(r1)d/2}q^r\right),`$
which implies immediately the theorem by means of Lemma 3.2. ∎
###### Remark 3.2.
When $`G`$ is trivial and $`Y`$ is an algebraic surface, Theorem 3.1 recovers the orbifold Hodge numbers for the symmetric product $`Y^n/S_n`$ which was calculated in . On the other hand, if we set $`x=y=1`$ we recover the orbifold Euler numbers for $`Y^n/G_n`$ which was first computed in for any topological space $`Y`$.
###### Remark 3.3.
In the above we have constrained ourselves to the case that the shift numbers are integers. Physicists are also interested in the case of fractional shift numbers (see e.g. Zaslow ). It is straightforward to generalize our result.
### 3.3. Some consequences
We assume that $`Y`$ is a quasi-projective surface acted upon by a finite group $`G`$, and that $`X`$ is a resolution of singularities of the orbifold $`Y/G`$. We denote by $`X^{[n]}`$ the Hilbert scheme of $`n`$ points on $`X`$. It is well known (cf. ) that the Hilbert-Chow morphism $`X^{[n]}X^n/S_n`$ is a resolution of singularities. Indeed it is crepant. We have the following commutative diagram
$$\begin{array}{ccc}X^{[n]}& & X^n/S_n\\ & & & & \\ Y^n/G_n& \stackrel{}{}& (Y/G)^n/S_n\end{array}$$
which implies that the Hilbert scheme $`X^{[n]}`$ is a resolution of singularities of the orbifold $`X^n/G_n`$.
As calculated in and , the Hodge numbers for the Hilbert scheme $`X^{[n]}`$ are given by the following formula:
$`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e(X^{[n]};x,y)q^n={\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{s,t}{}}{\displaystyle \frac{1}{(1x^sy^tq^r(xy)^{r1})^{e^{s,t}(X)}}}.`$
Here and below we use the cohomology with compact supports.
By comparing with Theorem 3.1 we obtain the following theorem which provides us a large class of higher dimensional examples which verify the orbifold Hodge number conjecture (cf. ). The assumption of the theorem is necessary by Remark 2.1.
###### Theorem 3.2.
Let $`Y`$ be a smooth quasi-projective surface which admits a $`G`$-action with only isolated fixed points. Assume that $`\pi :XY/G`$ is a resolution such that $`e(X;x,y)=e(Y,G;x,y)`$. Let $`X^{[n]}`$ be the Hilbert scheme of $`n`$ points of $`X`$. Then for all $`r,s`$ we have
$$h^{r,s}(X^{[n]})=h^{r,s}(Y^n,G_n).$$
###### Remark 3.4.
When $`G`$ is trivial and $`X`$ equals $`Y`$, we recover the theorem of . We will see later many interesting examples arise when $`G`$ is not trivial.
More generally if $`Y`$ has dimension greater than $`2`$, there is no such a favorable resolution as Hilbert scheme for $`Y^n/G_n`$. Nevertheless we have the following interesting corollary of Theorem 3.1. Here we assume that the shifts are integers for the orbifold $`Y/G`$ so that its orbifold Hodge numbers are well defined.
###### Corollary 3.1.
Let $`Y`$ be a smooth variety of even dimension acted on by a finite group $`G`$ of automorphisms, and $`\pi :XY/G`$ is a resolution such that $`h^{s,t}(X)=h^{s,t}(Y,G)`$ for all $`s,t`$, then for all $`s,t`$ we have
$$h^{s,t}(X^n,S_n)=h^{s,t}(Y^n,G_n).$$
## 4. Examples and applications
In this section we provide various concrete examples which satisfy the assumptions of Theorem 3.1 and Theorem 3.2. We also give explicit conjectures on the elliptic genera for the wreath product orbifolds.
### 4.1. Various examples
###### Example 4.1.
When $`G`$ is trivial and $`X`$ equals $`Y`$, this gives us the example of symmetric products .
###### Example 4.2.
$`Y`$ is $`^2`$, $`G`$ is a finite subgroup of $`SL_2()`$, and $`X`$ is the minimal resolution of $`^2/G`$. The exceptional fiber consists of $`|G_{}|1`$ irreducible components which are $`(2)`$-curves (cf. e.g. ). We have
$`h^{s,t}(X)=\{\begin{array}{cc}1,\hfill & s=t=0,\hfill \\ |G_{}|1,\hfill & s=t=1,\hfill \\ 0,\hfill & \text{otherwise}.\hfill \end{array}`$
On the other hand, for any non-trivial conjugacy class $`cG_{}`$, the corresponding shift is $`1`$ and thus makes a contribution to $`h^{1,1}(^2,G)`$ which results that $`h^{1,1}(^2,G)=|G_{}|1.`$ The other $`h^{s,t}(^2,G)`$ can be also seen to coincide with $`h^{s,t}(X)`$.
This example has played a key role in the connections between the wreath product orbifolds and the vertex representations of affine and toroidal Lie algebras .
###### Example 4.3.
(Bryan-Donagi-Leung ) Let $`Y`$ be an abelian surface (two dimensional torus). The $`_2`$-action induced by the involution $`\tau :xx`$ has $`16`$ fixed points, at each of which the shift $`F^\tau `$ is $`1`$. So the twisted sectors contribute an extra $`16`$ to $`h^{1,1}`$. Write $`Y=^2/L`$ for some lattice $`L`$, and let $`(z_1,z_2)`$ be the linear coordinates on $`^2`$. Then $`H^,(Y)`$ is generated by $`dz^1,d\overline{z}^1,dz^2,d\overline{z}^2`$. The action of $`\tau `$ just takes $`dz^j`$ to $`dz^j`$, etc. Hence it is clear that
$`H^,(Y)^_2`$ $``$ $`dz^1dz^2(_{j,k=1}^2dz^jd\overline{z}^k)`$
$`d\overline{z}^1d\overline{z}^2dz^1dz^2d\overline{z}^1d\overline{z}^2.`$
Therefore,
$`h^{s,t}(Y,_2)=\{\begin{array}{cc}1,\hfill & s=t=0,\hfill \\ 20,\hfill & s=t=1,\hfill \\ 1,\hfill & s=2,t=0\text{or}s=0,t=2,\hfill \\ 1,\hfill & s=t=2,\hfill \\ 0,\hfill & \text{otherwise}.\hfill \end{array}`$
The minimal resolution $`XY/\pm 1`$ is a crepant resolution, where $`X`$ is a K3 surface. This is the famous Kummer construction. By the well known Hodge numbers of a K3 surface, we have $`h^{s,t}(X)=h^{s,t}(Y,_2)`$ for all $`s,t`$.
###### Example 4.4.
Let $`_3`$ act on $`_2`$ by
$$\alpha [z_0:z_1:z_2]=[\alpha z_0:\alpha ^1z_1:z_2],$$
where $`\alpha `$ is a generator of $`_3`$ and identified with a cubic root of unity on the right-hand side. This action has three fixed points: $`p_0=[1:0:0]`$, $`p_1=[0:1:0]`$, and $`p_2=[0:0:1]`$. At these point, the weights of the action are $`(1,2)`$, $`(2,1)`$, and $`(1,2)`$ respectively. It is then straightforward to see that for $`g1`$ we have
$$F^g=\frac{1}{3}+\frac{2}{3}=1.$$
Therefore,
$`H^,(_2,_3)=H^,(_2)^_3{\displaystyle \underset{j=0}{\overset{2}{}}H^,(p_j)^_3\{1\}\underset{j=0}{\overset{2}{}}H^,(p_j)^_3\{1\}},`$
and hence
$`h^{s,t}(_2,_3)=\{\begin{array}{cc}1,\hfill & s=t=0,\hfill \\ 7,\hfill & s=t=1,\hfill \\ 1,\hfill & s=t=2,\hfill \\ 0,\hfill & \text{otherwise}.\hfill \end{array}`$
The minimal resolution $`X=\widehat{_2/_3}`$ is obtained by replacing each singular point by a string of two $`(2)`$-curves, each of which contributes $`1`$ to $`h^{1,1}`$, hence $`h^{1,1}`$ of $`\widehat{_2/_3}`$ is $`7`$. This resolution is a crepant resolution.
###### Example 4.5.
Let $`n>2`$ be an odd number, Consider the action of $`_n`$ on $`_3`$ given by
$$\alpha [z_0:z_1:z_2:z_3]=[z_0:z_1:\alpha z_2:\alpha ^1z_3],$$
where $`\alpha `$ is a generator of $`_n`$. It has a fixed line $`\{[z_0:z_1:0:0]\}`$ and two isolated fixed points $`[0:0:1:0]`$, and $`[0:0:0:1]`$. Let $`Y_{m,n}`$ be the Fermat surface defined by
$$z_0^{mn}+z_1^{mn}+z_2^{mn}+z_3^{mn}=0$$
in $`_3`$. The above action preserves $`Y_{m,n}`$, with $`mn`$ isolated fixed points:
$$[1:e^{(2k+1)\pi \sqrt{1}/(mn)}:0:0],k=0,\mathrm{},mn1.$$
Note the action is semi-free, i.e. the stabilizers are either trivial or the whole group $`_n`$. Near each of the fixed points, say $`[1:e^{\pi \sqrt{1}/(mn)}:0:0]`$, $`Y_{m,n}`$ is given by the equation
$$1+u_1^n+u_2^n+u_3^n=0,$$
where $`u_j=z_j/z_0`$. We can use $`(u_2,u_3)`$ as local coordinates, then $`_n`$ acts with weight $`(1,1)`$, i.e., $`_n`$ acts locally by matrices in $`SL(2,)`$. Therefore, $`Y_{m,n}/Z_n`$ admits a crepant resolution obtained by replacing each isolated singular point with a string of $`n1`$ copies of $`(2)`$-curves.
###### Example 4.6.
Denote now by $`\beta `$ a generator of $`_4`$. Consider the $`_4`$-action on $`_3`$ given by
$$\beta [z_0:z_1:z_2:z_3]=[z_0:z_1:\sqrt{1}z_3:\sqrt{1}z_2].$$
Combined with the $`_n`$-action in Example 4.5, we get an action of the binary dihedral group $`D_n^{}`$ on $`_3`$ which preserves $`Y_{4m,n}`$. By the same method as in Example 4.5 one can find the fixed points and sees that $`Y_{4m,n}/D_n^{}`$ admits a crepant resolution.
###### Example 4.7.
The method of Example 4.5 and Example 4.6 can be generalized to other finite subgroups of $`SL(2,)`$. Given such a group $`G`$, let it act on $`^4`$ on the last two factors. This action induces an action on $`_3`$. Now consider a smooth hypersurface $`Y`$ defined by an equation of the form
$$f(z_0,z_1)+g(z_2,z_3)=0,$$
where $`f`$ and $`g`$ are two homogeneous polynomials of the same degrees, and $`g`$ is an invariant polynomial for $`G`$. Using the explicit description of the $`G`$-action on $`^2`$ and the invariant polynomials (see e.g. Klein ), one can find many examples which admits crepant resolutions. One should be able to find more examples by considering complete intersections in (weighted) projective spaces.
###### Example 4.8.
More complicated examples can be found in two papers by Barlow , e.g. the quotient of a Hilbert modular surface by $`_2`$ or $`D_{10}`$, or the quotient of a complete intersection of $`4`$ quadrics in $`_6`$ by a group of order $`16`$, or the quotient of a Godeaux-Reid surface by an involution.
### 4.2. Conjectures on elliptic genera of wreath product orbifolds
Let $`Y`$ be a compact Kähler manifold of complex dimension $`d`$, denote by $`TY`$ (resp. $`T^{}Y`$) its holomorphic tangent (resp. cotangent) bundle. Consider the formal power series of vector bundles:
$$E_{q,y}(Y)=y^{\frac{d}{2}}\underset{n1}{}\left(\mathrm{\Lambda }_{yq^{n1}}(T^{}Y)\mathrm{\Lambda }_{y^1q^n}(TY)S_{q^n}(T^{}Y)S_{q^n}(TY)\right).$$
If we write
$$E_{q,y}(Y)=\underset{m0,l}{}q^my^lE_{m,l}(Y),$$
we easily see that each $`E_{m,l}`$ is a holomorphic bundle of finite rank, hence one can consider its Riemann-Roch number
$$c(m,l)=\chi (E_{m,l}(Y))=\underset{k0}{}(1)^kdimH^k(Y,E_{m,l}(Y)).$$
The generating function
$$\chi (Y;q,y)=\underset{m0,l}{}q^my^l\chi (E_{m,l}(Y))=\chi (E_{q,y}(Y))$$
is called the elliptic genus of $`Y`$ (cf. ). In the very important special case when $`q=0`$, one recovers the Hirzebruch genus:
$`E_{0,y}(Y)=y^{\frac{d}{2}}\mathrm{\Lambda }_y(T^{}Y),`$
$`\chi (Y;0,y)=y^{\frac{d}{2}}\chi _y(Y)=y^{\frac{d}{2}}{\displaystyle \underset{s,t0}{}}(1)^t(y)^sh^{s,t}(Y).`$
We do not know of a good mathematical formulation of elliptic genera for orbifolds. However physicists have interpreted elliptic genera as partition functions of supersymmetric sigma models, which makes sense also for orbifolds (cf. and references therein). Based on physical arguments and the description of fixed-point sets for the symmetric group action on $`Y^n`$, Dijkgraaf et al derived a formula for the elliptic genera of the symmetric products $`S^n(Y)`$ in terms of that of $`Y`$. In the case of a K3 surface or an abelian surface, they also conjectured that the same formula should compute the elliptic genera of the Hilbert schemes. Their method, if can be made mathematically rigorous, should also provide the proof of the following conjectures with suitable modifications. (In the following we denote by $`\chi (Y,G;q,y)`$ the elliptic genera of an orbifold $`Y/G`$. )
###### Conjecture 4.1.
Let $`Y`$ be a Kähler $`G`$-manifold. If we write the elliptic genus for $`Y/G`$ as $`\chi (Y,G;q,y)=_{m0,l}c(m,l)q^my^l`$, then the elliptic genus for the wreath product orbifold $`Y^n/G_n`$ is given by the following formula:
$$\underset{N=0}{\overset{\mathrm{}}{}}p^N\chi (Y^N,G_N;q,y)=\underset{n>0,m0,l}{}\frac{1}{(1p^nq^my^l)^{c(nm,l)}}.$$
###### Conjecture 4.2.
Let $`Y`$ be a Kähler $`G`$-surface. We assume that $`X`$ is a resolution of singularities of $`Y/G`$ such that $`\chi (Y,G;q,y)=\chi (X;q,y).`$ Then $`\chi (Y^n,G_n;q,y)=\chi (X^{[n]};q,y)`$ for all $`n`$.
When $`G`$ is trivial, one recovers the symmetric product situation as in . In this case, the $`q=0`$ version of Conjecture 4.1 has been verified in as a corollary of the calculation of orbifold Hodge numbers. Similarly, our results in Section 3 can be viewed as supporting evidence for the above conjectures in the general setup of wreath product orbifolds.
Note added. In a recent remarkable paper , Borisov and Libgober have introduced the mathematically rigorous notion of orbifold elliptic genera among other things, and verified our Conjecture 4.1.
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# 1 Introduction
## 1 Introduction
Over the past two years there has been a great deal of progress in the study of non-BPS solitonic states in string theory . Among these are the stable non-BPS D-branes of certain orbifold theories , and of the Type I (and Type IA) theories . All of these stable non-BPS states can be elegantly characterized in terms of K theory .
Most of the above non-BPS D-branes can be obtained by suitable projections from the corresponding unstable non-BPS states of the Type II theories . These unstable non-BPS D$`p`$-branes occur for the complementary values of the (stable) BPS branes, i.e. they have odd $`p`$ in Type IIA and even $`p`$ in Type IIB. However, in the case of Type I there exists also a $`𝖹𝖹_2`$ D-instanton (D$`(1)`$-brane), which has the same value of $`p`$ as the corresponding Type IIB D$`(1)`$-brane. (The same is also the case for the Type I D$`7`$-brane.) The presence of the Type I D-instanton is associated with the fact that the actual gauge group of Type I is $`SO(32)`$ rather than the gauge group of the perturbative Type I theory which is $`O(32)`$.
In this paper we will determine some explicit effects of the Type I D-instanton by studying a simple scattering process in M-theory that can be identified, via familiar dualities, with a Type I string theory process. We will start, in section 2, by considering dynamics in the eleven-dimensional geometry considered by Hořava and Witten, in which the eleven-dimensional M-theory space-time is taken to be $`M^{10}\times R`$, where $`M^{10}`$ is ten-dimensional Minkowski space and $`R`$ is the interval $`0x^{11}d`$ . The tree-level contribution to the scattering of two $`E_8`$ gauge particles that are confined to distinct boundaries of the eleven-dimensional space-time will be evaluated in section 3, as illustrated in figure 1. Since the external particles are localized on distinct boundaries the exchanged bulk states are singlets under both the $`E_8`$ gauge groups. We will verify that in the low-momentum limit this amplitude reproduces the same expression as that obtained in the $`E_{8L}\times E_{8R}`$ heterotic string (which will be referred to as the $`HE`$ theory) in the same limit, where the subscripts $`L`$ and $`R`$ refer to the two boundaries (left and right) at $`x^{11}=0`$ and $`x^{11}=d`$. This is true independent of the separation of the boundaries.
In section 4 we will consider the theory compactified on a circle, $`S^1`$, in the $`x^9`$ direction with the symmetry broken to $`SO(16)_L\times SO(16)_R`$. We will be particularly interested in the scattering of massive spinor states localized on the two nine-dimensional boundaries that are modes of the massless ten-dimensional states with Kaluza–Klein charge $`\pm 1/2`$. The four-point function for these states follows very simply from the ten-dimensional expression. We will also describe the interpretation of this amplitude in terms of the Type IA theory in which the background is that of the Type IIA theory on the $`𝖹𝖹_2`$ orbifold of a circle. In this description the scattering spinor states are D-particles ($`D_S`$) or anti D-particles ($`\overline{D}_S`$) stuck to the two orientifold planes. The scattering amplitudes that will be considered are ‘tree-level’ interactions. Processes of the form $`D_S+D_SD_S+D_S`$, $`D_S+\overline{D}_SD_S+\overline{D}_S`$ are described by the exchange of closed strings (states of zero Kaluza–Klein charge) between the separated orientifold planes <sup>1</sup><sup>1</sup>1Here the notation $`L+RL^{}+R^{}`$ denotes a scattering process in which $`L`$ and $`R`$ are incoming states on the left and right orientifolds, respectively, while the corresponding outgoing states are $`L^{}`$ and $`R^{}`$.. At low energies the former process is BPS, and the amplitude vanishes in the zero velocity limit. In the impact parameter description, the resulting force is proportional to $`v^4`$ which is in agreement with the string theory calculation that can be performed following . The second process is non-BPS, and the amplitude is singular as $`v0`$, reflecting the instability of the system to decay into a non-BPS D1-brane.
The process $`D_S+\overline{D}_S\overline{D}_S+D_S`$ is described by the exchange of a $`D0`$-brane (the state of Kaluza–Klein charge $`1`$) between the separated orientifold planes. The effect of a massive particle exchange generates an exponentially small amplitude. After a T-duality in the $`x^{11}`$ direction this is reinterpreted as the effect of the Type I D-instanton. This is analogous to the description of the Type IIB D-instanton as the dual of the world-line of a Type IIA D-particle which also has an origin in eleven-dimensional supergravity . As in that case, the supergravity Feynman diagram calculation automatically takes into account the appropriate quantum measure, including fermionic zero modes. These modes are interpreted in terms of massless string states in the Type IA picture. As expected, the D-instanton contribution respects $`SO(16)_L\times SO(16)_R`$ but is not invariant under $`O(16)_L\times SO(16)_R`$ or $`SO(16)_L\times O(16)_R`$, as will be described in some detail in section 4.2. Finally, the discussion in section 5 contains some speculations about the situation in which the four particles in the scattering amplitude are located on one of the two boundaries in the nine-dimensional theory.
## 2 Feynman diagrams in the Hořava-Witten geometry
We are interested in calculating the tree-level scattering amplitude illustrated in figure 1, in which the propagator joins the two boundaries which are separated by a distance $`d`$. The scattering particles with momenta $`k_r`$ ($`r=1,\mathrm{},4`$) are the states confined to the ten-dimensional boundaries which are massless $`E_8`$ gauge particles in the unbroken theory. Later we will consider compactification on a circle, in which case the gauge group may be broken and the scattering particles may carry non-zero Kaluza–Klein charge and be massive states. When the external particles are bosons (as will be the case in this paper) the propagating intermediate state in the figure may be either a graviton, $`h_{\mu \nu }`$, or a third-rank antisymmetric tensor potential, $`C_{\mu \nu \rho }^{(3)}`$, and the full amplitude is obtained by summing over both these contributions (where the eleven-dimensional indices span the range $`\mu ,\nu ,\rho =0,1,\mathrm{},9,11`$).<sup>2</sup><sup>2</sup>2Recently a similar calculation was undertaken by Krause but the three-form exchange was omitted so his result is therefore not complete.
The Hořava–Witten geometry may be thought of as a $`𝖹𝖹_2`$ orbifold, where the generator of the orbifold group acts on the eleventh dimension, $`x^{11}`$, as $`x^{11}x^{11}`$, as well as acting on the three-form field $`C`$ as $`CC`$. In the resulting theory the momentum in the eleventh dimension, $`p_{11}`$, is not conserved, whilst it is still conserved in the tangential ten dimensions. Furthermore, the boundary conditions on the bulk fields require that $`h_{M\mathrm{\hspace{0.17em}11}}=0=C_{MNP}`$ on the boundaries $`x^{11}=0`$ and $`x^{11}=d`$, where $`M,N,P=0,1,\mathrm{},9`$ are the ten-dimensional indices.<sup>3</sup><sup>3</sup>3We are here following the ‘upstairs’ formalism of , in which the fields span the covering space $`R^{1,9}\times S^1`$ and the orbifold conditions are imposed by modding out by the action of $`𝖹𝖹_2`$.
We will consider the simple example of the propagator for a scalar field in this background by first considering M-theory compactified on a circle with background geometry $`M_{10}\times S^1`$, where $`S^1`$ is a circle of radius $`d/\pi =R_{11}l_p`$ (where $`R_{11}`$ is the dimensionless radius and $`l_p`$ is the eleven-dimensional Planck length). The momentum-space scalar propagator is
$$\stackrel{~}{G}(p,p_{11})=\frac{1}{p^2+p_{11}^2},$$
(2.1)
where $`p_{11}=m(R_{11}l_p)^1=\pi md^1`$ and $`p^M`$ ($`M=0,1,\mathrm{},9`$) is the arbitrary momentum in the ten dimensions tangential to the boundaries. The scalar propagator evaluated between two points, $`x^{11}`$ and $`y^{11}`$, on the circle is therefore given by the Fourier sum,
$$G(p^M;x^{11}y^{11})=\frac{1}{2d}\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}e^{i\frac{\pi m}{d}(x^{11}y^{11})}\frac{1}{p^2+\frac{\pi ^2m^2}{d^2}}.$$
(2.2)
We are really interested in defining the propagator on the $`𝖹𝖹_2`$ orbifold of the circle. If we choose the scalar field to be symmetric under the action of the $`𝖹𝖹_2`$ then the propagator must be invariant under $`x^{11}x^{11}`$ as well as $`x^{11}x^{11}+2d`$ (with similar conditions on $`y^{11}`$). The propagator satisfying these conditions is then given by
$$𝒢(p^M;x^{11},y^{11})=G(p^M;x^{11}y^{11})+G(p^M;x^{11}+y^{11}).$$
(2.3)
Since the external particles are constrained to the planes $`x^{11}=0`$ and $`x^{11}=d`$, the propagator that enters in the scattering amplitude is obtained by substituting (2.2) into (2.3) and setting $`x^{11}=0,y^{11}=d`$, which gives,
$`𝒢(p^M;0,d)`$ $`=`$ $`{\displaystyle \frac{1}{d}}{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^m}{p^2+\frac{\pi ^2m^2}{d^2}}}={\displaystyle \frac{1}{d}}{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}(1)^m{\displaystyle _0^{\mathrm{}}}𝑑\sigma e^{\sigma (p^2+\frac{\pi ^2m^2}{d^2})}`$ (2.4)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{\pi }}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}𝑑\sigma \sigma ^{\frac{1}{2}}e^{\sigma p^2\frac{d^2}{\sigma }(n+\frac{1}{2})^2},`$
where the last line follows after performing a Poisson resummation that replaces the integer Kaluza–Klein charge $`m`$ by the integer $`2n+1`$ that labels the winding number of the propagator around the $`x^{11}`$ direction. (In particular, the sum in (2.4) only contains contributions with odd winding number; this reflects the fact that the states in the propagator run from $`x^{11}=0`$ to $`x^{11}=d`$.) The $`\sigma `$ integral is standard (the saddle point approximation is exact) and leads to the result
$$𝒢(p^M;0,d)=\frac{1}{\sqrt{p^2}}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}e^{|2n+1|\sqrt{p^2}d}=\frac{1}{\sqrt{p^2}}\frac{1}{\mathrm{sinh}(\sqrt{p^2}d)}.$$
(2.5)
In deriving (2.5) we have assumed that $`p^2`$ is positive (so that the integral in (2.4) converges). However it is clear from the definition of $`𝒢(p^M;0,d)`$ in (2.4) that $`𝒢`$ is a meromorphic function of $`p^2`$; since the right-hand-side of (2.5) is also meromorphic, the final result will therefore hold in general (i.e. not only for positive $`p^2`$). We will see later how the terms in the sum over $`n`$ correspond to the contributions of D-instantons to the Type I theory in nine dimensions.
In the amplitude calculations that follow, the vertices will involve momentum factors. In the case of the exchange of the three-form potential this will result in a term in which a factor of $`p_{11}^2`$ is inserted into the numerator of the propagator, so that the sum in (2.4) is replaced by
$`\widehat{𝒢}(p^M;0,d)`$ $`=`$ $`{\displaystyle \frac{1}{d}}{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^mp_{11}^2}{p^2+p_{11}^2}}={\displaystyle \frac{1}{d}}{\displaystyle \underset{m𝖹𝖹}{}}(1)^m{\displaystyle _0^{\mathrm{}}}𝑑\sigma e^{\sigma p^2}{\displaystyle \frac{d}{d\sigma }}e^{\frac{\pi ^2m^2}{d^2}\sigma }`$ (2.6)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{\pi }}}{\displaystyle \underset{n𝖹𝖹}{}}{\displaystyle _0^{\mathrm{}}}𝑑\sigma e^{\sigma p^2}{\displaystyle \frac{d}{d\sigma }}\left(\sigma ^{\frac{1}{2}}e^{\frac{d^2}{\sigma }(n+\frac{1}{2})^2}\right)`$
$`=`$ $`\sqrt{p^2}{\displaystyle \frac{1}{\mathrm{sinh}(\sqrt{p^2}d)}},`$
where a Poisson resummation has been performed in the second step, and we have integrated by parts and used (2.5) to obtain the final line.
The propagators that enter into the tree diagrams of interest are those of the metric and the three-form potential. Writing the metric in non-compact eleven-dimensional space-time as $`g_{\mu \nu }=\eta _{\mu \nu }+h_{\mu \nu }`$, where $`\eta _{\mu \nu }`$ is the eleven-dimensional background Minkowski metric, the $`h_{\mu _1\nu _1}h_{\mu _2\nu _2}`$ propagator (in the de Donder gauge) is given by
$$\stackrel{~}{G}_{\mu _1\nu _1;\mu _2\nu _2}(p^M,p_{11})=\kappa ^2\left(\eta _{\mu _1\mu _2}\eta _{\nu _1\nu _2}+\eta _{\mu _1\nu _2}\eta _{\nu _1\mu _2}\frac{2}{9}\eta _{\mu _1\nu _1}\eta _{\mu _2\nu _2}\right)\frac{1}{p^2+p_{11}^2}.$$
(2.7)
Similarly, writing $`C_{\mu \nu \rho }^{(3)}=c_{\mu \nu \rho }^{(3)}+\widehat{C}_{\mu \nu \rho }^{(3)}`$, where $`c_{\mu \nu \rho }^{(3)}=0`$ is the background antisymmetric potential, the $`\widehat{C}_{\mu _1\nu _1\rho _1}^{(3)}\widehat{C}_{\mu _2\nu _2\rho _2}^{(3)}`$ propagator is given (in Feynman gauge) by
$$\stackrel{~}{G}_{\mu _1\nu _1\rho _1;\mu _2\nu _2\rho _2}(p^M,p_{11})=\frac{\kappa ^2}{(3!)^2}\eta _{\mu _1[\mu _2}\eta _{|\nu _1|\nu _2}\eta _{|\rho _1|\rho _2]}\frac{1}{p^2+p_{11}^2}.$$
(2.8)
The vertices that couple these propagators to the gauge particles in the boundaries are determined by the supergravity action in the Hořava–Witten geometry which is given in and reproduced in appendix A. The tree-level amplitudes involve the cubic vertices coupling a pair of gauge particles to a graviton or an antisymmetric potential. The graviton coupling $`AAh`$ is given by
$$S_{YM}^{AAh}=\frac{1}{(4\pi )^{5/3}\kappa ^{4/3}}_{R^{1,9}}d^{10}x\left(_MA_N^a_PA_Q^ah_{RS}\right)\left(\frac{1}{2}\eta ^{M[Q}\eta ^{P]N}\eta ^{RS}+\eta ^{M[P}\eta ^{Q]S}\eta ^{RN}+\eta ^{M[R}\eta ^{Q]N}\eta ^{PS}\right).$$
(2.9)
The three-form coupling $`AA\widehat{C}`$ is extracted from the $`G^2`$ term in the bulk action. This interaction arises because the Bianchi identity for the bulk field strength $`G_{MNP\mathrm{\hspace{0.17em}11}}`$ receives a boundary contribution in a manner that is determined by requiring local supersymmetry as discussed in ,
$$G_{MNP\mathrm{\hspace{0.17em}11}}=4!_{[M}\widehat{C}_{NP\mathrm{\hspace{0.17em}11}]}+\frac{\kappa ^{2/3}}{\sqrt{2}(4\pi )^{5/3}}\left(\delta (x^{11})+\delta (x^{11}d)\right)\omega _{MNP},$$
(2.10)
where $`\omega `$ is the Chern-Simons three-form defined by
$$\omega _{MNP}=2\mathrm{Tr}\left(A_M_{[N}A_{P]}+\frac{1}{3}A_M[A_N,A_P]+\text{cyclic perms.}\right).$$
(2.11)
The other components of $`G`$ are eliminated on the boundary since the three-form field is anti-symmetric under $`x^{11}x^{11}`$. The $`AAC`$ interaction comes from the term linear in $`\omega _{MNP}`$ in $`G_{MNP\mathrm{\hspace{0.17em}11}}^2`$ and is given by
$$S^{AAC}=\frac{\sqrt{2}}{(4\pi )^{5/3}\kappa ^{4/3}}_{R^{1,9}\times S^1}d^{11}x\left(\delta (x^{11})+\delta (x^{11}d)\right)_{[M}\widehat{C}_{NP11]}\mathrm{Tr}(A^M^{[N}A^{P]}).$$
(2.12)
## 3 Scattering of boundary gauge particles
We will now consider the amplitude illustrated in fig. 1 that describes tree-level elastic scattering of an incoming gauge particle localized on one of the boundaries with a gauge particle localized on the other boundary. Since the exchanged bulk fields do not carry gauge quantum numbers the diagram has an overall group theory factor $`\mathrm{Tr}_L(T_1T_2)\mathrm{Tr}_R(T_3T_4)`$, where $`\mathrm{Tr}_L`$ denotes the trace over the 248 components of the adjoint representation of the $`E_{8L}`$ gauge group on one boundary while $`\mathrm{Tr}_R`$ denotes the trace over $`E_{8R}`$ on the other. We will verify that this amplitude reduces in the limit of small momenta to the same expression as that of the tree-level $`E_{8L}\times E_{8R}`$ heterotic string with coupling constant related to the radius $`R_{11}`$ in the usual manner. We will only consider the case in which all the external states are gauge bosons, in which case the exchanged states are the three-form potential and the graviton.
The momentum of the scattering states is restricted to the ten flat directions parallel to the boundaries, and the Mandelstam variables $`s,t,u`$ are defined in terms of the ten-dimensional momenta $`k_1,k_2,k_3,k_4`$ by
$$s=(k_1+k_2)^2,t=(k_1+k_4)^2,u=(k_1+k_3)^2,$$
(3.1)
with $`s+t+u=0`$, which follows from momentum conservation, $`_rk_r=0`$, and the mass-shell conditions, $`k_r^2=0`$, for the external physical states. These vector states have polarization vectors $`ϵ_r^M`$ which satisfy the transversality conditions $`ϵ_rk_r=0`$.
The non-standard aspect of the tree diagrams illustrated in fig. 1 is that the two vertices are constrained to lie in the boundaries at $`x^{11}=0`$ and $`x^{11}=d=\pi R_{11}l_p`$. This means that momentum is not conserved in the eleventh direction and we must use propagators of the form $`𝒢(p;0,d)`$, suitably generalized to account for the tensor indices on the propagating fields as in (2.7) and (2.8).
The expressions for the tree amplitude due to the graviton and the three-form exchange are given in appendix B by $`A_4^{(h)}`$ (B.2) and $`A_4^{(C)}`$ (B.3). The complete tree-level contribution is the sum of these two amplitudes, and it is given by
$$A_4=\frac{\kappa ^{2/3}}{4(4\pi )^{10/3}}\mathrm{Tr}(T_1T_2)\mathrm{Tr}(T_3T_4)\frac{1}{\sqrt{s}}\frac{1}{\mathrm{sinh}(\sqrt{s}d)}t_8F^4,$$
(3.2)
where
$`t_8F^4`$ $`=`$ $`\{2ut(ϵ_1ϵ_2)(ϵ_3ϵ_4)2st(ϵ_1ϵ_3)(ϵ_2ϵ_4)2su(ϵ_1ϵ_4)(ϵ_2ϵ_3)`$ (3.3)
$`+(ϵ_1ϵ_2)\left[4t(ϵ_3k_1)(ϵ_4k_2)+4u(ϵ_3k_2)(ϵ_4k_1)\right]`$
$`+(ϵ_3ϵ_4)\left[4t(ϵ_1k_3)(ϵ_2k_4)+4u(ϵ_1k_4)(ϵ_2k_3)\right]`$
$`+(ϵ_1ϵ_3)\left[4s(ϵ_2k_3)(ϵ_4k_1)+4t(ϵ_2k_1)(ϵ_4k_3)\right]`$
$`+(ϵ_2ϵ_4)\left[4s(ϵ_1k_4)(ϵ_3k_2)+4t(ϵ_1k_2)(ϵ_3k_4)\right]`$
$`+(ϵ_1ϵ_4)\left[4s(ϵ_2k_4)(ϵ_3k_1)+4u(ϵ_2k_1)(ϵ_3k_4)\right]`$
$`+(ϵ_2ϵ_3)[4s(ϵ_1k_3)(ϵ_4k_2)+4u(ϵ_1k_2)(ϵ_4k_3)]\}.`$
In the limit of low momenta (i.e. $`s,t,u`$ small) the above amplitude becomes
$$A_4\frac{\kappa ^{2/3}}{4(4\pi )^{10/3}}\mathrm{Tr}(T_1T_2)\mathrm{Tr}(T_3T_4)\frac{1}{sd}t_8F^4.$$
(3.4)
In order to compare this to the heterotic string we need to rewrite
$$\frac{\kappa ^{2/3}}{d}=\frac{\kappa ^{2/3}}{\pi R_{11}l_p}=\frac{(2\pi l_p)^3}{\pi R_{11}l_p}=\frac{2}{(2\pi l_p)^4R_{11}}$$
(3.5)
in string units, using the standard relation between $`l_p`$ and the heterotic string length scale, $`l_s^H`$,
$$l_p^2=R_{11}\left(l_s^H\right)^2.$$
(3.6)
We also need the relation between $`R_{11}`$ and the coupling constant, $`g_E=e^{\varphi ^E}`$, of the ten-dimensional heterotic $`E_8\times E_8`$ (HE) string theory ,
$$R_{11}=g_E^{2/3}=e^{\frac{2\varphi ^E}{3}}.$$
(3.7)
The amplitude in heterotic string units is then given by
$$A_4^{HE}(2\pi l_s^H)^4e^{2\varphi ^E}\mathrm{Tr}(T_1T_2)\mathrm{Tr}(T_3T_4)\frac{1}{s}t_8F^4,$$
(3.8)
which agrees with the tree-level heterotic result in the low-momentum limit .
This agreement should not seem surprising in view of the fact that the low-momentum tree-level heterotic string amplitude is equal to the tree-level amplitude of ten-dimensional $`N=1`$ supergravity coupled to $`E_{8L}\times E_{8R}`$ $`N=1`$ supersymmetric Yang–Mills theory. In the small-$`R_{11}`$ limit (the heterotic string weak coupling limit) the two Hořava–Witten boundaries coincide and our tree-level field theory calculation becomes ten-dimensional. The slightly subtle point is that in the Hořava–Witten case the full tree amplitude is given by the sum of ‘pole diagrams’ which only involve three-point vertices while the ten-dimensional supergravity calculation also involves the contribution of a ‘contact’ interaction of the form $`\omega _{MNP}\omega ^{MNP}`$ (the four gauge particle contact term does not contribute to the particular group theory factor that is of interest to us). In fact, it is easy to see that this contact term is reproduced precisely by the second contribution $`A_4^{(C)\mathrm{\hspace{0.17em}2}}`$ in appendix B.
## 4 Compactification on $`S^1`$
Upon compactification on a circle of radius $`R_9l_p`$ in the $`x^9`$ direction the gauge group may be broken by the introduction of Wilson lines. We will here consider the case in which the unbroken group is $`SO(16)_L\times SO(16)_R`$. The two nine dimensional heterotic strings with this gauge group are T-dual to each other. This system has a description in Type IA language in which the Hořava–Witten boundaries are represented by two orientifold eight-planes separated by $`\pi R_{11}l_p`$ with eight $`D8`$-branes and their images placed coincident with each of them.
The adjoint representation of $`E_8`$ decomposes into the adjoint and one of the two chiral spinor representations of the $`SO(16)`$ subgroup, $`\mathrm{𝟐𝟒𝟖}=\mathrm{𝟏𝟐𝟎}+\mathrm{𝟏𝟐𝟖}`$. The breaking of the $`E_8`$ symmetry is accompanied by the generation of a mass for the chiral spinor $`\mathrm{𝟏𝟐𝟖}`$ which is a state with Kaluza–Klein charge $`l=\pm 1/2,\pm 3/2,\mathrm{}`$. The chirality of these states is independent of the signs of the charges. In eleven-dimensional units, the mass of the spinor states is given by
$$M_S=\frac{|l|}{R_9l_p}.$$
(4.1)
We will mainly consider external states that are of lowest mass, and therefore have $`l=\pm 1/2`$. We shall also choose to satisfy $`M_S<<M_{planck}`$; in this case we may define a ‘low energy’ limiting effective field theory in which the spinor states survive but Planck-scale or string-scale excitations can be ignored. In the Type IA description a spinor state corresponds to a single D-particle or anti D-particle stuck on an orientifold planes. In terms of the eleven-dimensional moduli, the Type IA string coupling constant $`g_{IA}`$ is given as ,
$$g_{IA}=R_9^{\frac{3}{2}},$$
(4.2)
and the Type I string length, $`l_s^I`$, is given in terms of the Planck length by
$$\frac{l_s^I}{l_p}=g_{IA}^{\frac{1}{3}}=R_9^{1/2}.$$
(4.3)
Thus in terms of the Type IA theory, the mass $`M_S`$ for $`l=\pm 1/2`$ becomes
$$M_S=\frac{1}{2R_9l_p}=\frac{1}{l_s^I}\frac{1}{2g_{IA}},$$
(4.4)
which is indeed the mass of a stuck D-particle (i.e., of a $`D_S`$).
### 4.1 The different kinematical regimes
Since the $`SO(16)`$ spinor states carry Kaluza–Klein charge $`l=\pm 1/2`$ there are three distinct classes of four-point functions in the nine-dimensional theory. These are characterized by which of the incoming and outgoing particles are $`D_S`$ and which are $`\overline{D}_S`$. We are considering the scattering process as a T-channel process in which $`K_1`$ and $`K_4`$ are the nine-dimensional momenta of the incoming particles that scatter into outgoing particles with momenta $`K_2`$ and $`K_3`$. These processes are the following: (a) $`D_S+D_SD_S+D_S`$, which has total $`S`$-channel Kaluza–Klein charge $`l=0`$. At low velocity this becomes a BPS configuration that preserves supersymmetry. (b) $`D_S+\overline{D}_SD_S+\overline{D}_S`$ also has $`l=0`$ but is far from BPS at low velocity. (c) $`D_S+\overline{D}_S\overline{D}_S+D_S`$ has $`l=1`$. Although this is far from BPS at low T-channel velocity the S-channel process that is related by crossing symmetry is BPS at low S-channel velocity. In this process a bulk D0-brane with mass $`2M_S`$ is exchanged.
The processes obtained by interchanging all $`D_S`$’s with $`\overline{D}_S`$’s are trivially related to these three processes and need not be considered separately.
We will describe the kinematics in the nine-dimensional centre of mass frame with the outgoing states scattering through an angle $`\theta `$ relative to the incoming states. Writing the ten-dimensional momentum of the $`r`$’th particle as $`k_r=(K_r,q_9)`$, where $`q_9`$ is the Kaluza–Klein momentum in the circular $`x^9`$ direction, gives the explicit expressions
$$k_1=(E,\mathrm{𝟎},0,p,M_S)(K_1,M_S),$$
(4.5)
$$k_2=(E^{},\mathrm{𝟎},p\mathrm{sin}\theta ,p\mathrm{cos}\theta ,(1)^lM_S)(K_2,(1)^lM_S),$$
(4.6)
$$k_3=(E^{},\mathrm{𝟎},p\mathrm{sin}\theta ,p\mathrm{cos}\theta ,(1)^{l+r}M_S)(K_3,(1)^{l+r}M_S),$$
(4.7)
$$k_4=(E,\mathrm{𝟎},0,p,(1)^rM_S)(K_4,(1)^rM_S),$$
(4.8)
where $`\mathrm{𝟎}`$ represents the zero momentum in the six dimensions transverse to the scattering plane, $`p`$ is defined by
$$p=\frac{M_Sv}{(1v^2)^{1/2}},$$
(4.9)
and
$$E^2=\frac{M_S^2}{(1v^2)}=E^2.$$
(4.10)
The distinct choices $`r=0,1`$ and $`0lr`$ allow for the three processes defined above.
The nine-dimensional Mandelstam variables, $`S=(K_1+K_2)^2`$, $`T=(K_1+K_4)^2`$ and $`U=(K_1+K_3)^2`$ are given by
$$S=\frac{2M_S^2v^2}{(1v^2)}(1\mathrm{cos}\theta ),T=4M_S^2+\frac{4M_S^2v^2}{(1v^2)},U=\frac{2M_S^2v^2}{(1v^2)}(1+\mathrm{cos}\theta ).$$
(4.11)
These satisfy the mass-shell condition $`S+T+U=4M_S^2`$.
The corresponding expressions for the ten-dimensional Mandelstam invariants, $`s,t,u`$, depend importantly on which of the three kinds of process is being considered. Given these identifications, the nine-dimensional amplitude can be obtained directly from the ten-dimensional expression (3.2).
(a) $`l=0`$, $`r=0`$: $`D_S+D_SD_S+D_S`$
In this case the ten-dimensional kinematic invariants are
$$s=S,t=T4M_S^2=\frac{4M_S^2v^2}{(1v^2)},u=U.$$
(4.12)
This is the process that is near-BPS at low energy. Substituting (4.12) into (3.2) gives an amplitude for $`D_S`$ scattering proportional to
$$\frac{2M_S^2v^2}{(1v^2)}\left(2\frac{(1+\mathrm{cos}(\theta ))}{(1\mathrm{cos}(\theta ))}(ϵ_1ϵ_2)(ϵ_3ϵ_4)+2(ϵ_1ϵ_3)(ϵ_2ϵ_4)(1+\mathrm{cos}(\theta ))(ϵ_1ϵ_4)(ϵ_2ϵ_3)\right),$$
(4.13)
where we have chosen polarisation vectors that satisfy $`k_rϵ_s=0`$ for all $`r,s=1,\mathrm{},4`$. This expression vanishes in the limit of zero velocity, $`v0`$.
In this case the amplitude is closely related to the familiar Type IIA D-particle scattering amplitude studied in . In the string theory description the amplitude is defined in the impact parameter representation, where the back reaction of the interaction on the particle trajectories is ignored. This is valid if the D-particles are assumed to move in straight lines with slow relative velocity $`2v`$ and are separated by a transverse distance $`b`$. In the standard Type IIA theory the metric on the moduli space of two $`D0`$’s is flat and the leading velocity-dependent term in the amplitude is of order $`v^4`$.
We are here concerned with the situation in which the two D-particles are constrained to lie on different fixed orientifold planes so the distance of separation is at least $`\pi R_{11}l_p`$. In order to relate the Feynman diagram we considered earlier to the standard D-particle results of the amplitude must be expressed in terms of $`b`$ and $`v`$ instead of the Mandelstam variables $`S`$ and $`T`$. This amounts to replacing $`T`$ by its expression in terms of $`v`$ using (4.12) and Fourier transforming with respect to the momentum transfer in the seven directions transverse to the particle trajectories. In the special frame defined by (4.5)-(4.8) this momentum transfer has been rotated into the vector $`𝐩(\mathrm{𝟎},p\mathrm{sin}\theta )`$ and $`S|𝐩|^2`$. Thus, the impact parameter representation is defined by
$$\stackrel{~}{A}_4(v,b)const.𝑑S(S)^{5/2}A_4(S,T)e^{i\sqrt{S}b},$$
(4.14)
where the overall constant includes the volume of $`S^6`$ to account for the integration over angular components of the momentum transfer. The transverse separation of the two $`D_S`$ particles is $`|\pi R_{11}l_p+ib|`$ and for $`R_{11}0`$ the integral in (4.14) is dominated by the region in which $`0<<\sqrt{S}R_{11}l_p<<M_Sv`$, where the factor of $`1/\mathrm{sinh}(\pi \sqrt{S}R_{11}l_p)`$ in $`A_4`$ can be approximated by $`2e^{\pi \sqrt{S}R_{11}l_p}`$. Ignoring the deflection of the particles amounts to picking out the leading term in the $`1/|\pi R_{11}l_p+ib|`$ expansion of the amplitude and it is easy to see that this leading behaviour has a coefficient proportional to $`v^4`$ which comes from the factor of $`tu`$ associated with the $`ϵ_1ϵ_2ϵ_3ϵ_4`$ term.
This behaviour is in agreement with the string calculation that can be performed as in . In this case the relevant amplitudes consist of the sum of three kinds of world-sheets. The first diagram is the cylinder that describes the overlap between the two D-particle boundary states representing the two $`D_S`$’s. The second kind of diagram consists of the Möbius strips that describes the overlap involving one $`D0`$ boundary state and the crosscap associated with the mirror image of that particle in the other orientifold plane. Finally, there are the cylinder diagrams that describe the overlap of either D-particle with any of the sixteen $`D8`$-brane boundary states. It is easy to argue that only the $`D0`$-$`D0`$ contribution can depend on $`v`$ since the velocity is constrained to be parallel to the orientifold plane and $`D8`$. The $`D0`$-$`D0`$ contribution is therefore identical to the bulk term calculated in and behaves as $`v^4`$ as expected.<sup>4</sup><sup>4</sup>4The leading contribution in the case of a bulk $`D0`$ moving transverse to the orientifold planes is proportional to $`v^2`$. This reflects the reduced amount of supersymmetry in the system.
(b) $`l=0`$, $`r=1`$: $`D_S+\overline{D}_SD_S+\overline{D}_S`$
In this case the ten-dimensional kinematic invariants are given by
$$s=S,t=T,u=U4M_S^2=4M_S^2\frac{2M_S^2v^2}{(1v^2)}(1+\mathrm{cos}\theta ).$$
(4.15)
The non-BPS amplitude that results by substituting these expressions into (3.2) behaves to leading order in $`v`$ as
$$\frac{8M_S^2(1v^2)}{v^2(1\mathrm{cos}(\theta ))}(ϵ_1ϵ_2)(ϵ_3ϵ_4),$$
(4.16)
which diverges as $`v0`$. This divergence is presumably related to the instability of the system to decay into a non-BPS D-string which will not be discussed further here.
(c) $`l=1`$, $`r=1`$: $`D_S+\overline{D}_S\overline{D}_S+D_S`$
The ten-dimensional kinematic invariants are
$$s=S4M_S^2=4M_S^2\frac{2M_S^2v^2}{(1v^2)}(1\mathrm{cos}\theta ),t=T,u=U.$$
(4.17)
This is the case in which there is $`D0`$ exchange and the behaviour of the amplitude is qualitatively different from the cases with $`l=0`$ since $`s4M_S^2`$ is large in the low-velocity limit, $`v0`$. The denominator of (3.2) can then be expanded as
$$\frac{1}{\mathrm{sinh}(\pi \sqrt{s}R_{11}l_p)}=2e^{2\pi M_SR_{11}l_p}\underset{r=0}{\overset{\mathrm{}}{}}e^{4\pi rM_SR_{11}l_p}.$$
(4.18)
Using the expression for $`M_S`$ from (4.4), (4.18) then becomes
$$2\underset{r=0}{\overset{\mathrm{}}{}}e^{\frac{(2r+1)\pi R_{11}}{R_9}}=2\underset{r=0}{\overset{\mathrm{}}{}}e^{(2r+1)\pi \frac{R^{IA}}{g_{IA}}}=2\underset{r=0}{\overset{\mathrm{}}{}}e^{\frac{(2r+1)\pi }{g_I}},$$
(4.19)
where we have used the standard relation between the M-theory variables and the ten-dimensional Type IA coupling constant, $`g_{IA}=e^{\varphi ^{IA}}`$, and radius, $`R_{IA}`$,
$$R_{IA}=R_{11}R_9^{1/2},g_{IA}=R_9^{3/2},$$
(4.20)
together with the familiar T-duality relation between Type IA and Type I,
$$g_I=\frac{g_{IA}}{R_{IA}}=\frac{R_9}{R_{11}},R_I=\frac{1}{R_{IA}}=\frac{1}{R_{11}R_9^{1/2}}.$$
(4.21)
In order to rewrite the total amplitude (3.2) in terms of Type IA (or Type I) variables, we also have to re-express the dimensionful quantities in terms of the string scale. As before, the amplitude of the original external states is proportional to $`\kappa ^{2/3}`$ which is now equal to
$$\kappa ^{2/3}=\frac{1}{(2\pi l_p)^3}=\frac{1}{(2\pi l_s^I)^3}e^{\varphi ^{IA}}.$$
(4.22)
Together with $`\sqrt{s}=2M_S=(g_{IA}l_s^I)^1`$ the amplitude thus becomes
$$A_4^I=\frac{1}{(2\pi l_s^I)^2}\frac{1}{(4\pi )^{13/3}}\mathrm{Tr}(T_1T_2)\mathrm{Tr}(T_3T_4)t_8F^4\underset{r=0}{\overset{\mathrm{}}{}}e^{\frac{(2r+1)\pi }{g_I}}.$$
(4.23)
This expansion as a series of exponentially suppressed terms is relevant at weak Type I coupling, $`g_I=R_9/R_{11}<<1`$, where the leading term dominates. This term can be identified as the effect of a non-BPS D-instanton with action $`\pi /g_I`$.<sup>5</sup><sup>5</sup>5As we shall demonstrate in the next section, this process also exhibits the key characteristic of the D-instanton of Type I: it breaks $`O(32)`$ to $`SO(32)`$. According to the Type I D-instanton can be thought of as the sum of the Type IIB D-instanton and the anti-D-instanton (without any prefactors). The orientation-reversing operator $`\mathrm{\Omega }`$ maps the D-instanton boundary state to the anti-D-instanton boundary state, and therefore only the sum of the two boundary states is invariant. Under T-duality, the D-instanton becomes a $`D0`$-brane whose world-line stretches across the interval between the orientifold planes of the nine-dimensional Type IA theory, and the anti-D-instanton becomes the configuration in which the world-line has the opposite orientation. Again this can be represented by two boundary states that are mapped into one another under the T-dual of $`\mathrm{\Omega }`$, $`\mathrm{\Omega }_{11}`$. The D-particle world-line transfers $`D0`$ Ramond–Ramond charge from one orientifold plane to the other, while the image under $`\mathrm{\Omega }_{11}`$ transfers the opposite $`D0`$ Ramond-Ramond charge. In the process $`D0+\overline{D}0\overline{D}0+D0`$ only of these two processes contributes; the conjugate process $`\overline{D}0+D0D0+\overline{D}0`$ would pick out the other contribution.
In our conventions, the mass of a bulk D-particle in Type IA is $`1/g_{IA}`$, and comparison with Type IIA gives $`g_{IA}=g_{IIA}`$. On the other hand, under T-duality, the stuck D-particle of Type IA (with mass $`1/(2g_{IA})`$) becomes a D1-brane of Type I that wraps the circle $`R_I`$ once. This has a mass $`R_I/(2g_I)`$ and therefore $`2g_I=g_{IIB}`$. Using this, the action of the type I D-instanton, that is the modulus of the exponent of the $`r=0`$ term in (4.23), has the value
$$\frac{\pi }{g_I}=\frac{2\pi }{g_{IIB}},$$
(4.24)
in Type IIB units. This is precisely the action of a single D-instanton of Type IIB . The agreement is a consequence of the BPS nature of the exchanged type IIA D-particle.
The terms with $`r>0`$ in (4.23) describe processes where the D-particle is emitted on one boundary and absorbed at the other, but where it winds an integer number of times around the compact circle. These terms appear as multi-instanton contributions with odd D-instanton number. However, all of these multi-instanton terms carry the same topological $`𝖹𝖹_2`$ charge, and the contributions with $`r>0`$ are thus presumably ‘unstable’ and therefore only the leading contribution with $`r=0`$ should be taken seriously in the perturbative Type I limit. However, the complete series is crucial in determining the amplitude in the limit $`g_I\mathrm{}`$, which is the weak coupling limit of the $`SO(32)`$ heterotic string theory (the HO theory) compactified on a circle with appropriate Wilson lines to break the symmetry to $`SO(16)\times SO(16)`$. The heterotic coupling is given by $`g_{HO}=1/g_I`$ so that the amplitude (4.23) has an infinite power series expansion in heterotic string perturbation theory. After taking into account the rescaling of the string scale
$$l_s^I=l_s^H\left(\frac{R_{11}}{R_9}\right)^{1/2}=l_s^Hg_{HO}^{1/2},$$
(4.25)
needed to pass from the Type I to the heterotic $`SO(32)`$ theory, the amplitude has an expansion in powers of $`g_{HO}^2`$,
$$A_4^{HO}\frac{1}{g_{HO}^2(2\pi l_s^H)^2}\mathrm{Tr}(T_1T_2)\mathrm{Tr}(T_3T_4)t_8F^4(1+O(g_{HO}^2)),$$
(4.26)
where we have used
$$\underset{r=0}{\overset{\mathrm{}}{}}e^{\frac{(2r+1)\pi }{g_I}}=\frac{1}{2\mathrm{sinh}(\pi g_{HO})}=\frac{1}{2\pi g_{HO}}(1+O(g_{HO}^2)).$$
(4.27)
The first term in (4.26) is a tree-level term that can be determined directly from the tree-level heterotic string compactified on a circle with the appropriate Wilson lines to give the unbroken $`SO(16)\times SO(16)`$. The $`O(g_{HO}^2)`$ corrections appear to correspond to loop corrections to the amplitude, but these might also get contributions from M-theory loops that have not been considered in this paper.
The M-theory tree-level calculation can be trusted whenever both $`R_{11}`$ and $`R_9`$ are large. In particular, this implies that the ten-dimensional Type IA coupling is large, and that the radius of the Type I theory is small. On the other hand, the ten-dimensional Type I coupling only depends on the ratio of $`R_9`$ and $`R_{11}`$, and therefore can be small. However this does not imply that one can trust the Type I perturbation. In fact, $`R_I=g_I/g_{IA}<<g_I`$ and therefore the condensation of closed-string winding states cannot be ignored. This can also be seen from the fact that both the radius and the ten-dimensional coupling constant of the dual Type IA theory are large.
If we had been working in a regime in which perturbative string theory could be trusted there would have been a paradox. We would have been led to believe that the Type I D-instanton is only a ‘stable’ solution (i.e., a solution with no tachyonic modes) provided the radius $`R_I`$ is sufficiently large. Within the perturbative approximation, when $`R_I<\sqrt{2}`$ the open string beginning and ending on the D-instanton develops a tachyonic mode and the D-instanton should ‘decay’ into a non-BPS $`D0`$-$`\overline{D}0`$ pair whose world-lines stretch along the circle.<sup>6</sup><sup>6</sup>6The K-theory class that corresponds to the D-instanton is however non-trivial for all values of the radii. The question of stability is the question of whether the representative of the non-trivial K-theory class that has least action is a pointlike instanton or an extended (one-dimensional) object. Furthermore, these $`D0`$ and $`\overline{D}0`$ branes carry a relative $`𝖹𝖹_2`$ Wilson line. Under T-duality, the two $`𝖹𝖹_2`$ $`D0`$ branes become two $`𝖹𝖹_2`$ D-instantons of Type IA, and the relative Wilson line means that they are located at opposite orientifold planes. From the point of view of Type IA, the above superposition of $`D0`$ world-lines would therefore be unstable to decay into two non-BPS Type IA D-instantons provided that $`R_{IA}>1/\sqrt{2}`$. It is easy to see that such a configuration is not at all similar to the exchange of a D-particle that is contained in the M-theory amplitude.
However, the region of validity of our M-theory argument is one of strong coupling in the Type IA theory, in which the preceding stability arguments are not valid. The fact that the D-instanton action extracted from the spinor scattering process corresponds to that expected for large $`R_I`$ in the Type I theory suggests that the non-BPS Type I D-instanton remains stable in the non-perturbative region $`R_I<<g_I`$. The reason why the M-theory calculation reproduces the correct value for the instanton action is related to the fact that for the process we are considering, only one term in the superposition of the D-instanton and the anti-D-instanton contributes. The corresponding D-instanton action is therefore protected by supersymmetry since the bulk D0-brane that is exchanged is a BPS state of Type IIA (and the anti-D0-brane does not play a role). We thus expect that the exponent in the amplitude can be trusted beyond the original regime of validity. Although the limit which gives ten-dimensional weakly coupled Type I theory reproduces the correct normalisation of the instanton action it will not give the correct value for the coefficient of the exponential.
### 4.2 The Type I D-instanton and the breaking of $`O(32)`$
From its description in terms of K-theory, it is clear that the D-instanton is associated with the fact that the gauge group of Type I string theory is $`SO(32)`$ rather than the $`O(32)`$ that might have been expected on the basis of perturbation theory .<sup>7</sup><sup>7</sup>7The actual gauge group is obviously $`Spin(32)/𝖹𝖹_2`$. The issue of the correct spin cover will be discussed later. Here we will demonstrate this by considering the dual Type IA theory in the standard $`SO(16)\times SO(16)`$ vacuum. Consideration of the perturbative approximation to the low energy effective field theory would suggest that the symmetry group could be $`O(16)\times O(16)`$. However, we will see that the instantonic contribution to the scattering of spinor states is not invariant under the disconnected component of either of the two $`O(16)`$ groups.
We shall concentrate on the $`O(16)`$ group that is associated to the left orientifold plane since the argument for the other gauge group is identical. In the vacuum we are considering, eight $`D8`$-branes (plus eight mirror $`D8`$-branes) are located at each of the two orientifold planes. As we have seen, the instantonic contribution to the scattering amplitude comes from the diagram where a $`D0`$ (or $`\overline{D0}`$) is emitted from the left orientifold plane and absorbed at the right orientifold plane. In this process, the $`D0`$ (and its mirror partner) has to cross the eight $`D8`$-branes (and their mirror partners) that are localised at the fixed plane. However, whenever a $`D0`$ crosses a $`D8`$-brane, a fundamental string is created that stretches between the $`D0`$ and the $`D8`$-brane . This is an important effect as we shall show momentarily.
In order to obtain a clear picture, it is useful to consider the configuration where the eight $`D8`$-branes have been moved to $`x^{11}=ϵ`$, away from the orientifold fixed plane at $`x^{11}=0`$ (with corresponding displacements of the mirror eight-branes to $`x^{11}=ϵ`$). In this configuration, the ‘cosmological constant’ $`m`$ takes the value $`m=8`$ for $`0<x^{11}<ϵ`$ and $`m=0`$ for $`|x^{11}|>ϵ`$. As was shown by , in a background characterised by $`m`$, each $`D0`$ has to carry $`|m|`$ strings, where the sign of $`m`$ determines whether the string begins or ends on the $`D0`$. The $`D0`$ (and its mirror) that are emitted from the orientifold plane each carry eight strings in the regime $`|x^{11}|<ϵ`$. These strings can be thought of as stretching between the bulk $`D0`$ (that has been emitted from the orientifold plane) and the remaining fractional $`D0`$ that is still stuck. As the bulk $`D0`$ (and its mirror) crosses the eight $`D8`$-branes (and their mirrors), additional strings are created that join it to the $`D8`$-branes. These can recombine with the original strings to give eight strings that stretch between the stuck $`D0`$ and the eight $`D8`$-branes, as well as eight strings that stretch between the stuck $`D0`$ and the eight mirror $`D8`$-branes. In the final configuration the bulk $`D0`$ does not have any strings attached to it, which is consistent with the fact that $`m=0`$ in the interval between the two fixed planes.
The sixteen different strings that stretch between the stuck $`D0`$ and the eight plus eight $`D8`$-branes can be identified with the generators $`\gamma _i`$, $`i=1,\mathrm{},16`$ of the Clifford algebra
$$\{\gamma _i,\gamma _j\}=2\delta _{ij}$$
(4.28)
that give rise to the Lie algebra elements of $`so(16)`$ by
$$\sigma _{ij}=\frac{1}{4}[\gamma _i,\gamma _j].$$
(4.29)
Indeed, the elements of the gauge group correspond to 8-8 strings, and are therefore bilinear in the 0-8 strings that are described by the $`\gamma _i`$. The product of all sixteen 0-8 strings (that arises naturally in the above configuration) is then the chirality operator of $`O(16)`$,
$$\mathrm{\Gamma }=\underset{i=1}{\overset{16}{}}\gamma _i.$$
(4.30)
Thus it follows that the vertex that describes the emission of a single $`D0`$ (or $`\overline{D}0`$) is actually given by
$$\mathrm{Tr}\left(S_1S_2\mathrm{\Gamma }\right),$$
(4.31)
where $`S_i`$ describes the state in the spinor representation for the two fractional $`D0`$’s. More generally, if $`m`$ bulk $`D0`$’s are emitted, the vertex will be
$$\mathrm{Tr}\left(S_1S_2\mathrm{\Gamma }^m\right),$$
(4.32)
since each bulk $`D0`$ gives rise to one chirality operator.
The chirality operator $`\mathrm{\Gamma }`$ is invariant under conjugation by any element of $`SO(16)`$, but it changes sign under conjugation by an element $`g_0`$ in the disconnected component of $`O(16)`$, $`g_0\mathrm{\Gamma }g_0^1=\mathrm{\Gamma }`$. This implies that the vertex with odd $`m`$ only respects the $`SO(16)`$ subgroup of $`O(16)`$. Since the contributions with odd $`m`$ are precisely those that appear in the instantonic contribution to the amplitude, this exponentially suppressed contribution is associated with the breaking of $`O(16)`$ to $`SO(16)`$. In terms of Type I, $`SO(16)\times SO(16)`$ is a subgroup of $`SO(32)`$. The large $`O(32)`$ transformations that are not in $`SO(32)`$ can be taken to lie in $`O(16)\times SO(16)`$; these transformations are indeed not respected by the D-instanton, as was stressed in .
The actual gauge group of the ten-dimensional Type I theory is $`Spin(32)/𝖹𝖹_2`$. Upon compactification to nine dimensions with the above Wilson line, this is broken to $`S(Pin(16)\times Pin(16))/𝖹𝖹_2`$, where the $`S`$ outside the bracket indicates that either both the elements of $`Pin(16)\times Pin(16)`$ are in the disconnected component of $`Pin(16)`$ or neither of them are. As we have just explained, the effect of the D-instanton of Type I is to impose this constraint. The actual gauge group is somewhat smaller since the nine-dimensional Type IA theory also has a D-instanton that is T-dual to the wrapped D0-brane of Type I. This instanton is stuck on one of the two fixed planes, and it is responsible for breaking each $`Pin(16)`$ separately to $`Spin(16)`$. (This follows by essentially the same arguments that were used by Witten for the case of the ten-dimensional Type I D-instanton.) Thus the actual gauge group is
$$\left(Spin(16)\times Spin(16)\right)/𝖹𝖹_2,$$
(4.33)
where the element in the centre by which the group is divided is the element that does not allow states transforming in the vector representation of one $`Spin(16)`$ and the scalar representation of the other. It is not difficult to see that all allowed (and no other) conjugacy classes of representations of this gauge group are actually present in both nine dimensional heterotic theories. (This discussion is relevant for some issues raised in .)
## 5 Discussion
In this paper we have identified the D-instanton contribution to a certain amplitude involving spinor states of Type I, using the relation of nine-dimensional Type I to M-theory. This relation involves Type IA, and it is therefore natural to ask whether the $`𝖹𝖹_2`$ non-BPS D-instanton of Type IA (that is T-dual to the D0-brane of Type I with its world-line wrapped along $`x^9`$) can also be understood in terms of M-theory (compare also ). The wrapped D0-brane of Type I can be obtained from a D1-brane anti-D1-brane pair that are wrapped around $`x^9`$ and $`x^8`$ say, and that have a relative Wilson line along $`x^8`$. Under T-duality of the $`x^9`$-circle, we therefore obtain a D0-brane and an anti-D0-brane, both localised on the same orientifold plane and wrapped around $`x^8`$ with a relative Wilson line. One should therefore expect that the Type IA D-instanton contributes to the scattering diagram $`D_S+\overline{D}_S\overline{D}_S+D_S`$ where now all four D-particles are localised on the same boundary. An analysis of this process can be made that is analogous to that of the main part of the paper, with the important difference that the propagator (2.4) does not involve a $`(1)^m`$, and therefore becomes
$$\widehat{𝒢}(p^M;0,0)=\frac{1}{\pi R_{11}l_p}\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\frac{1}{p^2+\frac{m^2}{R_{11}^2l_p^2}}=\frac{1}{\sqrt{p^2}}\underset{n𝖹𝖹}{}e^{2|n|\pi \sqrt{p^2}R_{11}l_p}.$$
(5.34)
The term with $`n=0`$ should correspond to a part of the contribution of the Type IA $`𝖹𝖹_2`$ D-instanton. The terms with $`n0`$ come from multiple windings of the $`D0`$ world-line and, as in the earlier case, carry the same $`𝖹𝖹_2`$ topological charge as the $`n=0`$ term; the latter is therefore the dominant contribution. Thus the tree amplitude describing the exchange of the graviton and the three-form in Type IA variables is approximated by
$$A_4=\frac{1}{2(2\pi l_s^I)^2}\frac{1}{(4\pi )^{13/3}}t_8F^4\left[\mathrm{Tr}(T_1T_2)\mathrm{Tr}(T_3T_4)+\mathrm{}\right],$$
(5.35)
where we have only exhibited explicitly the term whose group structure is the same as in the main part of the paper. By essentially the same arguments as above, the coupling of the spinors to the bulk graviton and three-form should again involve a chirality operator and therefore break $`O(16)`$ to $`SO(16)`$. However, the Type IA instanton should contribute a term of the form $`e^{c/g_{IA}}`$ to the amplitude (for some constant, $`c`$). Since the M-theory calculation is justified only in the region of large large $`R_9,R_{11}`$, where $`g_{IA}>>1`$, the expression (5.35) can only represent the first term in an expansion of the exponent in powers of $`1/g_{IA}`$. The higher-order terms will depend on the undetermined loop diagrams of M-theory.
## Acknowledgements
We acknowledge partial support from the PPARC SPG programme, “String Theory and Realistic Field Theory”, PPA/G/S/1998/00613. TD is supported by a Scholarship of Trinity College, Cambridge, and MRG is supported by a Royal Society University Research Fellowship. MRG and MBG are grateful to the California Institute of Technology for hospitality during the final stages of this work. We also thank Oren Bergmann and Michael Dine for useful conversations.
## Appendix
## Appendix A The Hořava–Witten action
The spinor indices are written as $`\alpha ,\beta ,\gamma `$. The supergravity multiplet has the graviton $`g`$, the gravitino $`\psi _{\mu \alpha }`$, and a three-form $`C`$ with the field strength $`G`$ (in component form $`G_{\mu \nu \rho \sigma }=4!_{[\mu }C_{\nu \rho \sigma ]}`$). The spinors are Majorana and $`\overline{\psi }_\alpha `$ is defined by $`\overline{\psi }_\alpha =C_{\alpha \beta }\psi ^\beta `$ where $`C_{\alpha \beta }`$ is the charge conjugation matrix. The $`32\times 32`$ real Dirac matrices $`\mathrm{\Gamma }_\mu `$ satisfy the Clifford algebra $`\{\mathrm{\Gamma }_\mu ,\mathrm{\Gamma }_\nu \}=2g_{\mu \nu }`$, and we define $`\mathrm{\Gamma }^{\mu _1\mu _2\mathrm{}\mu _n}\frac{1}{n!}\mathrm{\Gamma }^{\mu _1}\mathrm{\Gamma }^{\mu _2}\mathrm{}\mathrm{\Gamma }^{\mu _n}\pm `$ permutations. With these conventions the $`D=11`$, $`N=1`$ bulk supergravity action is given by ,
$`S_{bulk}`$ $`=`$ $`{\displaystyle \frac{1}{\kappa ^2}}{\displaystyle _{R^{1,9}\times S^1}}d^{11}x\sqrt{g}\{{\displaystyle \frac{R}{2}}{\displaystyle \frac{1}{2}}\overline{\psi }_\mu \mathrm{\Gamma }^{\mu \nu \rho }D_\nu \psi _\rho {\displaystyle \frac{1}{48}}G_{\mu \nu \rho \sigma }G^{\mu \nu \rho \sigma }`$
$`{\displaystyle \frac{\sqrt{2}}{192}}(\overline{\psi }_\mu \mathrm{\Gamma }^{\mu \nu \rho \sigma \tau \lambda }\psi _\lambda +12\overline{\psi }^\nu \mathrm{\Gamma }^{\rho \sigma }\psi ^\tau )G_{\nu \rho \sigma \tau }{\displaystyle \frac{\sqrt{2}}{3456}}ϵ^{\mu _1\mu _2\mathrm{}\mu _{11}}C_{\mu _1\mu _2\mu _3}G_{\mu _4\mathrm{}\mu _7}G_{\mu _8\mathrm{}\mu _{11}}\}.`$
Here $`\kappa ^2=(2\pi l_p)^9`$, where $`l_p`$ is the eleven-dimensional Planck length. We are not considering higher order terms quartic in the gravitino. The Riemann tensor is the field strength of the spin connection $`\mathrm{\Omega }`$. On the other hand the D=10 supersymmetric Yang-Mills action on each boundary, containing the $`E_8`$ gauge field $`A^a`$ and the gluino $`\chi ^a`$, coupled to the bulk supergravity fields, is given by ,
$`S_{YM}`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^{5/3}\kappa ^{4/3}}}{\displaystyle _{R^{1,9}}}d^{10}x\sqrt{g}[{\displaystyle \frac{1}{4}}F_{MN}^aF^{aMN}{\displaystyle \frac{1}{2}}\overline{\chi }^a\mathrm{\Gamma }^MD_M\chi ^a`$ (A.2)
$`{\displaystyle \frac{1}{4}}\overline{\psi }_M\mathrm{\Gamma }^{NP}\mathrm{\Gamma }^MF_{NP}^a\chi ^a+\overline{\chi }^a\mathrm{\Gamma }^{MNP}\chi ^a\{{\displaystyle \frac{\sqrt{2}}{48}}G_{MNP11}+{\displaystyle \frac{1}{32}}\overline{\psi }_M\mathrm{\Gamma }_{NP}\psi _{11}+{\displaystyle \frac{1}{32}}\overline{\psi }^D\mathrm{\Gamma }_{DMNP}\psi _{11}`$
$`+{\displaystyle \frac{1}{128}}(3\overline{\psi }_M\mathrm{\Gamma }_N\psi _P\overline{\psi }_M\mathrm{\Gamma }_{NPQ}\psi ^Q{\displaystyle \frac{1}{2}}\overline{\psi }_Q\mathrm{\Gamma }_{MNP}\psi ^Q{\displaystyle \frac{13}{6}}\overline{\psi }^D\mathrm{\Gamma }_{QMNPR}\psi ^R)\}],`$
where $`g`$ is here the restriction of the eleven-dimensional metric to ten dimensions, and the Yang-Mills field strength is given by $`F_{MN}^a=_MA_N^a_NA_M^a+f_{bc}^aA_M^bA_N^c`$. The Yang-Mills coupling $`\lambda `$ has been expressed in terms of the gravitational coupling $`\kappa `$ by the relation $`\lambda ^2=4\pi (4\pi \kappa ^2)^{2/3}`$ ; the prefactor is therefore $`(4\pi )^{5/3}(2\pi l_p)^6`$. The expression for $`G_{MNP\mathrm{\hspace{0.17em}11}}`$ is modified in the manner described by (2.10).
## Appendix B The separate tree amplitudes
The tree amplitude due to graviton exchange is given by
$`A_4^{(h)}`$ $`=`$ $`{\displaystyle \frac{\kappa ^{2/3}}{(4\pi )^{10/3}d}}\mathrm{Tr}(T_1T_2)\mathrm{Tr}(T_3T_4){\displaystyle \underset{m𝖹𝖹}{}}(1)^m{\displaystyle \frac{1}{s+p_{11}^2}}`$ (B.1)
$`\times \left[ϵ_1^{N_1}k_1^{M_1}ϵ_2^{Q_1}k_2^{P_1}+ϵ_1^{Q_1}k_1^{P_1}ϵ_2^{N_1}k_2^{M_1}\right]`$
$`\times \left({\displaystyle \frac{1}{2}}\eta _{M_1[Q_1}\eta _{P_1]N_1}\eta _{R_1S_1}+\eta _{M_1[P_1}\eta _{Q_1]S_1}\eta _{R_1N_1}+\eta _{M_1[R_1}\eta _{Q_1]N_1}\eta _{P_1S_1}\right)`$
$`\times \left(\eta ^{R_1R_2}\eta ^{S_1S_2}+\eta ^{R_1S_2}\eta ^{S_1R_2}{\displaystyle \frac{2}{9}}\eta ^{R_1S_1}\eta ^{R_2S_2}\right)`$
$`\times \left({\displaystyle \frac{1}{2}}\eta _{M_2[Q_2}\eta _{P_2]N_2}\eta _{R_2S_2}+\eta _{M_2[P_2}\eta _{Q_2]S_2}\eta _{R_2N_2}+\eta _{M_2[R_2}\eta _{Q_2]N_2}\eta _{P_2S_2}\right)`$
$`\times \left[ϵ_3^{N_2}k_3^{M_2}ϵ_4^{Q_2}k_4^{P_2}+ϵ_3^{Q_2}k_3^{P_2}ϵ_4^{N_2}k_4^{M_2}\right].`$
The sum over $`m`$ can be performed directly using (2.5), and after a lengthy calculation the result simplifies to
$`A_4^{(h)}`$ $`=`$ $`{\displaystyle \frac{\kappa ^{2/3}}{4(4\pi )^{10/3}}}\mathrm{Tr}(T_1T_2)\mathrm{Tr}(T_3T_4){\displaystyle \frac{1}{\sqrt{s}}}{\displaystyle \frac{1}{\mathrm{sinh}(\sqrt{s}d)}}\{`$ (B.2)
$`2ut(ϵ_1ϵ_2)(ϵ_3ϵ_4)+s^2(ϵ_1ϵ_3)(ϵ_2ϵ_4)+s^2(ϵ_1ϵ_4)(ϵ_2ϵ_3)`$
$`+(ϵ_1ϵ_2)\left[4t(ϵ_3k_1)(ϵ_4k_2)+4u(ϵ_3k_2)(ϵ_4k_1)\right]`$
$`+(ϵ_3ϵ_4)\left[4t(ϵ_1k_3)(ϵ_2k_4)+4u(ϵ_1k_4)(ϵ_2k_3)\right]`$
$`+(ϵ_2ϵ_3)\left[+2s(ϵ_1k_4)(ϵ_4k_3)+2s(ϵ_1k_2)(ϵ_4k_1)2t(ϵ_1k_2)(ϵ_4k_3)\right]`$
$`+(ϵ_1ϵ_3)\left[+2s(ϵ_2k_4)(ϵ_4k_3)+2s(ϵ_2k_1)(ϵ_4k_2)2u(ϵ_2k_1)(ϵ_4k_3)\right]`$
$`+(ϵ_2ϵ_4)\left[+2s(ϵ_1k_3)(ϵ_3k_4)+2s(ϵ_1k_2)(ϵ_3k_1)2u(ϵ_1k_2)(ϵ_3k_4)\right]`$
$`+(ϵ_1ϵ_4)\left[+2s(ϵ_2k_3)(ϵ_3k_4)+2s(ϵ_2k_1)(ϵ_3k_2)2t(ϵ_2k_1)(ϵ_3k_4)\right]`$
$`4(ϵ_1k_2)(ϵ_2k_1)(ϵ_3k_4)(ϵ_4k_3)+4(ϵ_1k_4)(ϵ_2k_1)(ϵ_3k_2)(ϵ_4k_3)`$
$`+4(ϵ_1k_2)(ϵ_2k_4)(ϵ_3k_1)(ϵ_4k_3)+4(ϵ_1k_3)(ϵ_2k_1)(ϵ_3k_4)(ϵ_4k_2)`$
$`+4(ϵ_1k_2)(ϵ_2k_3)(ϵ_3k_4)(ϵ_4k_1)\}.`$
The three-form exchange diagram involves a vertex of the form $`S^{AAC}`$ (2.12) at each end of the propagator, and thus the total amplitude is of the form
$`A_4^{(C)}`$ $`=`$ $`{\displaystyle \frac{2\kappa ^{2/3}}{(4\pi )^{10/3}(3!)^2d}}\mathrm{Tr}(T_1T_2)\mathrm{Tr}(T_3T_4){\displaystyle \underset{m𝖹𝖹}{}}(1)^m{\displaystyle \frac{1}{s+p_{11}^2}}`$ (B.3)
$`\left[ϵ_1^{M_1}k_2^{N_1}ϵ_2^{P_1}+ϵ_2^{M_1}k_1^{N_1}ϵ_1^{P_1}\right]`$
$`\left(p^{[M_2}p_{[M_1}\eta _{N_1P_1}\eta _{11]}^{N_2}\eta ^{P_2\mathrm{\hspace{0.17em}11}]}\right)\left[ϵ_3^{M_2}k_4^{N_2}ϵ_4^{P_2}+ϵ_4^{M_2}k_3^{N_2}ϵ_3^{P_2}\right],`$
where $`p=(k_1+k_2,p_{11})`$, since the external gauge bosons do not have any momentum (or polarisation vectors) in the 11th direction. The propagator therefore has two contributions
$$\left(3k_{[M_2}k^{[M_1}\eta ^{N_1P_1]}\eta _{N_2P_2]}p_{11}^2\eta ^{[M_1N_1}\eta _{[M_2}^{P_1]}\eta _{N_2P_2]}\right),$$
(B.4)
where $`k=k_1+k_2=k_3+k_4`$. In the first term we can again do the sum over $`p_{11}`$ as before, and we obtain a contribution of the form
$`A_4^{(C)\mathrm{\hspace{0.17em}1}}`$ $`=`$ $`{\displaystyle \frac{\kappa ^{2/3}}{4(4\pi )^{10/3}}}\mathrm{Tr}(T_1T_2)\mathrm{Tr}(T_3T_4){\displaystyle \frac{1}{\sqrt{s}}}{\displaystyle \frac{1}{\mathrm{sinh}(\sqrt{s}d)}}\{`$ (B.5)
$`+(ut)(ϵ_1ϵ_4)(ϵ_2k_1)(ϵ_3k_4)+(tu)(ϵ_1ϵ_3)(ϵ_2k_1)(ϵ_4k_3)`$
$`+(tu)(ϵ_2ϵ_4)(ϵ_1k_2)(ϵ_3k_4)+(ut)(ϵ_2ϵ_3)(ϵ_1k_2)(ϵ_4k_3)`$
$`+(ϵ_1k_3)(ϵ_2k_1)(ϵ_3k_4)(ϵ_4k_1)(ϵ_1k_3)(ϵ_2k_1)(ϵ_3k_4)(ϵ_4k_2)`$
$`(ϵ_1k_4)(ϵ_2k_1)(ϵ_3k_4)(ϵ_4k_1)+(ϵ_1k_4)(ϵ_2k_1)(ϵ_3k_4)(ϵ_4k_2)`$
$`+(ϵ_1k_4)(ϵ_2k_1)(ϵ_3k_1)(ϵ_4k_3)(ϵ_1k_4)(ϵ_2k_1)(ϵ_3k_2)(ϵ_4k_3)`$
$`(ϵ_1k_3)(ϵ_2k_1)(ϵ_3k_1)(ϵ_4k_3)+(ϵ_1k_3)(ϵ_2k_1)(ϵ_3k_2)(ϵ_4k_3)`$
$`+(ϵ_1k_2)(ϵ_2k_3)(ϵ_3k_4)(ϵ_4k_2)(ϵ_1k_2)(ϵ_2k_3)(ϵ_3k_4)(ϵ_4k_1)`$
$`(ϵ_1k_2)(ϵ_2k_4)(ϵ_3k_4)(ϵ_4k_2)+(ϵ_1k_2)(ϵ_2k_4)(ϵ_3k_4)(ϵ_4k_1)`$
$`+(ϵ_1k_2)(ϵ_2k_4)(ϵ_3k_2)(ϵ_4k_3)(ϵ_1k_2)(ϵ_2k_4)(ϵ_3k_1)(ϵ_4k_3)`$
$`(ϵ_1k_2)(ϵ_2k_3)(ϵ_3k_2)(ϵ_4k_3)+(ϵ_1k_2)(ϵ_2k_3)(ϵ_3k_1)(ϵ_4k_3)\}.`$
In the second term, the sum over $`m`$ is evaluated using (2.6), and we obtain
$`A_4^{(C)\mathrm{\hspace{0.17em}2}}`$ $`=`$ $`{\displaystyle \frac{\kappa ^{2/3}}{4(4\pi )^{10/3}}}\mathrm{Tr}(T_1T_2)\mathrm{Tr}(T_3T_4)\sqrt{s}{\displaystyle \frac{1}{\mathrm{sinh}(\sqrt{s}d)}}\{`$ (B.6)
$`+(tu)(ϵ_1ϵ_3)(ϵ_2ϵ_4)+(ut)(ϵ_1ϵ_4)(ϵ_2ϵ_3)`$
$`+(ϵ_1ϵ_3)\left[(ϵ_2k_3)(ϵ_2k_4)\right]\left[(ϵ_4k_2)(ϵ_4k_1)\right]`$
$`+(ϵ_1ϵ_4)\left[(ϵ_2k_3)(ϵ_2k_4)\right]\left[(ϵ_3k_1)(ϵ_3k_2)\right]`$
$`+(ϵ_2ϵ_3)\left[(ϵ_1k_3)(ϵ_1k_4)\right]\left[(ϵ_4k_1)(ϵ_4k_2)\right]`$
$`(ϵ_2ϵ_4)[(ϵ_1k_3)(ϵ_1k_4)][(ϵ_3k_1)(ϵ_3k_2)]\}.`$
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# Weighing the universe with accelerators and detectors
## I Introduction
Despite the many successes of the standard cosmology, we still do not know the composition and the amount of energy density in our universe (for a review, see for example ). The existence of matter which does not emit much light is certain from the fact that stars and other luminous matter contribute only a tiny fraction (about 0.004) of the critical energy density (required for a spatially flat universe) while the orbits of stars around galaxies (see e.g. ) indicate that the gravitating energy density is about ten times larger.<sup>*</sup><sup>*</sup>*In this paper, we will use the standard notation of ratio of energy density in $`X`$ to the critical energy density as $`\mathrm{\Omega }_X`$. Since the critical energy density is determined by the Hubble expansion rate which is uncertain, we will use the usual parameterization $`H=h\text{100 km/s/Mpc}`$ and often write the energy density as $`\mathrm{\Omega }_Xh^2`$ value instead. From a cosmological point of view, a variety of observations, such as the peculiar motions of galaxies as well as the masses of clusters of galaxies, corroborate the existence of dark energy and indicate that $`\mathrm{\Omega }_M`$, the ratio of the average matter-energy density in our Hubble volume to the critical energy density, is about $`1/3`$ if the Hubble expansion rate is about 70 km/s/Mpc as the current observations indicate (see e.g. ). Although without other constraints this dark energy density could be in the form of ordinary baryonic matter (objects made of neutrons and protons), the standard model of big bang nucleosynthesis (BBN) calculations and the measurement of primeval abundance of deuterium gives strong evidence that the baryon density is $`\mathrm{\Omega }_Bh^2=(0.02\pm 0.002)`$ which implies $`\mathrm{\Omega }_B`$ of about $`0.04`$. Hence, most of the matter energy density in the universe seems to be in the form of nonbaryonic dark matter (NBDM).
Measurements of cosmological parameters will determine $`\mathrm{\Omega }_{NBDM}`$ accurately, but are unlikely to help us know the actual character of the NBDM, since they mainly depend on its gravitational interactions. Some information on NBDM collisional energy loss and clustering will also come from astrophysics information. For perhaps the most likely forms of NBDM we can hope to observe the actual particles in laboratory experiments, and calculate to a few percent accuracy the actual contribution to $`\mathrm{\Omega }`$ .
Remarkably, the more likely extensions to the standard model of particle physics provide candidates for the nonbaryonic dark matter, candidates which existed even before the need for NBDM was established. Indeed, given that these extensions to the standard model are theoretically compelling from particle physics considerations, without cosmological considerations, their provision of dark matter energy of the right order of magnitude abundance provides an independent hint at the existence of physics beyond the standard model. Among the various candidates for this nonbaryonic dark matter (see e.g. ), a typical unified supersymmetric theory with conserved R-parity will have a stable lightest supersymmetric particle (LSP) that will probably constitute most of the cold dark matter (CDM). In addition, the non-baryonic dark energy would most likely consist of neutrino hot dark matter ($`\nu `$HDM), axionic CDM (ACDM), and the cosmological constant ($`\mathrm{\Lambda }`$). Given a Lagrangian for the unified theory (UT) and thermal equilibrium initial conditions for the fluid determining the cosmological evolution,The cosmological initial conditions can be determined by an inflationary model, which in principle can also be determined by the UT if the initial conditions for inflation are set by an even more fundamental principle or UT offers a unique inflationary history. We assume that the underlying cosmology is known when the relic abundances are finally calculated. one can calculate the dark matter content of our universe by solving the Boltzmann equations.
Many other NBDM candidates have been proposed . In this paper we will only analyze the situations for neutrinos, axions, and LSPs in some detail. Similar conclusions hold for the rest. At appropriate places in the paper we will consider non-thermal evolution.
The neutrino situation is simpler. Once the cosmology is known and the neutrino masses are determined, one can compute the relic number of neutrinos, multiply by the masses, and obtain $`\mathrm{\Omega }_\nu `$. If neutrinos have anomalous interactions they have to be included.
The axion case is difficult. Assume axions are observed in terrestrial detectors, and perhaps even information about them comes from astrophysical data. Then the thermal density $`\mathrm{\Omega }_{a(t)}`$ can be computed precisely enough . However, coherent oscillations of axion zero modes give non-thermal contributions $`\mathrm{\Omega }_{a(nt)}`$ that depend strongly on a “misalignment” factor, related to the value of the axion field at the confinement transition. We are aware of no way to determine $`\mathrm{\Omega }_{a(nt)}`$, which can be the dominant contribution, so we conclude that it may not be possible to even know how much of the CDM is $`\mathrm{\Omega }_{axion}=\mathrm{\Omega }_{a(t)}+\mathrm{\Omega }_{a(nt)}`$.
The LSP case turns out to be difficult, but solvable. We first demonstrate that knowing the LSP mass, and even knowing its cross section on a nucleon, do not allow one to know its contribution to $`\mathrm{\Omega }_{LSP}`$ to better accuracy than an order of magnitude or so. Part of this uncertainty arises from a current lack of knowledge of the phases of the soft supersymmetry breaking Lagrangian ($`_{soft}`$), but even if one arbitrarily took all phases to be zero or $`\pi `$ at least a factor of six uncertainty in $`\mathrm{\Omega }_{LSP}`$ remains. And for an important question such as the composition of the CDM, one would not want to make unwarranted assumptions about the phases or other relevant parameters. We then show that knowing the relevant parameters of $`_{soft}`$, and $`\mathrm{tan}\beta `$, to about 5% allows a determination of $`\mathrm{\Omega }_{LSP}`$ to similar accuracy. Such measurements will be possible by combining data from hadron colliders, low energy (e.g. electric dipole moments) experiments, B-factories, and a lepton collider with a polarized beam and sufficient energy to produce several superpartners (and appropriate luminosity).
If WIMPs are observed in explicit (i.e., non-collider) detection experiments (direct underground WIMP scattering on nuclei, or space-based, or indirectly via $`\nu `$ interactions from WIMP annihilation in the sun or earth) there are additional large uncertainties which we do not consider in this paper. These include going from WIMP-nucleus to WIMP-quark cross sections, local density and velocity distributions of WIMPs, how long antiprotons or positrons can persist in the galaxy, and so on. All of our results hold even if these other factors could be controlled.
In the next two sections we discuss the LSP case in detail, first analytically and then numerically. Then we turn to some cosmological considerations, and examine the axion case. In all cases we assume that $`h^2`$ will be known accurately before it is needed to achieve a particle physics knowledge of $`\mathrm{\Omega }_{NBDM}`$, so we do not include errors in $`h^2`$. We also assume that the usual Boltzmann treatment is sufficient to estimate the possible uncertainties and that refinements of the calculational procedure necessary for more accurate calculation can be accomplished when appropriate. For examples of such refinements, see Ref. and references therein.
## II Analytic estimates
In this section, we give an analytic estimate for the uncertainties that can be expected in the calculation of LSP CDM, given that some of the parameters in the MSSM has been measured to a certain accuracy. The most common situation for most parts of the MSSM parameter space is when no other particle mass is within $`m_\chi /20`$ of the LSP mass $`m_\chi `$ . In that case, self-annihilation determines the relic abundance. Assuming as usual that the annihilation products are in chemical thermal equilibrium with the massless degrees of freedom in a radiation dominated flat FRW universe, the simplified Boltzmann equation governing the relic abundance can be written as
$$\frac{df}{dx}=m_\chi \sqrt{\frac{45M_{pl}^2}{4\pi ^3g_{}}}\sigma v(f^2f_0^2)$$
(1)
where $`f=n_\chi /T^3`$ is the LSP volume density scaled by the cube of the temperature which sets the scale for the volume density of the photons, $`\sigma v`$ is the thermal averaged self annihilation cross section, $`x=T/m_\chi `$ is the scaled temperature, $`g_{}`$ counts the degrees of freedom contributing to the entropy, and $`f_0=\frac{x^{3/2}}{\sqrt{2\pi ^2}}e^{1/x}`$ is the nonrelativistic approximation of the thermal equilibrium volume density of LSPs. Starting from a thermal equilibrium initial conditions, $`f`$ tracks $`f_0`$ until the freeze-out temeperature $`x_F`$ and then $`f`$ decouples from $`f_0`$. Approximating this decoupling to occur sharply at a particular temperature $`x_F`$, one finds an expression for the relic density as
$$\mathrm{\Omega }=T_X^3\sqrt{\frac{4\pi ^3g_{}}{45M_{pl}^2}}(_{x_0}^{x_F}\sigma v𝑑x)^1/\rho _c$$
(2)
where $`\rho _c`$ is the critical energy density that only depends on the cosmology. The approximate expression for the freezeout temperature can be written as
$`x_F{\displaystyle \frac{1}{\mathrm{ln}[m_\chi \xi \sigma v]+\frac{1}{2}\mathrm{ln}x_F}}`$
$`\xi {\displaystyle \frac{1}{(2\pi )^3}}\sqrt{{\displaystyle \frac{45M_{pl}^2}{2g_{}}}}`$Propagating the error in quadratures, the fractional error is characterized by
$`\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Omega }}{\mathrm{\Omega }}}\right)^2={\displaystyle \underset{i}{}}(\mathrm{\Delta }_i(\delta P_i))^2`$where $`\delta P_i`$ denotes the uncertainty in parameter $`P_i`$. Neglecting the $`T`$ and $`g_{}`$uncertainties, we can write
$`\mathrm{\Delta }_i={\displaystyle \frac{[x_F^2\frac{m_\chi }{P_i}\frac{\sqrt{g_{}}}{m_\chi }\sigma v+x_F^2\sqrt{g_{}}\frac{\sigma v}{P_i}\frac{\sigma v}{P_i}\sqrt{g_{}}𝑑x]}{_{x_0}^{x_F}\sqrt{g_{}}\sigma v𝑑x}}`$where we have used the fact that $`x_F1`$. We can approximate $`\frac{\sigma v}{P_i}\sigma v\frac{r_i}{P_i}`$where $`r_iO(1)`$ (which is a good approximation in the region of the parameter space where the $`P_i`$ dependence is analytic) and conclude
$`\mathrm{\Delta }_i{\displaystyle \frac{r_i}{P_i}}.`$Hence, we conclude that
$`\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Omega }}{\mathrm{\Omega }}}\right)^2={\displaystyle \underset{i}{}}r_i^2({\displaystyle \frac{\delta P_i}{P_i}})^2`$where the strength of the errror contribution is determined by $`r_i=\frac{\mathrm{ln}\sigma v}{\mathrm{ln}P_i}`$which is what we expect.
For example, consider a typical nonresonant self-annihilation cross section of an LSP going to two fermions through a sfermion exchange. We have
$`\sigma v{\displaystyle \frac{g_{f\stackrel{~}{f}\chi }^2}{64\pi }}\sqrt{1{\displaystyle \frac{m_f^2}{m_\chi ^2}}}\left[{\displaystyle \frac{c_1m_f^2+c_2xm_\chi ^2}{(m_{\stackrel{~}{f}}^2+m_\chi ^2m_f^2)^2}}\right]`$where $`c_1`$ and $`c_2`$ are $`O(1)`$ dimensionless constants and $`g_{f\stackrel{~}{f}\chi }^2`$ is the fermion-sfermion-neutralino coupling. Then considering the contribution of $`m_\chi `$ to the fractional error, we find to leading order in $`\frac{m_f}{m_\chi }\frac{m_f}{m_{\stackrel{~}{f}}}`$ ,
$`r_{m_\chi }={\displaystyle \frac{2\left(1\left(\frac{m_{\stackrel{~}{f}}}{m_\chi }\right)^2\right)}{1+\left(\frac{m_{\stackrel{~}{f}}}{m_\chi }\right)^2}}`$which is about $`\stackrel{<}{}1`$ . Considering the contribution of $`m_{\stackrel{~}{f}}`$ to the uncertainty, we find similarly
$`r_{m_{\stackrel{~}{f}}}={\displaystyle \frac{4}{1+\left(\frac{m_\chi }{m_{\stackrel{~}{f}}}\right)^2}}`$which has a magnitude of about $`\stackrel{>}{}2`$. Finally we see that the coupling will also contribute
$`r_{g_{f\stackrel{~}{f}\chi }}=2,`$showing that the abundance is most sensitive to the sfermion mass and the LSP-sfermion-fermion coupling. In addition we see that relative uncertainties in parameters are likely to lead to even larger relative uncertainties in $`\mathrm{\Omega }_{LSP}`$.
The standard scalar neutralino cross section on protons is expressed as
$$\sigma _p^{scalar}=\frac{4m_r^2}{\pi }f_p^2$$
(3)
where the effective scalar interaction coupling $`f_p`$ is usually dominated by the CP-even Higgs parton level exchange between quarks and neutralinos
$$f_p^{scalar}f_p^{(H)}=m_p\left[\underset{q=u,d,s}{}f_{T_q}^p\frac{f_q^H}{m_q}+\frac{2}{27}f_{TG}^p\underset{q=c,b,t}{}\frac{f_q^H}{m_q}\right].$$
(4)
The matrix element coefficients $`f_{T_q}^p`$ and $`f_{TG}^p=1_qf_{T_q}^p`$ can be extracted from pion-nucleon scattering using chiral perturbation theory and are subject to large uncertainties which are reflected in the neutralino-proton scattering cross section .
The parton level couplings, on the other hand, depend entirely on the SUSY Lagrangian parameters and Higgs masses
$`f_q^H=m_q{\displaystyle \underset{i=1,2}{}}{\displaystyle \frac{c_\chi ^{(i)}c_q^{(i)}}{m_{H_i}^2}}`$with the (in general complex) neutralino couplings to the light CP-even Higgs
$`c_\chi ^{(1)}={\displaystyle \frac{1}{2}}(gN_{12}^{}g^{}N_{11}^{})(N_{13}^{}\mathrm{sin}\alpha +N_{14}^{}\mathrm{cos}\alpha )`$and to the heavy Higgs
$`c_\chi ^{(2)}={\displaystyle \frac{1}{2}}(gN_{12}^{}g^{}N_{11}^{})(N_{14}^{}\mathrm{sin}\alpha N_{13}^{}\mathrm{cos}\alpha ),`$where $`\alpha `$ is the Higgs mixing angle. The quark-Higgs couplings depend on weak isospin quantum number and we have for the up type quarks
$`c_u^{(1)}={\displaystyle \frac{g}{2m_W}}{\displaystyle \frac{\mathrm{cos}\alpha }{\mathrm{sin}\beta }},`$ $`c_u^{(2)}={\displaystyle \frac{g}{2m_W}}{\displaystyle \frac{\mathrm{sin}\alpha }{\mathrm{sin}\beta }}`$ (5)
and for the down type quarks
$`c_d^{(1)}={\displaystyle \frac{g}{2m_W}}{\displaystyle \frac{\mathrm{sin}\alpha }{\mathrm{cos}\beta }},`$ $`c_d^{(2)}={\displaystyle \frac{g}{2m_W}}{\displaystyle \frac{\mathrm{cos}\alpha }{\mathrm{cos}\beta }}.`$ (6)
In the large $`m_A`$ limit $`\alpha \beta \pi /2`$ and the $`c_d^{(2)}`$ coupling is enhanced by a factor of $`\mathrm{tan}\beta `$. As a result, for small values of $`\mathrm{tan}\beta `$ (less than 4) the cross section is dominated by the light Higgs exchange while for large values of $`\mathrm{tan}\beta `$ the heavy Higgs exchange contribution prevails.
In analogy to the analysis of the relic density we can determine accuracy of the cross section calculation as a function of the variation in individual parameters
$`\left({\displaystyle \frac{\mathrm{\Delta }\sigma _p}{\sigma _p}}\right)^2={\displaystyle \underset{i}{}}s_i^2({\displaystyle \frac{\mathrm{\Delta }P_i}{P_i}})^2`$where $`s_i`$ is the variation coeficient corresponding to the uncertainty $`\delta P_i`$ in parameter $`P_i`$. Assuming that there is indeed a dominant contribution from one of the Higgs bosons, we find that the Higgs and quark coupling both contribute to the uncertainty in a similar way
$`s_{c_q}s_{c_\chi }2`$ (7)
and the Higgs mass contribution to the uncertainty is
$`s_{m_H}4.`$ (8)
while the cross section is largely insensitive to variations of the neutralino mass since the reduced mass of the system is very close to the mass of the proton.
It is important to realize that the neutralino-Higgs couplings depend crucially on the neutralino mixing matrix. Both gaugino and Higgsino components are required to participate if the couplings are not to vanish. The neutralino mass matrix depends on complex quantities $`M_1`$, $`M_2`$ and $`\mu `$ with potentially large phases which can significantly modify the lightest neutralino composition and subsequently its couplings to the Higgs bosons. It is important to include these complex phases in the general analysis of the relic density and the neutralino proton cross section in order to be able to estimate the overall uncertainty in these quantities which can be achieved once the SUSY Lagrangian parameters including the phases are measured.
## Numerical Results
We illustrate the general behavior of the neutralino relic density and proton elastic cross section on a characteristic set of the MSSM parameters which can provide neutralino relic abundance in the cosmologically relevant region and a cross section small enough to be allowed by direct detection experiments. Since the neutralino scattering cross section grows with $`\mathrm{tan}\beta `$ models with low and moderate values of $`\mathrm{tan}\beta `$ can easily satisfy direct detection constraints and still provide a significant neutralino abundance. In all of our numerical calculations we use the following reference set of parameters – $`M_1=80\mathrm{GeV}`$, $`m_A=250\mathrm{GeV}`$, $`\mathrm{tan}\beta =3`$, $`\phi _1=\phi _\mu =0`$, $`m_{\stackrel{~}{\nu }}=110\mathrm{GeV}`$, $`m_{\stackrel{~}{\mathrm{}}_L}=125\mathrm{GeV}`$, $`m_{\stackrel{~}{\mathrm{}}_R}=110\mathrm{GeV}`$, $`m_{\stackrel{~}{q}_L}=420\mathrm{GeV}`$, $`m_{\stackrel{~}{u}_R}=400\mathrm{GeV}`$, $`m_{\stackrel{~}{d}_R}=380\mathrm{GeV}`$.
This set has been chosen so that the resulting neutralino is predominantly a bino and the effects of the light Higgs and Z pole neutralino annihilation pole as well as any co-annihilation effects are minimized. Our choice of parameters leads to the values of $`\mathrm{\Omega }h^20.148`$ and $`\sigma _p11.4\times 10^9pb`$. Since we are working in a general parametric framework, the soft Higgs masses can be chosen so that electroweak symmetry breaking conditions are satisfied. These numbers are not special, and only illustrate typical results.
First let us turn to the discussion of the neutralino observables sensitivity to the CP-conserving parameters appearing in the supersymmetric Lagrangian. Obviously, the most important ones are $`M_1`$, $`M_2`$, $`\mu `$ and $`\mathrm{tan}\beta `$ which enter into neutralino mass matrix and determine both the mass of the lightest neutralino $`m_\chi `$ and its Higgs and lepton-slepton couplings. Another significant variation comes from the change in the mass of the CP-odd Higgs boson which influences the neutralino-Higgs couplings and the heavy CP-even Higgs boson mass. The dominant slepton exchange contribution to the neutralino annihilation cross section depends crucially on the slepton mass while the scattering cross section is insensitive to the scalar masses.
Fig. 1 shows the dependence of the neutralino relic density and proton cross section on the relative variation of the important parameters within a $`\pm 20\%`$ range. As shown in frame a), the relic density is mostly sensitive to variations in the slepton masses appearing in the dominant annihilation diagram – as the slepton mass increases the annihilation cross section is more suppressed and $`\mathrm{\Omega }h^2`$ increases. All the other parameters enter the neutralino sector and their variations are reflected in the variations of the LSP mass and the LSP-slepton-lepton coupling. The most significant is the variation of the density with the bino mass $`M_1`$ which determines $`m_\chi `$ and subsequently the dominant $`s`$-wave contribution to the neutralino annihilation cross section. It is relatively stable with respect to changes in $`M_2`$, $`\mu `$ and $`\mathrm{tan}\beta `$ since the annihilation cross section depends weakly on the neutralino composition. In summary, the overall variation of $`\mathrm{\Omega }h^2`$ is at most $`\pm 25\%`$ for any single parameter in the given range of the input SUSY parameters.
On the other hand, the spin independent part of the neutralino-proton scattering cross section depends crucially on the gaugino-Higgsino composition of the lightest neutralino since both components take part in the neutralino-Higgs interaction. This fact is reflected in frame b) which shows a comparable sensitivity to $`M_1`$, $`M_2`$ and $`\mu `$. It is obvious that for the given range of $`\delta =\mathrm{\Delta }c/c`$ the change in $`\sigma _p`$ can be as big as a factor of two. Note that in our particular case ($`\mathrm{tan}\beta =3`$) the cross section is dominated by the light Higgs boson exchange and consequently the sensitivity to $`m_A`$ is limited. As $`\mathrm{tan}\beta `$ increases, the heavy Higgs exchange takes over the cross section and the sensitivity to $`m_A`$, which determines the heavy Higgs boson mass, increases correspondingly.
In order to estimate the accuracy with which we can determine the relic density and the proton scattering cross section once we have measured the SUSY parameters at a collider experiment we plot the range of both quantities when the SUSY parameters are all varied within a 20 % range from the central value. Fig. 2 shows the resulting region in the $`\mathrm{\Omega }h^2`$$`\sigma _p`$ plane. In part of parameter space the variations of Fig. 1 can combine, and lead to much larger ranges of variation in $`\mathrm{\Omega }h^2`$. In particular, the points in the upper left region with small relic density result from non-linear effects associated with the presence of an s-channel light Higgs exchange annihilation pole, which can drastically increase the annihilation cross section, and parameter points on the tail of the Breit-Wigner resonance have small neutralino relic density. This is a fairly generic situation since the allowed range for neutralino relic abundance imposes a limit on the neutralino mass. Note that the 5% error region allows a reasonable determination of $`\mathrm{\Omega }h^2`$. To understand the weak constraints from direct detection experiments, imagine a horizontal line across Fig.2, with a height uncertainty coming from the nuclear physics and astrophysics ambiguities in extracting $`\sigma _p`$, and notice the resulting uncertainties in $`\mathrm{\Omega }h^2`$.
Next we examine the effects of soft SUSY phases on these questions. It has been demonstrated recently that all soft SUSY phases can be large. Even if it turned out that some are smaller, e.g. from future electric dipole moment experiments, others relevant here may not be. Also, phases have two effects – not only do they directly enter the calculations of $`\mathrm{\Omega }h^2`$ and $`\sigma _p`$, they also make measurement of other parameters such as $`\mathrm{tan}\beta `$ more difficult, and impossible at LEP or hadron colliders .
We work in a phase parametrization consistent with and choose $`\phi _2=Arg(M_2)=0`$. In Fig. 3 we show the range of $`\mathrm{\Omega }h^2`$ (frame a)) and $`\sigma _p`$ frame (b)) as the two relevant phases entering the neutralino mass matrix $`\phi _1`$ and $`\phi _\mu `$ are varied in their full range while all other parameters are kept constant. It turns out that the relic abundance has two local minima and maxima and the range is bigger than an order of magnitude between the minimum and maximum. The neutralino proton cross section also varies by more than an order of magnitude but it is monotonous in the $`\phi _\mu `$\- $`\phi _1`$ plane. From Fig. 3 it is clear that by neglecting the phases one can be missing a crucial part of the particle physics information needed to determine the LSP abundance.
## Cosmological uncertainties and axions
Although the focus of our paper is examining the uncertainties in our knowledge of dark matter coming from particle physics, here we will briefly comment on some of the cosmological uncertainties related to the dark matter determination. Specifically, we would like to emphasize a point that can easily be overlooked: only with a firm determination of the cosmological history can the particle physics data weigh the universe. To make this point concrete, let us consider some possible difficulties with the usual approach to the dark matter calculation which we have assumed in our paper. First we will consider the effect of energy density behaving like a cosmological constant today. Then, we will consider the situation when the dominant contribution to zero pressure dark energy is not a thermal relic. In such “nonthermal” cases, particle physics data generically cannot determine the “matter” energy density contribution because “history” or boundary conditions of the particle physics may not be determined by a fundamental principle or dynamics. We will choose axionic dark matter to illustrate this point. Such uncertainties are relevant to all forms of dark matter. There can also be extra entropy production which has to be included.
Let us first consider how the existence of energy density with a negative pressure (which we will for simplicity call $`\rho _\mathrm{\Lambda }`$ ) affects the relic LSP density calculations. First, consider the Boltzmann equations. Suppose the spatial curvature is negligible as will be the case after inflation. Then, we can start with the same Boltzmann equation
$`{\displaystyle \frac{dn}{dt}}+3Hn=\sigma v(n^2n_0^2)`$Since conservation of entropy is still valid, we will have
$`{\displaystyle \frac{df}{dt}}=\sigma vT^3(f^2f_0^2)`$where $`f\frac{n}{T^3}.`$ In the usual scenario, we assume radiation dominance and use
$`t={\displaystyle \frac{1}{2H}}`$This is still a good approximation because the temperature at which the relic density froze out is $`O(\text{GeV})T_{\text{nucleosynthesis}}`$ and we require that $`\rho _\mathrm{\Lambda }`$ not destroy one of the pillars of standard cosmology and that $`\mathrm{\Omega }_\mathrm{\Lambda }`$ be a monotonic function of time. Hence, the rest of the formalism follows as usual for the determination of the freezeout temperature. As far as the evolution of the relic density after freezeout is concerned, there is no change from before again because the nucleosynthesis temperature is much lower than the freezeout temperature.
Of course, if we do not require monotonic variation of $`\rho _\mathrm{\Lambda }`$ then the above relation between the time and expansion rate may no longer valid and the relic abundance calculation may have to be modified. However, this is unlikely because to preserve the successes of BBN, the models of $`\rho _\mathrm{\Lambda }`$ then must be fine tuned such that its energy is significant at the time of dark matter freezeout while being negligible at the time of BBN . In the unlikely event that such tuned history is that of our universe, then the dark matter energy density implied by the particle physics measurements may be significantly different from what we have calculated in this paper.
Axion scenarios offer an elegant solution to the strong CP problem. Furthermore, pseudoscalars resembling axions are generically expected from effective theories arising from compactifications of string theories.However, most embeddings of axions in string compactification models run into trouble with scales although there may be some ways to circumvent the problem. Unlike LSPs, the axions are expected to naturally contribute significantly to the cosmological energy density even when the thermal relic abundance does not contribute significantly. Even with the assumption of inflation and sufficiently low reheating temperature rendering possible axionic topological defect related contributions irrelevant, one must account for the fact that generically there may be coherent oscillations of the axion zero modes (commonly called misalignment contribution) giving a nonthermal contribution to the cosmological energy density with the magnitude
$`\mathrm{\Omega }_{a(nt)}h^2=0.13\times 10^{\pm 0.4}({\displaystyle \frac{\mathrm{\Lambda }_{QCD}}{200\mathrm{MeV}}})^{0.7}f(\theta _1)\theta _1^2({\displaystyle \frac{m_a}{10^5\mathrm{eV}}})^{1.18}`$where $`m_a`$ is the axion mass, $`f`$ is some known monotonically increasing function accounting for anharmonic effects ($`f(0)=1)`$, and $`\theta _1`$ is the “amplitude” of the oscillations of the axion field. The value of $`\theta _1`$ is essentially equal to the value of the axion field at the confinement transition. Because there is no direct way to measure $`\theta _1`$, even with a precise determination of $`m_a`$, there will be a great uncertainty in the axionic contribution. This is in contrast with the most likely smaller thermal relic contribution of the axion which can be quite precisely determined once the mass of the axion is measured. Indeed, most of the relic axion detection experiments such as electromagnetic spectrum observation of nearby clusters, as well as Sikivie-type resonant microwave cavities immersed in strong magnetic fields, mostly measure the axion mass, and there is no obvious direct handle on the “misalignment” angle $`\theta _1`$. Indeed, $`\theta _1`$ may be most likely determined randomly from the space of all possible values, since above the confinement transition the axions have essentially zero mass. Hence, in such “nonthermal” cosmological boundary condition dependent situations, the determination of the dark matter from particle physics data seems impossible.
## Conclusion
Understanding the composition of our universe is one of the most fundamental issues we can study. We demonstrate that we can only be fully confident we have learned the answer if we can calculate from first principles and laboratory data for each type of matter X what $`\mathrm{\Omega }_X`$ is, and find agreement with cosmological measurement of $`\mathrm{\Omega }_{matter}`$. For neutrinos this can be done if their masses can be measured, or shown to be small compared to $`1\mathrm{eV}`$. For axions we have emphasized there does not seem to be a way to do this.
For the LSP it is possible to calculate $`\mathrm{\Omega }_{LSP}`$ if measurements of soft supersymmetry phases, and $`\mathrm{tan}\beta `$ are made. Unless such measurements are done, it is simply incorrect to suggest that detection of the LSP at a collider, or in a direct or indirect WIMP experiment, means that the cold dark matter has been observed. Indeed, there is a loose correlation of easier detection with smaller relic density. All measurements, of course, help constrain the situation, and the detection of the LSP will mean that some of the CDM has been observed – but perhaps only a fraction.
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# Measuring cosmological bulk flows via the kinematic Sunyaev-Zeldovich effect in the upcoming cosmic microwave background maps.
## Abstract
We propose a new method to measure the possible large-scale bulk flows in the Universe from the cosmic microwave background (CMB) maps from the upcoming missions, MAP and Planck. This can be done by studying the statistical properties of the CMB temperature field at many X-ray cluster positions. At each cluster position, the CMB temperature fluctuation will be a combination of the Sunyaev-Zeldovich (SZ) kinematic and thermal components, the cosmological fluctuations and the instrument noise term. When averaged over many such clusters the last three will integrate down, whereas the first one will be dominated by a possible bulk flow component. In particular, we propose to use all-sky X-ray cluster catalogs that should (or could) be available soon from X-ray satellites, and then to evaluate the dipole component of the CMB field at the cluster positions. We show that for the MAP and Planck mission parameters the dominant contributions to the dipole will be from the terms due to the SZ kinematic effect produced by the bulk flow (the signal we seek) and the instrument noise (the noise in our signal). Computing then the expected signal-to-noise ratio for such measurement, we get that at the 95% confidence level the bulk flows on scales $`100h^1`$Mpc can be probed down to the amplitude of $`<200`$ km/sec with the MAP data and down to only $`30`$ km/sec with the Planck mission.
Cosmology: Cosmic Microwave Background. Large Scale Structure. Galaxies: Clusters.
1. Introduction.
Peculiar motions trace the overall mass distribution and it is important to determine the coherence scale and amplitude of bulk flows. Current measurements range from bulk flows as high as 700 km/sec on scale of $``$150$`h^1`$Mpc (Lauer & Postman 1994) to finding little peculiar motion on scales beyond $``$70$`h^1`$Mpc (see Willick 2000 for review). It is thus important to find alternative ways to test for such bulk flows. One alternative way to measure peculiar flows is via the kinematic Sunyaev Zeldovich (SZ) effect produced on the cosmic microwave background (CMB) photons from the hot X-ray emitting gas in clusters of galaxies (Zeldovich & Sunyaev 1969). Such program is already being undertaken in the SuZIE project which plans to measure motion of 40 clusters at $`z`$=0.1-0.3 and determine the peculiar velocity of each cluster to a precision of 700 km/sec.
In this letter we propose to use CMB data from the MAP and Planck missions to measure the bulk flows in a quick, cheap and efficient way. X-ray cluster catalogs to be available shortly based on ROSAT, ASCA and XMM measurements, will provide locations of the SZ sources on the CMB sky and a reasonable estimate of the cluster electron density distribution and temperature. If there are significant bulk motions they will leave an imprint via the cumulative kinematic SZ effect and can be uncovered by cross-correlating the temperature field at the cluster positions in the MAP and Planck CMB maps. Such bulk motions would produce a significant dipole component in the temperature field evaluated at the cluster locations. By averaging temperature fields at enough cluster positions, the thermal SZ and other noise components will integrate down enough to reveal possible bulk motions out to $``$ 300-400 km/s on scales $``$50-100$`h^1`$Mpc.
2. Dipole of the cumulative SZ kinematic effect.
Consider the CMB field at a beam centered on one isothermal X-ray emitting cluster at the angular position $`\stackrel{}{y}`$. If the cluster is moving with the line-of-sight velocity $`v_r`$ with respect to the CMB rest frame, the SZ CMB fluctuation at frequency $`\nu `$ at this position will be $`\delta _\nu (\stackrel{}{y})`$=$`\delta _{th}(\stackrel{}{y})G(\nu )`$+$`\delta _{kin}(\stackrel{}{y})H(\nu )`$, with $`\delta _{th}`$=$`\tau T_{\mathrm{virial}}/T_{e,ann}`$ and $`\delta _{kin}`$=$`\tau v_r/c`$ (e.g. Phillips 1995). Here $`\tau `$ is the projected optical depth due to Compton scattering, $`T_{\mathrm{virial}}`$ is the cluster virial temperature and $`k_BT_{e,ann}`$=0.5 MeV. (For expressions for the spectral dependence of the two SZ components, $`G(\nu ),H(\nu )`$ see e.g. Birkinshaw (1999).) Normally the thermal term dominates for individual clusters, but if averaged over many clusters moving at a significant bulk flow with respect to the CMB rest-frame, the former will integrate down $`1/N_{\mathrm{cluster}}`$, while the kinematic term will reflect coherent motions with amplitude $`V_{\mathrm{bulk}}`$.
In order to minimize the contribution from other sources, we will start with CMB maps from which the cosmological dipole component was subtracted down to $`\sigma _d`$. This uncertainty is already $`\sigma _d=7\mu `$K (Fixsen et al. 1996) and should be significantly smaller in MAP observations. The MAP radiometers produce a raw temperature measurement that is the difference between two points on the sky $`140^o`$ apart which introduces an error on the dipole determination, due to correlated noise, of the order of $`0.1\mu `$K; the dominant uncertainty in the dipole will be due to confusion from the galactic foreground (G. Hinshaw, private communication). After the cosmological CMB dipole subtraction the CMB fluctuation in band $`\nu `$ at position $`\stackrel{}{y}`$ centered on a known X-ray cluster will be $`\delta _\nu (\stackrel{}{y})`$=$`\delta _{th}(\stackrel{}{y})G(\nu )`$+$`[\delta _{kin}(\stackrel{}{y})`$+$`\delta _{\mathrm{CMB}}(\stackrel{}{y})]H(\nu )`$+$`r(\nu )`$. Here $`r(\nu )`$ is the instrument noise at frequency $`\nu `$ and $`\delta _{\mathrm{CMB}}(\stackrel{}{y})`$ is the cosmological CMB component whose dipole is now $`\sigma _d^2`$. Consider the dipole component of $`\delta _\nu (\stackrel{}{y})`$ with the dipole amplitude $`C_1`$ normalized so that a coherent motion at velocity $`V_{\mathrm{bulk}}`$ would lead to the dipole amplitude of $`V_{\mathrm{bulk}}^2/c^2`$. When computed from the total of $`N_{\mathrm{cluster}}`$ positions the dipole of the noise term becomes $`r^2(\nu )/N_{\mathrm{cluster}}`$. The cosmological signal gives rise to two different dipole contributions: 1) the cosmological dipole has not been perfectly removed so the temperature anisotropies at the cluster locations sample the residual dipole $`\sigma _d`$; and 2) even if all the cosmological dipole had been removed the intrinsic CMB temperature anisotropies could be seen as an extra dipole noise source. The latter contribute $`\sigma _{\mathrm{CMB}}^2/N_{\mathrm{cluster}}`$, with $`\sigma _{\mathrm{CMB}}`$ being the variance of the cosmological temperature field on the smallest angular scales probed by the experiment. Thus for $`N_{\mathrm{cluster}}1`$ the dipole of (3) becomes:
$$C_{1,\nu }C_{1,kin}H^2(\nu )+C_{1,th}G^2(\nu )+[\sigma _{\mathrm{CMB}}^2/N_{\mathrm{cluster}}+\sigma _d^2]H^2(\nu )+r^2(\nu )/N_{\mathrm{cluster}}$$
(1)
where it was assumed that for each individual X-ray cluster the thermal SZ term dominates.
3. Signal and noise terms.
We now estimate the amplitude of the signal, $`C_{1,kin}`$, and noise terms in eq. (1).
3.1. Kinematic component.
Assuming that cluster properties are independent of their velocities, this term is $`C_{1,kin}`$=$`T_0^2\tau _i^2\frac{V_{\mathrm{bulk}}^2}{c^2}`$. $`T_0`$ is the CMB temperature, and $`i`$ refers to the frequency band. The effective optical depth is affected by the beam dilution. To evaluate the expected mean optical depth accounting for the beam dilution effects we proceed as follows: for X-ray clusters the electron density profile can be approximated by the $`\beta `$-model. For isothermal and spherical X-ray gas distribution, the optical depth as a function of the angular distance from the cluster center would be $`\tau =\tau _0(1+\theta ^2D^2/r_{\mathrm{core}}^2)^{\frac{3\beta 1}{2}}`$. Here $`D`$ is the distance to the cluster and $`r_{\mathrm{core}}`$ is its core radius. For each individual cluster observed in a CMB search, this expression needs to be convolved with the beam profile. The effective optical depth becomes $`\tau _i\tau _0\mathrm{\Psi }_i(D)`$, where $`\mathrm{\Psi }_i`$ accounts for the beam dilution effects. To evaluate $`\mathrm{\Psi }_i(D)`$ we compiled a list of 37 X-ray clusters with measured $`\beta ,r_{\mathrm{core}}`$ and 2.4 KeV$`<`$$`T`$$`<`$14.6 KeV from Arnaud & Evrard (1999), Myers et al. (1997) and Neumann & Arnaud (1999). We then computed $`\mathrm{\Psi }_i(D)`$ for each individual cluster and the mean and the r.m.s. in $`\mathrm{\Psi }_i(D)`$ evaluated over the ensemble of clusters. Fig.1 plots the mean and the r.m.s. values of $`\mathrm{\Psi }_i(D)`$ vs $`D`$ for the largest and smallest MAP beams. The figure shows that for the purposes of estimating the magnitude of the kinematic dipole component we can assume that a “universal” profile for the cluster optical depth exists to within an uncertainty of $``$10-20%.
The value of $`\tau _0`$ can be estimated from the observed cluster properties and their X-ray luminosity function. Cooray (1999 - Table 1) compiled a list of 14 X-ray clusters with measured SZ thermal components, $`\mathrm{\Delta }T_{\mathrm{SZ}}`$, with the X-ray luminosity in the \[2-10\] KeV range, $`L_X`$(2-10 KeV), between 1.6$`\times `$$`10^{43}`$ and 3.6$`\times `$$`10^{45}h^2`$erg/sec. A linear regression fit to the data gives $`\tau _0`$=(4.8$`\pm `$1.0)$`\times `$$`10^3[L_X`$(2-10KeV)$`/10^{44}h^2\mathrm{erg}/\mathrm{sec}]^{\alpha _\tau }`$ with $`\alpha _\tau `$=0.41$`\pm `$0.12. Note that for isothermal X-ray clusters emitting due to thermal Bremsstrahlung, $`L_Xn_e^2T_{\mathrm{virial}}^{1/2}`$, and obeying the observed X-ray luminosity - temperature relation, $`L_X`$$``$$`T_{\mathrm{virial}}^\gamma `$ with $`\gamma `$$``$2.5 (Mushotzky & Scharf 1997, Allen & Fabian 1998, Arnaud & Evrard 1998), one expects $`\tau `$$``$$`L_X^{0.4}`$ (e.g. Haenhelt & Tegmark 1996). These relations show little evolution out to $`z`$$``$0.5 (Mushotzky & Scharf 1997; Schindler 1999) and are valid for cluster catalogs of depth $`<`$200$`h^1`$Mpc needed for this project.
The mean optical depth can now be computed from measurements of the X-ray luminosity function (XLF). XLF has now been determined very accurately from the ROSAT BCS sample out to $`z`$$``$0.3 (Ebeling et al., 1997). The sample is 90% complete for fluxes $``$4.45$`\times `$$`10^{12}`$ erg/cm<sup>2</sup>/sec in \[0.1-2.4\] KeV band. The XLF is of the Schechter type $`n(L_X)dL_X`$=$`n_{}(L_X/L_{})^{\alpha _X}\mathrm{exp}(L_X/L_{})dL_X`$ with $`\alpha _X`$$``$$`1.8`$ and bolometric $`L_{}`$$``$$`9.3`$$`\times `$$`10^{44}h^2`$erg/sec or $`L_{}`$(2-10 KeV)$``$$`3.2`$$`\times `$$`10^{44}h^2`$erg/sec. For these XLF parameters and $`V_{\mathrm{bulk}}`$=600 km/sec and a constant lower limit on the absolute X-ray luminosity $`L_0`$$``$$`5`$$`\times `$$`10^3L_{}`$, corresponding to the \[0.1-2.4\]KeV flux $`4.5`$$`\times 10^{12}`$ erg/cm<sup>2</sup>/sec at 50$`h^1`$Mpc, we get $`\sqrt{C_{1,th}}`$$``$$`9\mu `$K dropping to $``$$`6\mu `$K for $`L_0`$$``$$`2`$$`\times `$$`10^3L_{}`$.
3.2. Thermal component
The residual dipole from the SZ thermal component comes from the finite number of Poisson-distributed X-ray clusters and would integrate $`1/N_{\mathrm{cluster}}`$. Since cluster properties are independent of position, the dipole contribution is $`C_{1,th}`$$``$$`(\mathrm{\Delta }T_{SZ})^2𝒟_{\mathrm{cluster}}^2`$. Here $`(\mathrm{\Delta }T_{SZ})`$ is the mean amplitude of the thermal SZ temperature fluctuation produced by the observed X-ray clusters and $`𝒟_{\mathrm{cluster}}`$=$`\mathrm{cos}\theta `$, with $`\theta `$ being the cluster azimuthal angle, is the mean dipole of the cluster distribution out to the depth on which the bulk motions are probed. Assuming the scaling of $`\tau `$ vs $`L_X`$ above and integrating it over $`L_X`$$``$$`L_0`$=const gives $`(\mathrm{\Delta }T_{SZ})`$=12$`\mu `$K for $`L_0`$=$`0.002L_{}`$ and $`(\mathrm{\Delta }T_{SZ})`$=20$`\mu `$K for $`L_0`$=$`0.005L_{}`$. We used the Abell/ACO catalog (Abell 1958, Abell, Corwin & Olowin 1989) to verify that the angular distribution of clusters has $`𝒟_{\mathrm{cluster}}^2`$$``$$`N_{\mathrm{cluster}}^1`$. Statistically, one expects $`𝒟_{\mathrm{cluster}}^2N_{\mathrm{cluster}}`$=1/3. For the ACO catalogue out to $`200h^1`$Mpc we get $`𝒟_{\mathrm{cluster}}^2N_{\mathrm{cluster}}`$$``$0.2-0.3. The numbers are consistent with related dipole and monopole parameters of the gravitational field due to clusters ($`D_S,M_S`$=$`r^2\mathrm{cos}\theta ,r^2`$) from Fig.1 of Scaramella et al. (1991), when corrected for the fact that $`D_S,M_S`$ are dominated by more nearby clusters. The different spectral dependence, $`G(\nu )`$, of thermal SZ can be used to reduce this term further. Thus the contribution of this term to the dipole noise term in eq. (1) would be $`\sqrt{C_{1,th}}<10N_{\mathrm{cluster}}^{1/2}\mu `$K. (The thermal component dipole can also be seen to be small from extrapolation to $`l`$=1 from Figs. 1,2 in theoretical computation of the thermal SZ power spectrum by Atrio-Barandela & Mücket 1999, or Fig. 8 in Refreiger et al. 2000).
3.3. Dipole cosmological components.
There will be two independent contributions to the dipole noise from cosmological terms: from the residual dipole uncertainty and from the CMB fluctuations leaving a residual dipole when evaluated over a finite number of points on the sky $`(N_{\mathrm{cluster}})`$.
The first contribution will come from the cosmological dipole which can be eliminated from the CMB maps down to the 68% uncertainty of $`\sigma _d`$ (=7$`\mu `$K for FIRAS). With $`N_{\mathrm{cluster}}`$ we will approach this systematic component as $`\sigma _d^2[1+\mathrm{O}(N_{\mathrm{cluster}}^1)]`$. Because the correlation angle of the temperature anisotropies is larger than the pixel size, this component can be further decreased using a low-pass filter by performing the analysis on $`\delta (\stackrel{}{y}_{\mathrm{cluster}})`$$`\delta (\stackrel{}{y}_{\mathrm{ngb}})`$, where $`\stackrel{}{y}_{\mathrm{ngb}}`$ is the neighboring pixel that does not contain another cluster and is $`\mathrm{\Delta }\theta _{\mathrm{ngb}}`$ away from the original pixel. The dipole noise will then be reduced to $`\sigma _d(\mathrm{\Delta }\theta _{\mathrm{ngb}}/180^\mathrm{o})`$, while the noise variance will increase by only $`1/N_{\mathrm{ngb}}`$. The second contribution will come from the cosmological temperature anisotropy at each cluster location. The r.m.s. temperature anisotropy $`\sigma _{\mathrm{CMB}}`$ on the smallest scales probed by MAP is model dependent, but could be of the same order of magnitude as the pixel noise or even larger and it has the same frequency dependence as the kinematic SZ effect. Its contribution is $`\sigma _{\mathrm{CMB}}^2/N_{\mathrm{cluster}}`$, and can be reduced by the low-pass filtering discussed above. If necessary, $`\sigma _{\mathrm{CMB}}`$ can be reduced further with Wiener filtering (cf. Haenhelt & Tegmark 1996) designed to minimize the difference between the filtered signal, $`\stackrel{~}{\delta }_{\mathrm{CMB}}`$, and the instrument noise, $`r`$. We computed the residual (Wiener-filtered) variance for two cold-dark-matter (CDM) models, the standard CDM with $`\mathrm{\Omega }`$=1, and the cosmological constant dominated CDM with $`\mathrm{\Omega }`$=0.3 with $`\sigma _{\mathrm{CMB}}`$=(57–93)$`\mu `$K in the MAP bands. After the filtering with the noise multipoles of MAP the numbers reduce to $`<`$$`\frac{1}{2}`$$`r^2^{1/2}`$; this component adds in quadrature to the instrument noise.
3.4. Instrument noise component.
If the instrument noise in the given band is $`r(\nu )`$, its contribution to the dipole term in (1) is $`r^2(\nu )/N_{\mathrm{cluster}}^1`$. The MAP mission (http://map.gsfc.nasa.gov) has 5 bands from 22 to 90 GHz and the r.m.s. noise at the end of two years should reach $`r^2(\nu )^{1/2}`$=35$`\mu `$K per 0.3$`\times `$0.3 deg pixel (the MAP beams range from 0.93 to 0.2 deg and its total lifetime would be at least 27 months). The PLANCK satellite (http://astro.estec.esa.nl/Planck) has Low and High Frequency Instruments (LFI and HFI). The LFI has four bands from 30 to 100 GHz and the beams from 33 to 10 arcmin. Its noise at 4 $`\mu `$K at 30 GHZ to 12 $`\mu `$K at 100 GHz is significantly lower than that of MAP. The six HFI bands cover 100 to 857 GHz and adding them would increase the signal-to-noise of the measurement of bulk flows with Planck. Because the instrument noise dipole would be added in quadrature to the cosmological component, this term is expected to be the dominant dipole noise term for MAP instruments.
4. Results and strategies for measurement.
In order to apply this method we will take the available all-sky catalogs of imaged X-ray clusters and compute the dipole of the CMB temperature field at the cluster locations in the expectation that the noise terms will integrate down uncovering the bulk motion contribution to the dipole.
To measure the large-scale bulk flows by this method we need information on the location of $``$(100-300) X-ray clusters with reasonably measured $`\tau `$. A new all-sky catalog of imaged X-ray clusters should be available soon (Böhringer et al. 1999) from the ROSAT observations. The current catalog, known as BCS (Bright Cluster Sample), is $``$90% complete in the Northern hemisphere and $`|b_{\mathrm{Gal}}|`$$``$20 deg and contains 199 clusters with flux $``$4.45$`\times `$$`10^{12}`$erg/cm<sup>2</sup>/sec (Ebeling et al. 1997). With the XLF parameters from Ebeling et al. (1997) a sphere of radius 100$`h^1`$Mpc would contain $``$400 X-ray clusters with $`L_{\mathrm{X},\mathrm{bol}}`$$``$$`10^{42}h^2`$erg/sec (or $`10^3L_{}`$), $``$200 clusters with $`L_X`$$``$$`2.5`$$`\times `$$`10^3L_{}`$ and $``$100 clusters with $`L_X`$$``$5$`\times `$$`10^3L_{}`$. For flux-limited X-ray catalogs the numbers are similar: for flux limit of $`10^{12}`$erg/cm<sup>2</sup>/sec in the \[0.1-2.4\] KeV band of ROSAT, which is the flux limit of the NORA and REFLEX catalogs (Guzzo 2000), the number of clusters within 100$`h^1`$Mpc would be $``$200 or $``$120 for $`F`$$``$2$`\times `$$`10^{12}`$erg/cm<sup>2</sup>/sec.
ROSAT has already completed a Southern hemisphere catalog of X-ray clusters, ROSAT-ESO Flux-Limited X-ray cluster survey or REFLEX, (Böhringer et al. 1999, Guzzo et al. 2000). The flux limit in the \[0.1-2.4\] KeV ROSAT band of the REFLEX survey is $``$1.5$`\times `$$`10^{12}`$erg/cm<sup>2</sup>/sec corresponding to $`L_X(`$2-10KeV)$``$1.7$`\times `$$`10^{42}h^2`$erg/sec at the distance of 100$`h^1`$Mpc. (Such clusters have $`\tau `$$``$9$`\times `$$`10^4`$ and are quite numerous even at that depth avoiding the problem of significant shot-noise from the rare and very high-$`\tau `$ clusters). The Northern hemisphere catalog of the ROSAT X-ray clusters (NORA) should be completed shortly (Guzzo 2000). Altogether this would result in approximately 1,500 X-ray clusters down to the limiting flux $`F`$$``$1$`\times `$$`10^{12}`$ erg/cm<sup>2</sup>/sec in the \[0.1–2.4\] KeV band; out of these $``$300 would lie within $``$100$`h^1`$Mpc. In addition, there are already 50 X-ray clusters within $``$100$`h^1`$Mpc from ASCA searches which have both central temperature and electron density profile measured (Baumgartner et al 2000); this number is expected to more than double within the next year or two (Mushotzky, private communication). The ROSAT clusters can further be imaged very efficiently with XMM in order to obtain the necessary catalog of X-ray clusters by the time the MAP mission is completed.
Our ability to determine a bulk flow of amplitude $`V_{\mathrm{bulk}}`$ is limited by the instrumental noise term. The signal-to-noise ratio, $`\chi `$, in such a measurement out to the depth $`D`$ would be given by:
$$\chi ^2=\underset{i}{}\frac{_0^DC_{1,kin}𝑑N_{\mathrm{cluster}}}{r_i^2}=T_0^2\frac{V_{\mathrm{bulk}}^2}{c^2}N_{\mathrm{cluster}}\left[\frac{3}{D^3}\underset{i}{}\frac{_0^D\tau _0^2\mathrm{\Psi }_i^2(D)D^2𝑑D}{r_i^2}\right]$$
(2)
The integral on the right-hand-side of eq. (2) accounts for the fact that the number of clusters contributing to the reduction in the dipole from the instrument noise in a thin shell $`[D;D+dD]`$, where the beam dilutes the central optical depth by $`\mathrm{\Psi }_i(D)`$, is $`dN_{\mathrm{cluster}}`$=$`4\pi [_{L_0}^{\mathrm{}}n(L_X)𝑑L_X]D^2dD`$. For clusters selected from an absolute luminosity limited catalogs the term $`\tau _0^2`$ in the integral on the RHS of eq.(2) would be independent of $`D`$; for flux-limited catalog the dependence on $`D`$ would come via $`L_0`$$``$$`D^2`$. In computing $`\chi `$ we assumed the mean shape of $`\mathrm{\Psi }(D)`$ plotted in Fig. 1.
If we want to determine just the amplitude of the bulk flow, we need the amplitude of the dipole SZ kinematic component. If we are to determine the direction of the flow, we need to measure all three dipole components: $`a_{1,1},a_{1,0},a_{1,1}`$. Thus to determine the amplitude of the bulk velocity within a given $`D`$ at the 95% c.l. we need $`\chi ^2`$=3.84; if we want to resolve all three components of the dipole at the 95% c.l., we need $`\chi ^2`$=7.81. The error bar on the dipole components translates into an uncertainty on the measurement of the three components of the velocity flow, i.e., on the bulk flow direction. If all dipole components have been measured with the same precision, and the velocity of the flow is measured with an uncertainty $`\sigma _V`$, then the direction of the flow will be determined with an angular precision $`\delta \alpha =\sqrt{2}\sigma _V/V_{\mathrm{bulk}}`$.
Fig. 2a shows the smallest $`V_{\mathrm{bulk}}`$ that can be determined at 95% c.l. with MAP and Planck/LFI data at $`D`$=50,100,150$`h^1`$Mpc for a catalog with fixed lower bound of the bolometric X-ray luminosity, $`L_{0,\mathrm{bol}}`$. The bulk velocity at $`150h^1`$Mpc from Lauer & Postman (1994) is also shown for comparison. There is only a weak dependence in $`\chi ^2`$ on the (fixed) lower bound on the X-ray luminosity of the cluster catalog. Thus this method gives a very realistic way to measure the possible bulk flows on large scales or to significantly constrain their amplitude if the latter is small. Fig. 2b shows $`\chi ^2`$ vs the depth of the cluster catalog for the bulk flow of $`V_{\mathrm{bulk}}`$=600 km/sec for MAP (lower set of lines) and Planck/LFI parameters (upper set of lines). Since $`\chi ^2`$$``$$`V_{\mathrm{bulk}}^2`$, at the 95% c.l. the 2-year MAP data can probe in this way bulk flows with the amplitude as low as $``$250 km/sec on scale of $``$$`100h^1`$Mpc; with Planck data such bulk flows can be probed to even much lower amplitudes and scales. Concerning the direction, MAP and Planck/LFI would determine á la Lauer-Potsman bulk flow direction with a 95% c.l. error of $`41^o`$ and $`16^o`$, respectively.
Finally, because the upcoming X-ray catalogs, such as the NORA/REFLEX catalog, extend to significantly beyond $`100h^1`$Mpc this method also would probe bulk flows on much larger scales. Fig.2c shows the amplitude of the bulk flow that can be determined at 95% c.l. vs the depth of the flux-limited X-ray cluster catalog. Bulk flows $`>`$100km/sec on scale $``$300$`h^1`$Mpc can be probed with this method with MAP data; with Planck/LFI this reduces to $``$30km/sec. If MAP operates longer than 2 years, this number would decrease $``$$`t_{\mathrm{operation}}^{1/2}`$. If both LFI and HFI Planck channels are combined the bulk flows on large scales can be probed to even lower limits than with the Planck/LFI alone; e.g. on scale of $`100h^1`$Mpc the amplitude of the bulk flows can be determined at 95% c.l. down to $``$30km/sec and on scale of $``$300$`h^1`$Mpc to $``$15-20km/sec.
AK acknowledges the hospitality and visiting support of the University of Salamanca. FAB acknowledges financial support of the Junta de Castilla y León (project SA 19/00B) and Ministerio de Educación y Cultura (project PB 96-1306). We thank Chuck Bennett, Gary Hinshaw, John Mather and Richard Mushotzky for fruitful discussions on CMB and X-ray parts of this work.
FIGURE CAPTIONS
Fig. 1: The dilution factor $`\mathrm{\Psi }(D)`$ vs the depth $`D`$ for two MAP bands. Solid lines are for Gaussian beam corresponding to Band 1 (22 GHz, FWHM = 0.93 deg). Dashed lines are for Band 5 (90 GHz, FWHM = 0.21 deg). Thick lines correspond to the cluster profile averaged over 37 X-ray clusters as described in the main text. Thin lines correspond to the r.m.s. value of $`\mathrm{\Psi }(D)`$ over the same sample.
Fig. 2: (a) The amplitude of $`V_{\mathrm{bulk}}`$ that can be determined at 95% c.l. ($`\chi ^2=3.84`$) is plotted vs the lower limit on the bolometric X-ray luminosity, $`L_{0,\mathrm{bol}}`$. Thin lines correspond to MAP and thick lines to Planck/LFI data parameters. Solid lines correspond to X-ray clusters out to $`D=50,100,150h^1`$Mpc from top to bottom respectively. From top to bottom, dashed lines show the value of $`V_{\mathrm{bulk}}`$ when the three components of the dipole can be determined at the 95% c.l. at $`D=50,100,150h^1`$Mpc, respectively. Open square with error bar shows the amplitude of the bulk flow at $`150h^1`$Mpc from Lauer & Postman (1994). (b) Solid lines plot $`\chi ^2`$ from eq. (2) vs the depth of the X-ray cluster catalog, $`D`$, with a cutoff on the X-ray luminosity of $`L_0/L_{}=(1,2,5)\times 10^3`$ from top to bottom respectively. Dotted lines correspond to flux-limited X-ray cluster catalog with the \[0.1-2.4\] KeV flux $`(1,2)\times 10^{12}`$erg/cm<sup>2</sup>/sec. Lower set of lines is for MAP data and the upper set of lines is for the Planck/LFI instrument parameters. In the latter case over the range of the plot the two lines for the \[0.1-2.4\] KeV flux $`(1,2)\times 10^{12}`$erg/cm<sup>2</sup>/sec coincide with each other and with the $`L_X/L_{}10^3`$ case. Thick solid horizontal line corresponds to $`\chi ^2=3.84`$, when the amplitude of the bulk flow can be determined at 95% confidence level; thick dashed horizontal line corresponds to $`\chi ^2=7.81`$, when the three components of the dipole can be determined at 95% confidence level. (c) The amplitude of the bulk flow that can be determined at 95.4% c.l. is plotted vs the depth of the flux-limited X-ray cluster catalog. Solid lines correspond to the MAP data, thin dotted lines to Planck/LFI and thick solid lines to Planck/LFI and Planck/HFI together. Two sets of lines correspond to X-ray flux $`(1,2)\times 10^{12}`$erg/cm<sup>2</sup>/sec from top to bottom respectively.
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# 1 Introduction
## 1 Introduction
The study of non-BPS D-branes<sup>1</sup><sup>1</sup>1For some recent reviews see . has provided a better understanding of several aspects of string theory. The classification of D-branes in terms of K-theory can be physically understood by thinking of D-branes as solitonic solutions of the tachyon field of higher dimensional brane-anti-brane systems . The study of the tachyon potential using string field theory has given strong evidence to the conjecture that a brane-anti-brane system annihilates into the vacuum. Due to the presence of a single decay mode unstable non-BPS D-branes have been interpreted as sphaleron solutions of string theory and some understanding of bosonic D-branes has now also been developed. The confinement process of the unbroken gauge group boson on the world-volume of annihilating brane-anti-brane systems has been better understood .
Another aspect of recent developments has involved the study of stable non-BPS D-branes through the boundary state formalism . This was first considered in a specific context in . Over the past year the boundary state formalism has been used to construct many more examples of non-BPS D-branes . These constructions have tested the connection between D-branes and K-theory as well as testing S-dualities beyond the BPS constraints. In particular the dualities between heterotic and Type II theories and heterotic and Type I theories have been investigated. Due to the fact that certain non-BPS D-branes are the lightest states carrying a particular charge it has been possible to identify their duals and compare regions of stability as well as interaction properties in the dual theories.
The integrally (rather than torsion) charged non-BPS D-branes have so far been constructed on $`𝖹𝖹_2`$ orbifolds <sup>2</sup><sup>2</sup>2In this had been generalised somewhat to a Calabi-Yau $`𝖹𝖹_2\times 𝖹𝖹_2^{}`$ orbifold, where $`𝖹𝖹_2^{}`$ is freely acting. and have been interpreted in the blow-up as coming from certain volume minimising non-supersymmetric cycles of the manifold which in the orbifold limit shrink to one-cycles. Here we discuss a $`𝖹𝖹_2\times 𝖹𝖹_2`$ orbifold of Type IIA and IIB. In particular we study BPS and non-BPS D-branes on such an orbifold. The orbifold is a particular limit of a Calabi-Yau three-fold, and as such preserves N=2, D=4 supersymmetry. The action of the two generators $`g_1`$ and $`g_2`$ will be taken as
$`g_1(x^0,\mathrm{},x^9)`$ $`=`$ $`(x^0,x^1,x^2,x^3,x^4,x^5,x^6,x^7,x^8,x^9)`$ (1.1)
$`g_2(x^0,\mathrm{},x^9)`$ $`=`$ $`(x^0,x^1,x^2,x^3,x^4,x^5,x^6,x^7,x^8,x^9)`$ (1.2)
and we label $`g_3=g_1g_2`$. Thus both $`𝖹𝖹_2`$’s have fixed points. This orbifold has been studied previously . In particular in fractional D-branes on it were discussed and it was noted that there are two kinds of fractional D-branes, those that live at the fixed hyperplanes of one of the $`g_i`$ (we shall refer to these as singly fractional D-branes) and those that live at the fixed hyperplanes of all $`g_i`$ (we shall refer to these as totally fractional D-branes).
In this paper we use the boundary state formalism to investigate the full D-brane spectrum, BPS and non-BPS on this orbifold and study the stability regions of non-BPS $`\widehat{\text{D}}`$-branes on the orbifold. We give a boundary state description of both kinds of fractional D-branes as well as bulk D-branes wrapping the special Lagrangian three-cycle of the Calabi-Yau orbifold. We find non-BPS $`\widehat{\text{D}}`$-branes similar to those of . We also find that there are new kinds of non-BPS $`\widehat{\text{D}}`$-branes on this orbifold, and we refer to these too as truncated $`\widehat{\text{D}}`$-branes . The truncated $`\widehat{\text{D}}`$-branes break all supersymmetry and are charged under twisted R-R fields. The relevant K-theory groups are evaluated and complete agreement between the D-branes we have constructed and K-theory is found.
Some of the new $`\widehat{\text{D}}`$-branes have very unusual domains of stability and besides coupling to twisted R-R sectors they couple to certain twisted NS-NS sectors as well. We investigate the decay channels of the various $`\widehat{\text{D}}`$-branes. Besides the usual decays into brane-anti-brane pairs of BPS fractional D-branes there are other decay channels in which non-BPS $`\widehat{\text{D}}`$-branes decay into one another. There are also certain decay channels whose decay products are unknown. When the decay products are known the masses and charges of the relevant objects are the same at critical radii, indicating that the process is a marginal deformation in the conformal field theory .
The construction of boundary states describing the D-branes is given in section 2, as well as the appendices. In section 3 we compute the relevant equivariant K-theory groups and find complete agreement with section 2. In section 4 we consider the compact orbifold and we discuss in detail the minimal charge configurations of the D-branes. The stability regions of the truncated non-BPS $`\widehat{\text{D}}`$-branes are analysed in section 5. The new kind of truncated $`\widehat{\text{D}}`$-brane exhibits some remarkable stability properties as given in equation (4.8). In section 6 we identify most of the decay channels for the various $`\widehat{\text{D}}`$-branes. Finally section 7 discusses D-branes on the T-dual orbifold. A new K-theory group suitable for such an orbifold is defined.
## 2 Non-compact orbifold
In this section we discuss the non-compact orbifold, while later on the directions $`x^3,\mathrm{},x^8`$ will be compactified on circles. We first briefly discuss the massless spectrum, and then construct relevant D-brane boundary states.
The $`𝖹𝖹_2\times 𝖹𝖹_2`$ orbifold discussed here has been studied in the past . In particular its closed string spectrum has been worked out in some detail . Here we point out some of the points of present importance. Each $`g_i`$, $`i=1,2,3`$ gives rise to a twisted sector. Both the twisted NS and R sectors are massless and have zero modes, those of the R sector are in the directions unaffected by $`g_i`$, while the zero modes in the NS sector are in the directions inverted by $`g_i`$. The lowest lying states in the twisted R-R sector transform as a tensor product of $`SO(4)`$ (which contains the spacetime SO(2)), and are further required to be invariant under the remaining orbifold projections. The twisted NS-NS sector ground states transform as bispinors of an internal $`SO(4)`$, and have to be invariant under the other elements of the orbifold group. They give rise to four dimensional scalars. We denote the $`g_i`$-twisted sectors by NS-NS,T$`g_i`$ and R-R,T$`g_i`$.
Boundary states corresponding to physical D-branes consist of linear combinations of boundary states from the various closed string sectors which are GSO and orbifold invariant. We refer to a D-brane as of type $`(r;𝐬)=(r;s_1,s_2,s_3)`$, if it is a $`r+s_1+s_2+s_3`$-brane, and extends along $`r+1`$, $`s_1`$, $`s_2`$, $`s_3`$ of the directions $`(x^0,\mathrm{},x^2,x^9)`$, $`(x^3,x^4`$), $`(x^5,x^6)`$, $`(x^7,x^8)`$, respectively.<sup>3</sup><sup>3</sup>3With this convention $`g_i`$ acts trivially on the $`s_i`$ directions. Due to the presence of fermionic zero modes in the various sectors requiring GSO and orbifold invariance places restrictions on $`r`$ and the $`s_i`$. These conditions are analysed in Appendix B.
A boundary state which corresponds to a physical D-brane has to satisfy certain consistency criteria. In particular the string with endpoints on the D-brane has to be a suitably projected open string. It turns out that in the $`𝖹𝖹_2\times 𝖹𝖹_2`$ orbifold there are four kinds of combinations of boundary states from the various closed string sectors which correspond to physical D-branes. These include a bulk D-brane
$$|B(r;𝐬)_{\text{NS-NS }}+\epsilon |B(r;𝐬)_{\text{R-R }},$$
(2.1)
where $`\epsilon =\pm 1`$, and a fractional D-brane
$$|B(r𝐬;)_{\text{NS-NS }}+\epsilon |B(r;𝐬)_{\text{R-R }}+\underset{i=1}{\overset{3}{}}\epsilon _i(|B(r;𝐬)_{\text{NS-NS,T}g_i}+\epsilon |B(r;𝐬)_{\text{R-R,T}g_i}),$$
(2.2)
for which $`\epsilon _i\epsilon _j=\epsilon _k`$ for $`ijki`$; both these objects are BPS as can be seen for example from the vanishing of the cylinder amplitudes computed in appendix A. In the above $`\epsilon ,\epsilon _i`$ indicate the sign of the various R-R charges. Further there are truncated non-BPS $`\widehat{\text{D}}`$-branes charged under one type of twisted R-R field
$$|B(r;𝐬)_{\text{NS-NS }}+\epsilon |B(r;𝐬)_{\text{R-R,T}g_i},i=1,2,3.$$
(2.3)
Their boundary state is very similar to that of the $`\widehat{\text{D}}`$-branes in . Finally there is a second type of truncated non-BPS $`\widehat{\text{D}}`$-brane whose boundary state is
$`|B(r;𝐬)_{\text{NS-NS }}+\epsilon _i|B(r;𝐬)_{\text{R-R,T}g_i}+\epsilon _i\epsilon _j|B(r;𝐬)_{\text{NS-NS,T}g_k}+\epsilon _j|B(r;𝐬)_{\text{R-R,T}g_j},`$
$`i,j,k=1,2,3,\text{s.t. }ϵ_{ijk}0.`$ (2.4)
This is a new kind of $`\widehat{\text{D}}`$-brane which couples to a twisted NS-NS sector as well as to twisted R-R sectors.
A second consistency condition is that a string beginning on any one of the above branes and ending on a different one must describe an open string. In this condition ensured that for any given $`(r,s)`$ if a fractional and truncated $`(r,s)`$-branes existed then the fractional brane would be of smaller mass and charge rendering the truncated brane unstable. Applying the condition to the branes above we again conclude that if a fractional and truncated brane exist for some $`(r;𝐬)`$ then it is the fractional object that is minimal and stable. Similarly one shows that if for some $`(r;𝐬)`$ both truncated objects exist it is the one in equation (2.4) that is minimally charged and so is fundamental.
Given these consistency conditions the D-brane spectrum can then be determined. The detailed analysis can be found in appendix B. Bulk and fractional D-branes exist for $`(r;𝐬)`$ of the form
$$(r;0,0,0),(r;0,0,2),(r;0,2,0),(r;2,0,0),(r;0,2,2),(r;2,0,2),(r;2,2,0),(r;2,2,2),$$
(2.5)
where $`r`$ is even/odd for Type IIA/IIB. The elementary objects above are totally fractional branes which live on the fixed point of the whole orbifold group ($`x^3=\mathrm{}=x^8=0`$) and are charged under all three twisted R-R sectors. The corresponding K-theory group should be $`𝖹𝖹^4`$. The totally fractional branes correspond to a single brane in the covering space. Two such D-branes with the same bulk and $`g_i`$-twisted R-R charges but opposite remaining charges can come together to form a singly fractional brane, which is stuck to the fixed plane of $`g_i`$. Two singly fractional branes with opposite twisted R-R charge can come together and form a bulk brane. Further bulk D-branes exist for $`(r;1,1,1)`$. Since these are charged under the untwisted R-R field, the corresponding K-group should be $`𝖹𝖹`$.
Truncated $`\widehat{\text{D}}`$-branes with a boundary state given by equation (2.3) exist for $`(r;𝐬)`$ of the form<sup>4</sup><sup>4</sup>4We disregard here the tachyon that arises when considering decompactified $`s_i>0`$ truncated branes , as on compactification the D-brane will stabilise for certain values of the compactification radii.
$$(r;0,1,1),(r;1,0,1),(r;1,1,0),(r;2,1,1),(r;1,2,1),(r;1,1,2),$$
(2.6)
with $`r`$ even/odd for Type IIA/IIB. These $`\widehat{\text{D}}`$-branes are stuck at the fixed points of the $`g_i`$ under which they are charged, and one expects the corresponding K-group to be $`𝖹𝖹`$. Truncated $`\widehat{\text{D}}`$-branes with a boundary state given by equation (2.4) exist for $`(r;𝐬)`$ of the form
$`(r;0,0,1),(r;1,0,0),(r;0,1,0),(r;2,0,1),(r;1,2,0),(r;2,1,0),`$
$`(r;0,2,1),(r;1,0,2),(r;0,1,2),(r;2,2,1),(r;1,2,2),(r;2,1,2).`$ (2.7)
Here $`r`$ is even/odd for Type IIA/IIB. The basic such branes are stuck at the fixed points of the whole orbifold group and are charged under two twisted R-R fields suggesting that the K-group should be $`𝖹𝖹𝖹𝖹`$. However, as with the fractional branes there are also $`\widehat{\text{D}}`$-branes with the above $`(r;𝐬)`$ which are charged under only $`g_i`$-twisted R-R fields and as such are only stuck to the fixed points of $`g_i`$.
## 3 K-theory analysis in uncompactified theory
The K-groups relevant to this orbifold are the $`𝖹𝖹_2\times 𝖹𝖹_2`$-equivariant K-groups with compact support
$$K_{𝖹𝖹_2\times 𝖹𝖹_2}^{}(\mathrm{I}\mathrm{R}^{a;b,c,d}),$$
(3.1)
where $`a=0,\mathrm{},4`$ and $`b,c,d=0,1,2`$. In the above the directions $`a`$ are left invariant by $`G=𝖹𝖹_2\times 𝖹𝖹_2`$, $`c`$ and $`d`$ are inverted by the first $`𝖹𝖹_2`$ and $`b,c`$ are inverted by the second $`𝖹𝖹_2`$. These groups exhibit complex Bott periodicity in $`a,b,c`$ and $`d`$ thus the answer depends only on the parity of $`a,b,c`$ and $`d`$. Further there is a symmetry between $`b,c,d`$. For example $`K_G^{}(\mathrm{I}\mathrm{R}^{a;b,c,d})=K_G^{}(\mathrm{I}\mathrm{R}^{a;c,b,d})`$. As a result we need only compute very few terms. Since the representation ring of $`𝖹𝖹_2\times 𝖹𝖹_2`$ is $`𝖹𝖹^4`$ we have
$$K_G^{}(\mathrm{I}\mathrm{R}^{a;b,c,d})=\{\begin{array}{cc}𝖹𝖹^4& a+\text{ even}\\ 0& a+\text{ odd}\end{array}$$
(3.2)
for $`b,c,d`$ even. Fractional branes in the decompactified theory can be charged under three twisted R-R fields as well as the untwisted R-R field, the above K-groups confirm the presence of all fractional branes. Since $`g_1`$ acts trivially on $`\mathrm{I}\mathrm{R}^{0;1,0,0}`$ we have
$$K_G^{}(\mathrm{I}\mathrm{R}^{0;1,0,0})=K_{𝖹𝖹_2}^{}(\mathrm{I}\mathrm{R}^{0,1})R[𝖹𝖹_2],$$
(3.3)
where on the right hand side $`𝖹𝖹_2`$ inverts the line $`\mathrm{I}\mathrm{R}^{0,1}`$. The right-hand side groups have been obtained in . Thus we have for $`c,d`$ even
$$K_G^{}(\mathrm{I}\mathrm{R}^{a;1,c,d})=\{\begin{array}{cc}𝖹𝖹^2& a+\text{ even}\\ 0& a+\text{ odd}\end{array}$$
(3.4)
Similar results hold for permutations of $`1,c,d`$. This is in agreement with the presence of truncated $`\widehat{\text{D}}`$-branes in (2.7). In order to compute $`K_G^{}(\mathrm{I}\mathrm{R}^{0;1,1,2k})`$ we re-write it as
$$K_G^{}(\mathrm{I}\mathrm{R}^{0;1,1,0})=K_G^{}(\mathrm{I}\mathrm{R}^{0;1,0,0}\times D^1,\mathrm{I}\mathrm{R}^{a;1,0,0}\times S^0),$$
(3.5)
where $`S^0D^1\mathrm{I}\mathrm{R}^{0;0,1,0}`$ are the zero-sphere (two points) and the one-disc (the interval). By homotopy equivalence we have
$$K_G^{}(X\times D^1)=K_G^{}(X),$$
(3.6)
for any $`X`$ and
$$K_G^{}(\mathrm{I}\mathrm{R}^{0;1,0,0}\times S^0)=K_{𝖹𝖹_2}^{}(\mathrm{I}\mathrm{R}^{0,1}),$$
(3.7)
where on the right hand side $`𝖹𝖹_2`$ inverts the real line $`\mathrm{I}\mathrm{R}^{0,1}`$. We may now use the long exact sequence
$`\mathrm{}`$ $``$ $`K_G^1(\mathrm{I}\mathrm{R}^{0;1,0,0}\times S^0)K_G(\mathrm{I}\mathrm{R}^{0;1,0,0}\times \mathrm{I}\mathrm{R}^{0;0,1,0})K_G(\mathrm{I}\mathrm{R}^{0;1,0,0}\times D^1)`$ (3.8)
$``$ $`K_G(\mathrm{I}\mathrm{R}^{0;1,0,0}\times S^0)K_G^1(\mathrm{I}\mathrm{R}^{0;1,0,0}\times \mathrm{I}\mathrm{R}^{0;0,1,0})K_G^1(\mathrm{I}\mathrm{R}^{0;1,0,0}\times D^1)\mathrm{}`$
which in turn becomes
$$0K_G(\mathrm{I}\mathrm{R}^{0;1,1,0})𝖹𝖹𝖹𝖹𝖹𝖹K_G^1(\mathrm{I}\mathrm{R}^{0;1,1,0})0.$$
(3.9)
It is not difficult to see that the map $`𝖹𝖹𝖹𝖹𝖹𝖹`$ is surjective. This gives for $`d`$ even
$$K_G^{}(\mathrm{I}\mathrm{R}^{a;1,1,d})=\{\begin{array}{cc}𝖹𝖹& a+\text{ even}\\ 0& a+\text{ odd}\end{array}$$
(3.10)
Together with the permutations of $`1,1,d`$ this confirms the presence of truncated $`\widehat{\text{D}}`$-branes in equation (2.6). Finally we have
$$K_G^{}(\mathrm{I}\mathrm{R}^{0;1,1,1})=K_G^{}(\mathrm{I}\mathrm{R}^{0;1,1,0}\times D^1,\mathrm{I}\mathrm{R}^{a;1,1,0}\times S^0),$$
(3.11)
where $`S^0D^1\mathrm{I}\mathrm{R}^{0;0,0,1}`$. Using equation (3.6) the results above and
$$K_G^{}(\mathrm{I}\mathrm{R}^{0;1,1,0}\times 𝒮^0)=K_{𝖹𝖹_2}^{}(\mathrm{I}\mathrm{R}^{0,2}),$$
(3.12)
the following sequence is exact
$$0K_G(\mathrm{I}\mathrm{R}^{0;1,1,1})𝖹𝖹𝖹𝖹𝖹𝖹K_G^1(\mathrm{I}\mathrm{R}^{0;1,1,1})0.$$
(3.13)
One can show using this exact sequence that
$$K_G^{}(\mathrm{I}\mathrm{R}^{a;1,1,1})=\{\begin{array}{cc}𝖹𝖹& a+\text{ odd}\\ 0& a+\text{ even}\end{array}$$
(3.14)
This confirms the presence of bulk BPS branes with $`s_1=s_2=s_3=1`$. We have demonstrated complete agreement between K-theory and D-branes on this orbifold.
## 4 The compactified orbifold
From now on we take the directions $`x^3,\mathrm{},x^8`$ to be compactified on circles of radii $`R_3,\mathrm{},R_8`$. This introduces new fixed points at $`x^i=\pi R_i`$ and gives rise to new twisted sectors - in total 48 twisted R-R sectors and 48 twisted NS-NS sectors. D-branes with $`s_i>0`$ will be charged under the twisted sectors over which they stretch. The structure of the boundary states and in particular the normalisation is described in appendix A. The allowed charges for various D-branes are restricted by factorisation. In particular a brane charged under many R-R charges cannot have an arbitrary choice of minimal positive and negative charges. This was already encountered in where for example a $`\widehat{\text{D}}`$$`(r,2)`$-brane in the $`_n(1)^{F_l}`$ orbifolds was charged under four twisted R-R fields but the minimally charged branes had to have an even number of negative charges. This behaviour is typical and we encounter it here as well. The restriction arises as a result of cylinder-annulus consistency. The process of checking what sign freedom one has in the closed string sector in order to factorise on a consistent open string amplitude is laborious, here we summarise the results.
In the following subsections we shall discuss the D-brane spectrum of the compactified orbifold, paying attention to the allowed charges, as well as stability regions and decay products of the non-BPS branes. For simplicity we concentrate on the fractional D$`(r;0,0,0)`$ and D$`(r;0,0,2)`$-branes and truncated $`\widehat{\text{D}}`$$`(r;0,0,1)`$ and $`\widehat{\text{D}}`$$`(r;0,1,1)`$-branes. The extension of these results to the other branes is obvious.
### 4.1 The D$`(r;0,0,0)`$-brane
The fully fractional D$`(r;0,0,0)`$-brane’s consistent boundary state is given by
$`|D(r;0,0,0),\epsilon ,\epsilon _i`$ $`=`$ $`|D(r;0,0,0)_{\text{NS-NS }}+\epsilon |D(r;0,0,0)_{\text{R-R}}`$ (4.1)
$`+{\displaystyle \underset{i=1}{\overset{3}{}}}\epsilon _i\left(|D(r;0,0,0)_{\text{NS-NS,T}g_i}+\epsilon |D(r;0,0,0)_{\text{R-R,T}g_i}\right),`$
with $`\epsilon _3=\epsilon _1\epsilon _2=\pm 1`$. If we denote by $`[a;b,c,d]`$ the four charges under which the D-brane is charged (with $`a`$ corresponding to the bulk charge, and $`b,c,d`$ to the twisted charges) the allowed configurations of minimal charge are
$`[1;1,1,1],[1;1,1,1],[1;1,1,1],[1;1,1,1],`$
$`[1;1,1,1],[1;1,1,1],[1;1,1,1],[1;1,1,1].`$ (4.2)
From this it is easy to see that a singly fractional brane is a sum of two totally fractional branes. For example $`[2;0,2,0]=[1;1,1,1]+[1;1,1,1]`$ or $`[2;0,2,0]=[1;1,1,1]+[1;1,1,1]`$.
### 4.2 The D$`(r;0,0,2)`$-brane
We take the fully fractional D$`(r;0,0,2)`$-brane to extend along $`x^7`$ and $`x^8`$ and to be fixed at $`x^3=\mathrm{}=x^6=0`$. Its’ boundary state is
$`|D(r;0,0,2),\epsilon ,\epsilon _i,\theta _j`$ $`=`$ $`{\displaystyle \underset{w_7,w_8}{}}e^{i(\theta _7w_7+\theta _8w_8)}(|D(r;0,0,2),w_7,w_8_{\text{NS-NS }}+\epsilon |D(r;0,0,2),w_7,w_8_{\text{R-R }})`$ (4.3)
$`+{\displaystyle \underset{i\{1,2\}}{}}(\epsilon _i(|D(r;0,0,2)_{\text{NS-NS,T}g_{i,1}}+\epsilon |D(r;0,0,2)_{\text{R-R,T}g_{i,1}})`$
$`+\epsilon _ie^{i\theta _7}(|D(r;0,0,2)_{\text{NS-NS,T}g_{i,2}}+\epsilon |D(r;0,0,2)_{\text{R-R,T}g_{i,2}})`$
$`+\epsilon _ie^{i\theta _8}(|D(r;0,0,2)_{\text{NS-NS,T}g_{i,3}}+\epsilon |D(r;0,0,2)_{\text{R-R,T}g_{i,3}})`$
$`+\epsilon _ie^{i(\theta _7+\theta _8)}(|D(r;0,0,2)_{\text{NS-NS,T}g_{i,4}}+\epsilon |D(r;0,0,2)_{\text{R-R,T}g_{i,4}}))`$
$`+\epsilon _1\epsilon _2{\displaystyle \underset{w_7,w_8}{}}e^{i(\theta _7w_7+\theta _8w_8)}(|D(r;0,0,2),w_7,w_8_{\text{NS-NS,T}g_3}`$
$`+\epsilon |D(r;0,0,2),w_7,w_8_{\text{R-R,T}g_3})`$
From this it is apparent that D$`(r;0,0,2)`$-branes can only have an even number of negative R-R,T$`g_1`$ and R-R,T$`g_2`$ charges, and further that the sign of the R-R,T$`g_3`$ charge is fixed by the other twisted charges. A D$`(r;0,0,2)`$-brane with all positive R-R,T$`g_1`$ and bulk R-R charges, but negative remaining charges (take $`\epsilon =\epsilon _1=+1`$, $`\theta _7=\theta _8=0`$, $`\epsilon _3=1`$) can be added to a D$`(r;0,0,2)`$-brane with all positive charges, to give a brane which is only charged under bulk R-R and R-R,T$`g_1`$ fields, and so is a singly fractional brane (it can move off the $`x^3=x^4=0`$ fixed plane).
### 4.3 The $`\widehat{\text{D}}`$$`(r;0,0,1)`$-brane
We consider the $`\widehat{\text{D}}`$$`(r;0,0,1)`$-brane to stretch along the $`x^8`$ direction and to be fixed at $`x^3=\mathrm{}=x^7=0`$. In order to be consistent with the open string partition function the $`\widehat{\text{D}}`$$`(r;0,0,1)`$-brane boundary state must be written in the form
$`|\widehat{D}(r;0,0,1),\epsilon _i,\theta `$ $`=`$ $`{\displaystyle \underset{w}{}}e^{i\theta w}|\widehat{D}(r;0,0,1),w_{\text{NS-NS }}+\epsilon _1|\widehat{D}(r;0,0,1)_{\text{R-R,T}g_{1,1}}`$
$`+\epsilon _1e^{i\theta }|\widehat{D}(r;0,0,1)_{\text{R-R,T}g_{1,2}}+\epsilon _2|\widehat{D}(r;0,0,1)_{\text{R-R,T}g_{2,1}}`$
$`+\epsilon _2e^{i\theta }|\widehat{D}(r;0,0,1)_{\text{R-R,T}g_{2,2}}+\epsilon _1\epsilon _2{\displaystyle \underset{w}{}}e^{i\theta w}|\widehat{D}(r;0,0,1),w_{\text{NS-NS,T}g_3}.`$
Denoting by $`[a,b;c,d]`$ the four charges a $`\widehat{\text{D}}`$$`(r;0,0,1)`$-brane carries (two from the $`g_1`$-twisted sector, say $`a,b`$ and two from the $`g_2`$ twisted sector, say $`c,d`$), the allowed $`\widehat{\text{D}}`$$`(r;0,0,1)`$-branes with minimal charge are
$`[1,1;1,1],[1,1;1,1],[1,1;1,1],[1,1;1,1],`$
$`[1,1;1,1],[1,1;1,1],[1,1;1,1],[1,1;1,1].`$ (4.5)
A singly charged $`\widehat{\text{D}}`$$`(r;0,0,1)`$-brane with minimal charge comes from $`[1,1;1,1]+[1,1;1,1]=[2,2;0,0]`$ or $`[1,1;1,1]+[1,1,1,1]=[2,2;0,0]`$.
Since the stability region of this $`\widehat{\text{D}}`$-brane is so different from the non-BPS branes previously encountered we discuss this case in detail. The partition function for a string with both end-points on this $`\widehat{\text{D}}`$-brane is
$`{\displaystyle \frac{dt}{2t}}`$ $`\mathrm{Tr}_{NSR}\left({\displaystyle \frac{1+(1)^Fg_1}{2}}{\displaystyle \frac{1+(1)^Fg_3}{2}}e^{2tH_o}\right)`$ (4.6)
$`=`$ $`{\displaystyle \frac{V_{r+1}}{4(2\pi )^{r+1}}}{\displaystyle \frac{dt}{2t}(2t)^{(r+1)/2}\left(\frac{f_3^8(\stackrel{~}{q})f_2^8(\stackrel{~}{q})}{f_1^8(\stackrel{~}{q})}\right)\underset{n_8𝖹𝖹}{}e^{2t\pi n_8^2/R_8^2}\underset{m=3}{\overset{7}{}}\underset{w_m𝖹𝖹}{}e^{2t\pi w_m^2R_m^2}}`$
$`{\displaystyle \frac{V_{r+1}}{(2\pi )^{r+1}}}{\displaystyle \frac{dt}{2t}(2t)^{(r+1)/2}\frac{f_3^4(\stackrel{~}{q})f_4^4(\stackrel{~}{q})}{f_1^4(\stackrel{~}{q})f_2^4(\stackrel{~}{q})}\underset{m=3,4}{}\underset{w_m𝖹𝖹}{}e^{2t\pi w_m^2R_m^2}}`$
$`{\displaystyle \frac{V_{r+1}}{(2\pi )^{r+1}}}{\displaystyle \frac{dt}{2t}(2t)^{(r+1)/2}\frac{f_3^4(\stackrel{~}{q})f_4^4(\stackrel{~}{q})}{f_1^4(\stackrel{~}{q})f_2^4(\stackrel{~}{q})}\underset{m=5,6}{}\underset{w_m𝖹𝖹}{}e^{2t\pi w_m^2R_m^2}}`$
$`+{\displaystyle \frac{V_{r+1}}{(2\pi )^{r+1}}}{\displaystyle \frac{dt}{2t}(2t)^{(r+1)/2}\frac{f_3^4(\stackrel{~}{q})f_4^4(\stackrel{~}{q})}{f_1^4(\stackrel{~}{q})f_2^4(\stackrel{~}{q})}\underset{w_7𝖹𝖹}{}e^{2t\pi w_7^2R_7^2}\underset{n_8𝖹𝖹}{}e^{2t\pi n_8^2/R_8^2}}`$
$``$ $`{\displaystyle \frac{V_{r+1}}{4(2\pi )^{r+1}}}{\displaystyle }{\displaystyle \frac{dt}{2t}}(2t)^{(r+1)/2}\stackrel{~}{q}^1(\stackrel{~}{q}^{2(n_8/R_8)^2}+\stackrel{~}{q}^{2(w_7R_7)^2}+{\displaystyle \underset{i=3}{\overset{7}{}}}\stackrel{~}{q}^{2(w_iR_i)^2}(\stackrel{~}{q}^{2(n_8/R_8)^2}+\stackrel{~}{q}^{2(w_7R_7)^2})`$
$`+\stackrel{~}{q}^{2(w_3R_3)^2+2(w_5R_5)^2}+\stackrel{~}{q}^{2(w_3R_3)^2+2(w_6R_6)^2}+\stackrel{~}{q}^{2(w_4R_4)^2+2(w_5R_5)^2}+\stackrel{~}{q}^{2(w_4R_4)^2+2(w_6R_6)^2}).`$
We have expanded to relevant order in the last line of the amplitude. It is then clear that the tachyon cancels provided<sup>5</sup><sup>5</sup>5It is straightforward to see how this generalises to other $`\widehat{\text{D}}`$$`(0;2k,2k^{},1)`$-branes. For example for a $`\widehat{\text{D}}`$$`(0;2,0,1)`$-brane stretching along the directions $`x^3,x^4,x^8`$ the stability region is $`R_7{\displaystyle \frac{1}{\sqrt{2}}},R_8\sqrt{2},`$ $`{\displaystyle \frac{1}{R_3^2}}+R_5^2{\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{R_4^2}}+R_5^2{\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{R_3^2}}+R_6^2{\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{R_4^2}}+R_6^2{\displaystyle \frac{1}{2}}.`$
$`R_7{\displaystyle \frac{1}{\sqrt{2}}},R_8\sqrt{2},`$ (4.7)
$`R_3^2+R_5^2{\displaystyle \frac{1}{2}},R_4^2+R_5^2{\displaystyle \frac{1}{2}},R_3^2+R_6^2{\displaystyle \frac{1}{2}},R_4^2+R_6^2{\displaystyle \frac{1}{2}}.`$ (4.8)
Note that by fixing say $`R_3,R_4`$ to suitable values $`R_5,R_6`$ can take on any value and the truncated D-brane will still be stable.
In order to analyse the decay products of this $`\widehat{\text{D}}`$-brane we take it to carry a unit of positive twisted R-R charge in each of the four sectors under which it is charged (two R-R,T$`g_1`$ and two R-R,T$`g_2`$). As can be seen from the above equations there are three distinct decay channels: $`R_8`$ can increase beyond $`\sqrt{2}`$, $`R_7`$ can decrease below $`1/\sqrt{2}`$, and the bounds for $`R_3,R_4,R_5,R_6`$ in equation (4.8) can be violated.
For $`R_8>\sqrt{2}`$ the $`\widehat{\text{D}}`$$`(r;0,0,1)`$-brane decays, conserving mass and charge, into a pair of totally fractional D$`(r;0,0,0)`$-branes, with opposite bulk and R-R,T$`g_3`$ charge. In the notation of equation (4.1) one of the fractional D-branes has $`\epsilon =\epsilon _1=\epsilon _2=\epsilon _3=1`$, and the other has $`\epsilon =\epsilon _1=\epsilon _2=\epsilon _3=1`$. The easiest way to compare the mass and charges of the original truncated brane with those of the decay products is by perusal of the relevant normalisation constants computed in appendix A. The normalisation constants of the $`\widehat{\text{D}}`$$`(r;0,0,1)`$-brane are (cf. equations (A.36) and (A.37))
$`𝒩_{td,U}^2`$ $`=`$ $`{\displaystyle \frac{V_{r+1}}{(2\pi )^{r+1}}}{\displaystyle \frac{1}{64}}{\displaystyle \frac{R_8}{_{m=3}^7R_m}},`$ (4.9)
$`𝒩_{td,Tg_1,i}^2`$ $`=`$ $`{\displaystyle \frac{V_{r+1}}{(2\pi )^{r+1}}}{\displaystyle \frac{2^3}{64}}{\displaystyle \frac{1}{R_3R_4}},i=1,2,`$ (4.10)
$`𝒩_{td,Tg_3,i}^2`$ $`=`$ $`{\displaystyle \frac{V_{r+1}}{(2\pi )^{r+1}}}{\displaystyle \frac{2^3}{64}}{\displaystyle \frac{1}{R_5R_6}},i=1,2,`$ (4.11)
while those of the two D$`(r;0,0,0)`$-branes (cf. equations (A.32) and (A.33))
$`𝒩_{f,U}^2`$ $`=`$ $`4{\displaystyle \frac{V_{r+1}}{(2\pi )^{r+1}}}{\displaystyle \frac{1}{128}}{\displaystyle \frac{1}{_{m=3}^8R_m}},`$ (4.12)
$`𝒩_{f,Tg_1,i}^2`$ $`=`$ $`{\displaystyle \frac{V_{r+1}}{(2\pi )^{r+1}}}{\displaystyle \frac{2^4}{128}}{\displaystyle \frac{1}{R_3R_4}},i=1,2,`$ (4.13)
$`𝒩_{f,Tg_3,i}^2`$ $`=`$ $`{\displaystyle \frac{V_{r+1}}{(2\pi )^{r+1}}}{\displaystyle \frac{2^4}{128}}{\displaystyle \frac{1}{R_5R_6}},i=1,2.`$ (4.14)
The factor of 4 comes from the fact that there are two BPS branes. These agree at the critical radius.
The second decay channel occurs when we decrease $`R_7`$ below $`1/\sqrt{2}`$. The decay products are a pair of totally fractional D$`(r;0,0,2)`$-branes. In the notation of section 4.3 one of the D-branes has $`ϵ,ϵ_1,ϵ_3=+1`$, and $`\theta _7,\theta _8=0`$, while the other brane has $`ϵ,ϵ_1,ϵ_3=1`$ and $`\theta _7,\theta _8=0`$. Comparing the normalisation constants confirms that charge and mass are conserved in the decay.
The third decay channel seems much more complicated. Naively one might expect a decay into a brane-anti-brane combination of fractional branes stretching along $`x^8`$ and one of $`x^3,x^4,x^5,x^6`$. However, there are no such fractional branes. If the dacay is to be into branes with two internal directions, then by considering the Dirac quantisation condition these will have to be fractional branes, which presumably will be somehow bent away from the axes, and not have a boundary state description of the type analysed in this paper.
It is also possible to analyse the stability of the $`\widehat{\text{D}}`$$`(r;0,0,1)`$-brane charged only under,say, the $`g_1`$-twisted R-R sector. Its stability region is
$$R_5,R_6,R_7\frac{1}{\sqrt{2}},R_8\sqrt{2},$$
(4.15)
and the $`R_7`$ and $`R_8`$ decay channels follow easily from the decay channels of the $`\widehat{\text{D}}`$$`(r;0,0,1)`$-brane charged under the $`g_1`$ and $`g_2`$-twisted R-R sectors discussed above.
### 4.4 The $`\widehat{\text{D}}`$$`(r;0,1,1)`$-brane
Finally, we discuss a $`\widehat{\text{D}}`$$`(r;0,1,1)`$-brane which we take to extend along $`x^8`$ and $`x^6`$, and be fixed at $`x^5=x^7=0`$ and to consist of two identical copies at $`(x^3,x^4)`$ and $`(x^3,x^4)`$, as explained in appendix A. The allowed boundary states are of the form
$`|\widehat{D}(r;0,1,1),\epsilon _i,\theta _j`$ $`=`$ $`{\displaystyle \underset{w_6,w_8}{}}e^{i(\theta _6w_6+\theta _8w_8)}|\widehat{D}(r;0,1,1),w_6,w_8_{\text{NS-NS}}`$ (4.16)
$`+\epsilon _1|\widehat{D}(r;0,1,1)_{\text{R-R,T}g_{1,1}}+\epsilon _1e^{i\theta _6}|\widehat{D}(r;0,1,1)_{\text{R-R,T}g_{1,2}}`$
$`+\epsilon _1e^{i\theta _8}|\widehat{D}(r;0,1,1)_{\text{R-R,T}g_{1,3}}+\epsilon _1e^{i(\theta _6+\theta _8)}|\widehat{D}(r;0,1,1)_{\text{R-R,T}g_{1,4}}.`$
Hence, such truncated branes can only carry an even number of negative twisted R-R charges. The stability of this brane is very similar to those encountered in $`𝖹𝖹_2`$ orbifolds, we simply state the results here. It is stable for
$$R_5,R_7\frac{1}{\sqrt{2}},R_6,R_8\sqrt{2}.$$
(4.17)
The $`\widehat{\text{D}}`$-brane can decay in two different ways: by increasing $`R_6`$ or $`R_8`$ beyond $`\sqrt{2}`$, or by decreasing $`R_5`$ or $`R_7`$ below $`1/\sqrt{2}`$. In the first the decay is into a pair of singly charged $`\widehat{\text{D}}`$$`(r;0,0,1)`$-branes, which carry the same charges as the $`\widehat{\text{D}}`$$`(r;0,1,1)`$-brane. It is straightforward to check that at the critical radius ($`R_6=\sqrt{2}`$) mass and charge are conserved. The second decay channel corresponds to decreasing $`R_5`$ below $`1/\sqrt{2}`$. The $`\widehat{\text{D}}`$-brane decays into a pair of singly charged $`\widehat{\text{D}}`$$`(r;0,2,1)`$-branes, whose charges are the same at the four fixed points on which the $`\widehat{\text{D}}`$$`(r;0,1,1)`$-brane ends and opposite at the other four fixed points. Again both mass and charge are conserved at the transition. The decay products of the $`\widehat{\text{D}}`$$`(r;0,1,1)`$-branes are different from the decay products previously encountered. In particular the brane does not decay into a brane-anti-brane pair of separately BPS objects, rather the decay is into other non-BPS objects.
## 5 T-duality
In this section we discuss the T-dual of the $`𝖹𝖹_2\times 𝖹𝖹_2`$ orbifold. We shall T-dualise along the $`x^5`$ direction. The orbifold we consider will then be a $`𝖹𝖹_2\times 𝖹𝖹_2`$ orbifold where the generators are $`g_1=_{5678}(1)^{F_l}`$ and $`g_2=_{3478}`$, where $`F_l`$ is the left moving space-time fermion number. One imposes GSO and orbifold invariance of the boundary states in the different closed string sectors. As in the $`_4(1)^{F_l}`$ orbifolds there is an ambiguity regarding which part of the twisted sector states to keep: the $`_4(1)^{F_l}`$ odd states or the $`_4(1)^{F_l}`$ even states? For the case of the invariance under $`g_i`$ for the $`g_i`$-twisted sector we follow the prescription of . In particular we require the $`g_1`$\- and $`g_3`$\- twisted sectors to be odd under $`g_1`$ and $`g_3`$ respectively, while the $`g_2`$-twisted sectors are to be even under $`g_2`$. We further shall require the $`g_1`$-twisted sector to be $`g_2`$ even and $`g_3`$ odd, the $`g_2`$-twisted sector to be $`g_1`$ and $`g_3`$ even and the $`g_3`$-twisted sector to be $`g_2`$ even and $`g_1`$ odd. This ensures that we can construct fractional branes. With these clarifications one obtains table 1 for Type IIB
With these restrictions it is not difficult to identify bulk, fractional and both kinds of truncated branes. As before if two kinds of branes can exist for a given $`(r;𝐬)`$it is the fractional rather than truncated (or the truncated with the NS-NS twisted sector rather than the other truncated brane) that will be fundamental and of minimal charge. For Type IIB we then have bulk wrapped branes for $`r,s_1,s_3`$ odd $`s_2`$ even (only charged under the untwisted R-R sector), and bulk and fractional branes (latter charged under four kinds of charges in the decompactified theory) for $`r,s_1,s_3`$ even $`s_2`$ odd. Truncated branes charged under R-R,T$`g_1`$ and R-R,T$`g_2`$ fields exist for $`r,s_1`$ even $`s_2,s_3`$ odd. Truncated branes charged under R-R,T$`g_1`$ and R-R,T$`g_3`$ fields exist for $`r,s_1,s_2,s_3`$ even. Truncated branes charged under R-R,T$`g_2`$ and R-R,T$`g_3`$ fields exist for $`r,s_3`$ even $`s_1,s_2`$ odd. Truncated branes charged only under R-R,T$`g_1`$ exist for $`r,s_1,s_2`$ even $`s_3`$ odd. Truncated branes charged only under R-R,T$`g_2`$ exist for $`r`$ even $`s_1,s_2,s_3`$ odd. Truncated branes charged only under R-R,T$`g_3`$ exist for $`r,s_2,s_3`$ even $`s_1`$ odd. For Type IIA in the above the parity of $`r`$ changes.
One can develop a K-theory understanding of this. We define a new kind of K-theory which we call $`K_{𝖹𝖹_2\times 𝖹𝖹_2^\pm }^{}`$. Elements of this are pairs of isomorphism classes of bundles $`(E,F)`$, which are equivariant under the action of the first $`𝖹𝖹_2`$ and we are given an isomorphism between $`(E,F)`$ and $`(g^{}(E),g^{}(F))`$ where in this case $`g=_{5678}`$ and $`g^{}(E)`$ is the pullback of $`E`$ by $`g`$. Further we are given an isomorphism which maps $`(E,F)`$ to $`(h^{}(F),h^{}(E))`$, where $`h^{}(E)`$ is the pullback of $`E`$ by $`h`$, the generator of the geometric part of the second $`𝖹𝖹_2`$ (in other words in this case $`h=_{3478}`$). A version of Hopkins’ formula then is
$$K_{𝖹𝖹_2\times 𝖹𝖹_2^\pm }^{}(\mathrm{I}\mathrm{R}^{a;b,c,d})K_{𝖹𝖹_2\times 𝖹𝖹_2}^{}(\mathrm{I}\mathrm{R}^{a+1;b,c+1,d}),$$
(5.18)
where the second group is just the usual equivariant K-group. We have computed the latter in section 3, and using the above relation exact agreement between K-theory and the string results follows. The stability and decay of the the truncated branes can easily be obtained from the results of the T-dual scenario described in detail in the previous sections.
## Appendix A Construction and normalisation of boundary states
In this appendix we determine the normalisation constants of the boundary states for the orbifold theories under consideration.
### A.1 The uncompactified case
In each (bosonic) sector of the theory we can construct the boundary state
$$|B(r;𝐬),k,\eta =\mathrm{exp}\left(\underset{l>0}{\overset{\mathrm{}}{}}\left[\frac{1}{l}\alpha _l^\mu S_{\mu \nu }\stackrel{~}{\alpha }_l^\nu \right]+i\eta \underset{m>0}{\overset{\mathrm{}}{}}\left[\psi _m^\mu S_{\mu \nu }\stackrel{~}{\psi }_m^\nu \right]\right)|B(r;𝐬),k,\eta ^{(0)},$$
(A.1)
where, depending on the sector, $`l`$ and $`m`$ are integer or half-integer, and $`k`$ denotes the momentum of the ground state. We shall always work in light-cone gauge with light-cone directions $`x^0`$ and $`x^9`$; thus $`\mu `$ and $`\nu `$ take the values $`1,\mathrm{},8`$. We shall also drop the dependence on $`\alpha ^{}`$.
The parameter $`\eta =\pm `$ describes the two different spin structures , and the matrix $`S`$ encodes the boundary conditions of the Dp-brane which we shall always take to be diagonal
$$S=\text{diag}(1,\mathrm{},1,1,\mathrm{},1),$$
(A.2)
where $`p+1`$ entries are equal to $`1`$, $`7p`$ entries are equal to $`+1`$, and $`p=r+s_1+s_2+s_3`$. If there are fermionic zero modes, the ground state in (A.1) satisfies additional conditions discussed in the following appendix.
In order to obtain a localised D-brane, we have to take the Fourier transform of the above boundary state, where we integrate over the directions transverse to the brane,
$$|B(r,𝐬),y,\eta =\left(\underset{\mu \text{transverse}}{}dk^\mu e^{ik^\mu y_\mu }\right)𝑑k^0e^{ik^0y_0}𝑑k^9e^{ik^9y_9}|B(r;𝐬),k,\eta ,$$
(A.3)
$`y`$ denotes the location of the boundary state, and in the $`g_i`$-twisted sector the momentum integral only involves transverse directions that are not inverted by the action of $`g_i`$. In the following we shall typically consider (without loss of generality) the case $`y=0`$ in which case the boundary state is denoted by $`|B(r;𝐬),\eta `$.
The invariance of the boundary state under the GSO-projection always requires that the physical boundary state is a linear combination of the two states corresponding to $`\eta =\pm `$. Using the conventions of appendix B in , these linear combinations are of the form
$`|B(r;𝐬)_{\text{NS-NS}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(|B(r;𝐬),+_{\text{NS-NS }}|B(r;𝐬),_{\text{NS-NS }}\right),`$ (A.4)
$`|B(r;𝐬)_{\text{R-R}}`$ $`=`$ $`2i\left(|B(r;𝐬),+_{\text{R-R }}+|B(r;𝐬),_{\text{R-R }}\right),`$ (A.5)
$`|B(r;𝐬)_{\text{NS-NS,T}g_i}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(|B(r;𝐬),+_{\text{NS-NS,T}g_i}+|B(r;𝐬),_{\text{NS-NS,T}g_i}\right),`$ (A.6)
$`|B(r;𝐬)_{\text{R-R,T}g_i}`$ $`=`$ $`i\left(|B(r;𝐬),+_{\text{R-R,T}g_i}+|B(r;𝐬),_{\text{R-R,T}g_i}\right),`$ (A.7)
where, depending on the theory in question, these states are actually GSO-invariant provided that $`r`$ and $`s_i`$ satisfy suitable conditions discussed in section 2 ($`i=1,2,3`$). The normalisation constants have been introduced for later convenience.
In order to solve the open-closed consistency condition the actual D-brane state is a linear combination of physical boundary states from different sectors. There are three elementary cases to consider, fully fractional, and the two truncated D-branes. (cf. equations 2.2)-(2.4) In the fully fractional case, the D-brane state can be written as
$`|D(r;𝐬)`$ $`=`$ $`𝒩_{f,U}\left(|B(r;𝐬)_{\text{NS-NS }}+ϵ|B(r;𝐬)_{\text{R-R }}\right)`$ (A.8)
$`+{\displaystyle \underset{i=1}{\overset{3}{}}}ϵ_i𝒩_{f,Tg_i}\left(|B(r;𝐬)_{\text{NS-NS,T}g_i}+ϵ|B(r;𝐬)_{\text{R-R,T}g_i}\right),`$
where $`ϵ=\pm `$ determines the sign of the charge with respect to the untwisted R-R sector charge, while $`ϵ_i=\pm `$, $`i=1,2,3`$ determines the sign of the charge with respect to the R-R,T$`g_i`$ charge. For consistency when $`ijki`$ we have $`ϵ_i=ϵ_jϵ_k`$. The closed string cylinder diagram is then of the form
$`𝒜`$ $`=`$ $`{\displaystyle 𝑑lB(r;𝐬)\left|e^{lH_c}\right|B(r;𝐬)}`$ (A.9)
$`=`$ $`{\displaystyle \frac{1}{2}}𝒩_{f,U}^2{\displaystyle 𝑑ll^{(p9)/2}\left(\frac{f_3^8(q)f_4^8(q)f_2^8(q)}{f_1^8(q)}\right)}`$
$`+{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{1}{2}}𝒩_{f,Tg_i}^2{\displaystyle 𝑑ll^{(r+s_i5)/2}\left(\frac{f_3^4(q)f_2^4(q)f_2^4(q)f_3^4(q)}{f_1^4(q)f_4^4(q)}\right)},`$
where the functions $`f_i`$ are defined as in , $`q=e^{2\pi l}`$, and the closed string Hamiltonian is given by
$$H_c=\pi k^2+2\pi \underset{\mu =1}{\overset{8}{}}\left[\underset{l>0}{\overset{\mathrm{}}{}}(\alpha _l^\mu \alpha _l^\mu +\stackrel{~}{\alpha }_l^\mu \stackrel{~}{\alpha }_l^\mu )+\underset{m>0}{\overset{\mathrm{}}{}}m(\psi _m^\mu \psi _m^\mu +\stackrel{~}{\psi }_m^\mu \stackrel{~}{\psi }_m^\mu )\right]+2\pi C_c.$$
(A.10)
Here the constant $`C_c`$ is $`1`$ in the NS-NS sector, and zero in all other sectors. The corresponding open string amplitude is obtained by the modular transformation $`t=1/2l`$, $`\stackrel{~}{q}=e^{\pi t}`$,
$`𝒜`$ $`=`$ $`2^{(7p)/2}𝒩_{f,U}^2{\displaystyle \frac{dt}{2t}t^{(p+1)/2}\left(\frac{f_3^8(\stackrel{~}{q})f_2^8(\stackrel{~}{q})f_4^8(\stackrel{~}{q})}{f_1^8(\stackrel{~}{q})}\right)}`$ (A.11)
$`+{\displaystyle \underset{i=1}{\overset{3}{}}}2^{(3rs_i)/2}𝒩_{f,Tg_i}^2{\displaystyle \frac{dt}{2t}t^{(r+s_i+1)/2}\left(\frac{f_3^4(\stackrel{~}{q})f_4^4(\stackrel{~}{q})f_4^4(\stackrel{~}{q})f_3^4(\stackrel{~}{q})}{f_1^4(\stackrel{~}{q})f_2^4(\stackrel{~}{q})}\right)}.`$
This is to be compared with the open string one-loop diagram,
$`{\displaystyle \frac{dt}{2t}}`$ $`\mathrm{Tr}_{NSR}\left({\displaystyle \frac{1+(1)^F}{2}}{\displaystyle \frac{1+g_1+g_2+g_3}{4}}e^{2tH_o}\right)`$
$`=`$ $`{\displaystyle \frac{V_{p+1}}{(2\pi )^{p+1}}}2^{(p+7)/2}{\displaystyle \frac{dt}{2t}t^{(p+1)/2}\left(\frac{f_3^8(\stackrel{~}{q})f_4^8(\stackrel{~}{q})f_2^8(\stackrel{~}{q})}{f_1^8(\stackrel{~}{q})}\right)}`$
$`+{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{V_{r+s_i+1}}{(2\pi )^{r+s_i+1}}}2^{(r+s_i+3)/2}{\displaystyle \frac{dt}{2t}t^{(r+s_i+1)/2}\left(\frac{f_3^4(\stackrel{~}{q})f_4^4(\stackrel{~}{q})f_4^4(\stackrel{~}{q})f_3^4(\stackrel{~}{q})}{f_1^4(\stackrel{~}{q})f_2^4(\stackrel{~}{q})}\right)},`$
where $`V_{p+1}`$ is the (infinite) $`p+1`$ dimensional volume of the brane, whilst $`V_{r+s_i+1}`$ is the volume of the projection onto the directions unaffected by $`g_i`$. The open string Hamiltonian is given by
$$H_o=\pi p^2+\pi \underset{\mu =1}{\overset{8}{}}\left[\underset{l>0}{\overset{\mathrm{}}{}}\alpha _l^\mu \alpha _l^\mu +\underset{m>0}{\overset{\mathrm{}}{}}m\psi _m^\mu \psi _m^\mu \right]+\pi C_o,$$
(A.13)
where, in the R sector, $`l`$ and $`m`$ run over the positive integers for NN and DD directions, and over positive half integers for ND directions. In the NS sector, the moding of the fermions (and therefore the values for $`m`$) are opposite to those in the R sector. $`C_o`$ is zero in the R sector and is $`\frac{4t}{8}`$ in the NS sector, where $`t`$ is the number of ND directions. Comparison of equations (LABEL:theopenone) and (A.11) then gives
$`𝒩_{f,U}^2`$ $`=`$ $`{\displaystyle \frac{V_{p+1}}{(2\pi )^{p+1}}}{\displaystyle \frac{1}{128}},`$ (A.14)
$`𝒩_{f,Tg_i}^2`$ $`=`$ $`{\displaystyle \frac{V_{r+s_i+1}}{(2\pi )^{r+s_i+1}}}{\displaystyle \frac{1}{8}}.`$ (A.15)
Consider next the singly fractional D-branes charged under the R-R,T$`g_i`$ field. These can be thought of as superpositions of two totally fractional D-branes with opposite twisted charges in the other two twisted sectors. These pairs can move off the fixed points of $`g_j`$ for $`ji`$, such that they lie at positions $`x^\mu `$ and $`x^\mu `$ for the directions transverse to the brane fixed by $`g_j`$ and not by $`g_i`$. In particular their boundary states will now look like
$`|D(r;𝐬)`$ $`=`$ $`𝒩_{f,U}\left(\right|B(r;𝐬),x^\mu _{\text{NS-NS }}+|B(r;𝐬),x^\mu _{\text{NS-NS }}`$ (A.16)
$`+ϵ(|B(r;𝐬),x^\mu _{\text{R-R }}+|B(r;𝐬),x^\mu _{\text{R-R }}))`$
$`+ϵ_i𝒩_{f,Tg_i}\left(\right|B(r;𝐬),x^\mu _{\text{NS-NS,T}g_i}+|B(r;𝐬),x^\mu _{\text{NS-NS,T}g_i}`$
$`+ϵ|B(r;𝐬),x^\mu _{\text{R-R,T}g_i}+ϵ|B(r;𝐬),x^\mu _{\text{R-R,T}g_i}).`$
The above normalisation is consistent with the fact that the cylinder diagram for such D-branes factorises on
$$2\frac{dt}{2t}\mathrm{Tr}_{NSR}\left(\frac{1+(1)^F}{2}\frac{1+g_i}{2}e^{2tH_o}\right).$$
(A.17)
Such singly fractional D-branes correspond to, in the covering space, two Type II D-branes stuck at the fixed point of $`g_i`$. There are four kinds of open strings stretching between two such D-branes so one would expect a factor of four in front of the trace above. However, the orbifold identifies certain strings leaving only two independent ones (namely the string that has endpoint on the same brane and the string that has ends on opposite branes). Similarly a bulk D-brane is a configuration of four totally fractional D-branes whose twisted charges cancel. It’s boundary state is a sum over four NS-NS and four R-R boundary states at positions mapped to one another under the orbifold action. The cylinder diagram for such a D-brane factorises on
$$4\frac{dt}{2t}\mathrm{Tr}_{NSR}\left(\frac{1+(1)^F}{2}e^{2tH_o}\right).$$
(A.18)
The analysis for the case of the truncated $`\widehat{\text{D}}`$-branes is similar. The truncated D-brane of the type given by equation (2.3) charged only under the R-R,T$`g_1`$ sector, say, has a boundary state
$`|\widehat{D}(r;𝐬)`$ $`=`$ $`𝒩_{ts,U}\left(|B(r;𝐬),x^\mu _{\text{NS-NS }}+|B(r;𝐬),x^\mu _{\text{NS-NS }}\right)`$ (A.19)
$`+ϵ𝒩_{ts,Tg_1}\left(|B(r;𝐬),x^\mu _{\text{R-R,T}g_1}+|B(r;𝐬),x^\mu _{\text{R-R,T}g_1}\right),`$
where $`ϵ=\pm `$ determines the sign of the R-R,T$`g_1`$ sector charge and $`x^\mu `$ denotes some of the directions $`x^3,x^4`$ which are transverse to the $`\widehat{\text{D}}`$-brane. The closed string tree diagram now only produces some of the terms of (A.9), and the corresponding open string amplitude is
$`4{\displaystyle \frac{dt}{2t}}`$ $`\mathrm{Tr}_{NSR}\left({\displaystyle \frac{1+g_1(1)^F}{2}}e^{2tH_o}\right)`$ (A.20)
$`=`$ $`4{\displaystyle \frac{V_{p+1}}{(2\pi )^{p+1}}}2^{(p+3)/2}{\displaystyle \frac{dt}{2t}t^{(p+1)/2}\left(\frac{f_3^8(\stackrel{~}{q})f_2^8(\stackrel{~}{q})}{f_1^8(\stackrel{~}{q})}\right)}`$
$`4{\displaystyle \frac{V_{r+s_1+1}}{(2\pi )^{r+s_1+1}}}2^{(1s_1r)/2}{\displaystyle \frac{dt}{2t}t^{(r+s_1+1)/2}\left(\frac{f_4^4(\stackrel{~}{q})f_3^4(\stackrel{~}{q})}{f_1^4(\stackrel{~}{q})f_2^4(\stackrel{~}{q})}\right)}.`$
Comparison with the corresponding closed string calculation then gives
$`𝒩_{t1,U}^2`$ $`=`$ $`{\displaystyle \frac{V_{p+1}}{(2\pi )^{p+1}}}{\displaystyle \frac{1}{8}},`$ (A.21)
$`𝒩_{t1,Tg_1}^2`$ $`=`$ $`{\displaystyle \frac{V_{r+s_1+1}}{(2\pi )^{r+s_1+1}}}2.`$ (A.22)
The truncated D-brane of the type given in equation (2.4) is charged under the R-R,T$`g_i`$ and R-R,T$`g_j`$ sectors $`(ij)`$ has a boundary state given by
$`|\widehat{D}(r;𝐬)`$ $`=`$ $`𝒩_{td,U}|B(r;𝐬)_{\text{NS-NS }}+ϵ_i𝒩_{td,Tg_i}|B(r;𝐬)_{\text{R-R,T}g_i}`$ (A.23)
$`+ϵ_j𝒩_{td,Tg_j}|B(r;𝐬)_{\text{R-R,T}g_j}+ϵ_iϵ_j𝒩_{td,Tg_k}|B(r;𝐬)_{\text{NS-NS,T}g_k}`$
where $`ϵ_i=\pm `$ determines the sign of the R-R,T$`g_i`$ sector charge and $`jki`$ with $`k=1,2,3`$. The closed string tree diagram produces only some of the terms of (A.9), and the corresponding open string amplitude is
$`{\displaystyle \frac{dt}{2t}}`$ $`\mathrm{Tr}_{NSR}\left({\displaystyle \frac{1+g_i(1)^F}{2}}{\displaystyle \frac{1+g_j(1)^F}{2}}e^{2tH_o}\right)`$ (A.24)
$`=`$ $`{\displaystyle \frac{V_{p+1}}{(2\pi )^{p+1}}}2^{(p+5)/2}{\displaystyle \frac{dt}{2t}t^{(p+1)/2}\left(\frac{f_3^8(\stackrel{~}{q})f_2^8(\stackrel{~}{q})}{f_1^8(\stackrel{~}{q})}\right)}`$
$`{\displaystyle \underset{\alpha \{i,j,k\}}{}}\rho (\alpha ){\displaystyle \frac{V_{r+s_\alpha +1}}{(2\pi )^{r+s_\alpha +1}}}2^{(r+s_\alpha +1)/2}{\displaystyle \frac{dt}{2t}t^{(r+s_\alpha +1)/2}\left(\frac{f_4^4(\stackrel{~}{q})f_3^4(\stackrel{~}{q})}{f_1^4(\stackrel{~}{q})f_2^4(\stackrel{~}{q})}\right)},`$
where $`\rho (i)=\rho (j)=\rho (k)=1`$. Comparison with the corresponding closed string calculation then gives
$`𝒩_{t2,U}^2`$ $`=`$ $`{\displaystyle \frac{V_{p+1}}{(2\pi )^{p+1}}}{\displaystyle \frac{1}{64}},`$ (A.25)
$`𝒩_{t2,Tg_\alpha }^2`$ $`=`$ $`{\displaystyle \frac{V_{r+s_\alpha +1}}{(2\pi )^{r+s_\alpha +1}}}{\displaystyle \frac{1}{4}}.`$ (A.26)
### A.2 The compactified case
The construction in the compactified case is essentially the same as in the above uncompactified case; however there are the following differences.
1. In the localised boundary state (A.3) the integral over compact transverse directions is replaced by a sum
$$𝑑k^\nu e^{ik^\nu y_\nu }\underset{m^\nu 𝖹𝖹}{}e^{im^\nu y_\nu /R_\nu },$$
(A.27)
where $`R_\nu `$ is the radius of the compact $`x^\nu `$ direction.
2. In the two untwisted sectors, the ground state is in addition characterised by a winding number $`w_\nu `$ for each compact direction that is tangential to the world-volume of the brane. In the $`g_i`$ twisted sectors, the ground states will also be characterised by winding numbers in each of the $`s_i`$ directions tangential to the world-volume of the brane. The localised bound state (A.3) then also contains a sum over these winding states
$$\underset{w_\mu }{}e^{i\theta ^\mu w_\mu },$$
(A.28)
where $`\theta ^\mu `$ is a Wilson line; as required by orbifold invariance, $`\theta ^\mu \{0,\pi \}`$.
3. For general $`s_i`$, the contribution in the twisted sectors consists of a sum of terms that are associated to $`2^{s_i}`$ of the $`16`$ different twisted sectors that define the endpoints of the world-volume of the brane in the internal space. For convenience we may assume that one of the $`2^{s_i}`$ fixed points is always the origin.
4. The open and closed string Hamiltonians, $`H_o`$ and $`H_c`$, each acquire an extra term $`1/4\pi (_\mu w_\mu ^2)`$.
Let us now construct in more detail the boundary state for a fully fractional D$`(r;𝐬)`$ brane. This is of the form
$`|D(r;𝐬)`$ $`=`$ $`𝒩_{f,U}\left(|B(r;𝐬)_{\text{NS-NS }}+\epsilon |B(r;𝐬)_{\text{R-R }}\right)`$ (A.29)
$`+{\displaystyle \underset{i=1}{\overset{3}{}}}\epsilon _i𝒩_{f,Tg_i}{\displaystyle \underset{\alpha _i=1}{\overset{2^{s_i}}{}}}e^{i\theta _{\alpha _i}}\left(|B(r;𝐬)_{\text{NS-NS,T}g_{\alpha _i}}+ϵ|B(r;𝐬)_{\text{R-R,T}g_{\alpha _i}}\right),`$
where $`\alpha _i`$ labels the different fixed points between which the brane stretches (where we choose the convention that T$`g_{i_1}`$ is the twisted sector at the origin), and $`\theta _{\alpha _i}`$ is the Wilson line that is associated to the difference of the fixed point $`\alpha _i`$ and the origin. The closed string tree diagram is now
$`𝒜_c`$ $`=`$ $`{\displaystyle 𝑑lB(r;𝐬)\left|e^{lH_c}\right|B(r;𝐬)}`$ (A.30)
$`=`$ $`{\displaystyle \frac{1}{2}}𝒩_{f,U}^2{\displaystyle 𝑑ll^{(r3)/2}\left(\frac{f_3^8(q)f_2^8(q)f_4^8(q)}{f_1^8(q)}\right)}`$
$`\times {\displaystyle \underset{m=1}{\overset{s_1+s_2+s_3}{}}}{\displaystyle \underset{w_{j_m}𝖹𝖹}{}}e^{l\pi R_{j_m}^2w_{j_m}^2}{\displaystyle \underset{m=1}{\overset{6s_1s_2s_3}{}}}{\displaystyle \underset{n_{k_m}𝖹𝖹}{}}e^{l\pi (n_{k_m}/R_{k_m})^2}`$
$`+{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{2^{s_1+s_2+s_3s_i}}{2}}𝒩_{f,Tg_i}^2{\displaystyle 𝑑ll^{(r3)/2}\left(\frac{f_3^4(q)f_2^4(q)f_2^4(q)f_3^4(q)}{f_1^4(q)f_4^4(q)}\right)}`$
$`\times {\displaystyle \underset{m=1}{\overset{s_i}{}}}{\displaystyle \underset{w_{j_m}𝖹𝖹}{}}e^{l\pi R_{j_m}^2w_{j_m}^2}{\displaystyle \underset{m=1}{\overset{2s_i}{}}}{\displaystyle \underset{n_{k_m}𝖹𝖹}{}}e^{l\pi (n_{k_m}/R_{k_m})^2},`$
where $`R_{j_m}`$, $`m=1,\mathrm{},s_1+s_2+s_3`$ are the radii of the circles that are tangential to the world-volume of the brane, and $`R_{k_m}`$, $`i=1,\mathrm{},6s_1s_2s_3`$ are the radii of the directions transverse to the brane. Upon the substitution $`t=1/2l`$, using the Poisson resummation formula (see for example ), this amplitude becomes
$`𝒜_c`$ $`=`$ $`𝒩_{f,U}^2{\displaystyle \frac{\underset{m=1}{\overset{6s_1s_2s_3}{}}R_{k_m}}{_{m=1}^{s_1+s_2+s_3}R_{j_m}}}2^{(7r)/2}{\displaystyle \frac{dt}{2t}t^{(r+1)/2}\left(\frac{f_3^8(\stackrel{~}{q})f_2^8(\stackrel{~}{q})f_4^8(\stackrel{~}{q})}{f_1^8(\stackrel{~}{q})}\right)}`$
$`\times {\displaystyle \underset{m=1}{\overset{s_1+s_2+s_3}{}}}{\displaystyle \underset{n_{j_m}𝖹𝖹}{}}e^{2t\pi n_{j_m}^2/R_{j_m}^2}{\displaystyle \underset{m=1}{\overset{6s_1s_2s_3}{}}}{\displaystyle \underset{w_{k_m}𝖹𝖹}{}}e^{2t\pi w_{k_m}^2R_{k_m}^2}`$
$`+{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{\underset{m=1}{\overset{2s_i}{}}R_{k_m}}{_{m=1}^{s_i}R_{j_m}}}2^{(1r)/2}2^{s_1+s_2+s_3s_i}𝒩_{f,Tg_i}^2{\displaystyle \frac{dt}{2t}t^{(r+1)/2}\left(\frac{f_3^4(\stackrel{~}{q})f_4^4(\stackrel{~}{q})f_4^4(\stackrel{~}{q})f_3^4(\stackrel{~}{q})}{f_1^4(\stackrel{~}{q})f_2^4(\stackrel{~}{q})}\right)}`$
$`\times {\displaystyle \underset{m=1}{\overset{s_i}{}}}{\displaystyle \underset{n_{j_m}𝖹𝖹}{}}e^{2t\pi n_{j_m}^2/R_{j_m}^2}{\displaystyle \underset{m=1}{\overset{2s_i}{}}}{\displaystyle \underset{w_{k_m}𝖹𝖹}{}}e^{2t\pi w_{k_m}^2R_{k_m}^2}.`$
This is to be compared with the open string amplitude
$`{\displaystyle \frac{dt}{2t}}`$ $`\mathrm{Tr}_{NSR}\left({\displaystyle \frac{1+(1)^F}{2}}{\displaystyle \frac{1+g_1+g_2+g_3}{4}}e^{2tH_o}\right)`$ (A.31)
$`=`$ $`{\displaystyle \frac{V_{r+1}}{8(2\pi )^{r+1}}}2^{(r+1)/2}{\displaystyle \frac{dt}{2t}t^{(r+1)/2}\left(\frac{f_3^8(\stackrel{~}{q})f_4^8(\stackrel{~}{q})f_2^8(\stackrel{~}{q})}{f_1^8(\stackrel{~}{q})}\right)}`$
$`\times {\displaystyle \underset{m=1}{\overset{s_1+s_2+s_3}{}}}{\displaystyle \underset{n_{j_m}𝖹𝖹}{}}e^{2t\pi n_{j_m}^2/R_{j_m}^2}{\displaystyle \underset{m=1}{\overset{6s_1s_2s_3}{}}}{\displaystyle \underset{w_{k_m}𝖹𝖹}{}}e^{2t\pi w_{k_m}^2R_{k_m}^2}`$
$`{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{V_{r+1}}{2(2\pi )^{r+1}}}2^{(r+1)/2}{\displaystyle \frac{dt}{2t}t^{(r+1)/2}\frac{f_3^4(\stackrel{~}{q})f_4^4(\stackrel{~}{q})f_4^4(\stackrel{~}{q})f_3^4(\stackrel{~}{q})}{f_1^4(\stackrel{~}{q})f_2^4(\stackrel{~}{q})}}`$
$`\times {\displaystyle \underset{m=1}{\overset{s_i}{}}}{\displaystyle \underset{n_{j_m}𝖹𝖹}{}}e^{2t\pi n_{j_m}^2/R_{j_m}^2}{\displaystyle \underset{m=1}{\overset{2s_i}{}}}{\displaystyle \underset{w_{k_m}𝖹𝖹}{}}e^{2t\pi w_{k_m}^2R_{k_m}^2}.`$
By comparison this then fixes the normalisation constants as
$`𝒩_{f,U}^2`$ $`=`$ $`{\displaystyle \frac{V_{r+1}}{(2\pi )^{r+1}}}{\displaystyle \frac{1}{128}}{\displaystyle \frac{\underset{m=1}{\overset{s_1+s_2+s_3}{}}R_{j_m}}{_{m=1}^{6s_1s_2s_3}R_{k_m}}},`$ (A.32)
$`𝒩_{f,Tg_i}^2`$ $`=`$ $`{\displaystyle \frac{V_{r+1}}{(2\pi )^{r+1}}}{\displaystyle \frac{2^{4(s_1+s_2+s_3s_i)}}{128}}{\displaystyle \frac{\underset{m=1}{\overset{s_i}{}}R_{j_m}}{_{m=1}^{2s_i}R_{k_m}}}.`$ (A.33)
The extension of this to the singly fractional and bulk D-branes is obvious.
The analysis for the truncated D-branes is almost identical. The boundary state of a truncated D-brane of the type given in equation (2.3) is the truncation of (A.29) to the untwisted NS-NS and the $`g_i`$-twisted R-R sectors with the addition of boundary states at mirror positions as in equation (A.19). The open string amplitude contains also only the corresponding terms. Furthermore, since the projection operator is now<sup>6</sup><sup>6</sup>6cf equation (A.20) $`4\times \frac{1}{2}(1+g_i(1)^F)`$ each of the terms that appears is sixteen times as large as in the fractional case above. This implies that the relevant normalisation constants are given as
$`𝒩_{t1,U}^2`$ $`=`$ $`{\displaystyle \frac{V_{r+1}}{(2\pi )^{r+1}}}{\displaystyle \frac{1}{8}}{\displaystyle \frac{\underset{m=1}{\overset{s_1+s_2+s_3}{}}R_{j_m}}{_{m=1}^{6s_1s_2s_3}R_{k_m}}},`$ (A.34)
$`𝒩_{t1,Tg_i}^2`$ $`=`$ $`{\displaystyle \frac{V_{r+1}}{(2\pi )^{r+1}}}{\displaystyle \frac{2^{4(s_1+s_2+s_3s_i)}}{8}}{\displaystyle \frac{\underset{m=1}{\overset{s_i}{}}R_{j_m}}{_{m=1}^{2s_i}R_{k_m}}}.`$ (A.35)
The second type of truncated $`\widehat{\text{D}}`$-brane (equation (2.4)) contains the untwisted NS-NS , twisted R-R,T$`g_i`$, R-R,T$`g_j`$ and NS-NS,T$`g_k`$ boundary states. The projection operator is $`\frac{1}{4}(1+g_i(1)^F)(1+g_j(1)^F)`$ thus fixing the normalisations to be
$`𝒩_{t2,U}^2`$ $`=`$ $`{\displaystyle \frac{V_{r+1}}{(2\pi )^{r+1}}}{\displaystyle \frac{1}{64}}{\displaystyle \frac{\underset{m=1}{\overset{s_1+s_2+s_3}{}}R_{j_m}}{_{m=1}^{6s_1s_2s_3}R_{k_m}}},`$ (A.36)
$`𝒩_{t2,Tg_\alpha }^2`$ $`=`$ $`{\displaystyle \frac{V_{r+1}}{(2\pi )^{r+1}}}{\displaystyle \frac{2^{4(s_1+s_2+s_3s_\alpha )}}{64}}{\displaystyle \frac{\underset{m=1}{\overset{s_\alpha }{}}R_{j_m}}{_{m=1}^{2s_\alpha }R_{k_m}}},`$ (A.37)
where $`\alpha \{i,j,k\}`$.
## Appendix B Consistency conditions of boundary states
In this appendix we discuss the invariance of the boundary states under the GSO and orbifold actions. Since the action of the orbifold group is purely geometrical the GSO projection is the same in twisted and untwisted sectors namely we have
$`\text{NS-NS }{\displaystyle \frac{1}{4}}(1+(1)^F)(1+(1)^{\stackrel{~}{F}}),`$ (B.1)
$`\text{R-R }{\displaystyle \frac{1}{4}}(1+(1)^F)(1(1)^{\stackrel{~}{F}}),`$ (B.2)
for Type IIA and IIB, respectively.
The GSO invariance of each of the sectors’ boundary states was computed in detail in . Here we simply restate those results in terms of $`r`$ and $`s_i`$. We shall work in the light-cone gauge with the directions $`x^0,x^9`$ always Dirichlet . In the untwisted NS-NS sector it is easy to see that
$$|B(r;𝐬)_{\text{NS-NS }}=\frac{𝒩}{2}(|B(r;𝐬),+_{\text{NS-NS }}|B(r;𝐬),_{\text{NS-NS }})$$
(B.3)
is GSO invariant for all $`r,s_i`$. The exact values of the normalisation constants of the boundary states will depend on the kind of D-brane the boundary state is a part of, and are computed in the Appendix. In the untwisted R-R sector
$$|B(r;𝐬)_{\text{R-R }}=\frac{4i𝒩}{2}(|B(r;𝐬),+_{\text{R-R }}+|B(r;𝐬),_{\text{R-R }})$$
(B.4)
is a GSO invariant boundary state if $`r+s_1+s_2+s_3`$ is even/odd for Type IIA/B, respectively. For the NS-NS,T$`g_1`$ sector we find that
$$|B(r;𝐬)_{\text{NS-NS,T}g_1}=𝒩_T(|B(r;𝐬),+_{\text{NS-NS,T}g_1}+|B(r;𝐬),_{\text{NS-NS,T}g_1})$$
(B.5)
is a GSO invariant boundary state provided $`s_2+s_3`$ is even, while
$$|B(r;𝐬)_{\text{R-R,T}g_1}=i𝒩_T(|B(r;𝐬),+_{\text{R-R,T}g_1}+|B(r;𝐬),_{\text{R-R,T}g_1})$$
(B.6)
is GSO invariant for $`r+s_1`$ even/odd for Type IIA/IIB. Similarly in the NS-NS,T$`g_2`$ sector $`s_1+s_3`$ is to be even and for the R-R,T$`g_2`$ sector $`r+s_2`$ has to be even/odd for Type IIA/IIB. Finally for the NS-NS,T$`g_3`$ $`s_1+s_2`$ is to be even and for the R-R,T$`g_2`$ sector $`r+s_3`$ has to be even/odd for Type IIA/IIB.
Next one requires the boundary states to be invariant under $`g_i`$. As usual this places no restrictions on the untwisted NS-NS sector, and in the untwisted R-R sector it requires for $`s_1+s_2`$, $`s_2+s_3`$ and $`s_1+s_3`$ to be all even. The R-R,T$`g_i`$ boundary state has no restrictions placed on it by $`g_i`$ since it has no zero modes in the directions which $`g_i`$ inverts, while the NS-NS,T$`g_i`$ boundary state’s invariance under $`g_i`$ is equivalent to the GSO condition on that boundary state .
New non-trivial restrictions arise by requiring the $`g_i`$-twisted sector’s boundary state be invariant under $`g_j`$, $`ji`$. We consider the invariance of NS-NS,T$`g_1`$ and R-R,T$`g_1`$ under $`g_2`$ in some detail. The other conditions will follow in a similar way. Since the NS-NS,T$`g_1`$ sector has zero modes in the directions $`x^7`$ and $`x^8`$ (as well as $`x^5`$ and $`x^6`$) $`g_2`$ will have a non-trivial representation on these zero-modes, given by<sup>7</sup><sup>7</sup>7We use the same convention here for $`g_2`$ as was used for $`_2`$ in . Since this is a supersymmetric theory this can be viewed as a condition for supersymmetry preservation by the orbifold action.
$$g_2=\underset{\mu =7,8}{}(\sqrt{2}\psi _0^\mu )\underset{\mu =7,8}{}(\sqrt{2}\stackrel{~}{\psi }_0^\mu ).$$
(B.7)
This operator squares to one. It is not difficult to see that
$$g_2|B(r;𝐬),\eta _{\text{NS-NS,T}g_1}^0=(1)^{s_3}|B(r;𝐬),\eta _{\text{NS-NS,T}g_1}^0$$
(B.8)
and hence that $`s_3`$ has to be even.<sup>8</sup><sup>8</sup>8Note that in the notation of equations (B.4) and (B.5) Appendix B in for the NS-NS,T$`g_1`$ state $`a=b=1`$. Similarly since the R-R,T$`g_2`$ sector has zero modes in directions $`x^3`$ and $`x^4`$, $`g_2`$ has a non-trivial representation on this sector as
$$g_2=\underset{\mu =3,4}{}(\sqrt{2}\psi _0^\mu )\underset{\mu =3,4}{}(\sqrt{2}\stackrel{~}{\psi }_0^\mu ),$$
(B.9)
and so
$$g_2|B(r;𝐬),\eta _{\text{R-R,T}g_1}^0=(1)^{s_1}|B(r;𝐬),\eta _{\text{R-R,T}g_1}^0.$$
(B.10)
Thus $`s_1`$ has to be even.<sup>9</sup><sup>9</sup>9Following equations (B.12) and (B.13) of $`\widehat{a}=\widehat{b}=1`$ for the R-R,T$`g_1`$ ground state. Performing a similar analysis for the other twisted sectors one finds that in Type IIB the following boundary states are GSO and orbifold invariant
$`|B(r;𝐬)_{\text{NS-NS}}`$ for all $`r`$ and $`s_i`$,
$`|B(r;𝐬)_{\text{R-R}}`$ for either $`r`$ odd and $`s_i`$ all even or $`r`$ even and $`s_i`$ all odd,
$`|B(r;𝐬)_{\text{NS-NS,T}g_i}`$ for $`s_j,s_k`$ even, $`j,ki`$,
$`|B(r;𝐬)_{\text{R-R,T}g_i}`$ for $`r`$ odd and $`s_i`$ even.
For Type IIA the conditions are the same for $`s_i`$ but $`r`$ has to be even.
## Acknowledgments
I am grateful to D.-E. Diaconescu, M.B. Green, F.Quevedo, A. Sen, B. Totaro and in particular to M.R. Gaberdiel and G. Segal for many helpful discussions and insights. B.S. is supported by the Cambridge Commonwealth Trust.
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# References
The $`\omega \omega `$ channel has been studied in several different production mechanisms. In $`\pi ^{}p`$ interactions the $`\omega \omega `$ final state has been studied by the NA12 and VES collaborations. In both experiments clear signals were observed at 1.6 and 1.9 GeV and were found to have $`J^{PC}`$ = $`2^{++}`$, called the $`f_2(1640)`$ and $`X(1910)`$ . In addition, the VES collaboration reported evidence for an $`\omega \omega `$ decay mode of the $`f_4(2050)`$ and, more recently, for a $`J^{PC}`$ = $`4^{++}`$ object in the 2.3 GeV region . In $`p\overline{p}`$ annihilations C. Baker et al., using the data from the Crystal Barrel experiment , have reported evidence for a structure similar to the $`f_2(1640)`$ in the $`\omega \omega `$ final state but have shown that this state can be interpreted as being due to the $`f_2(1565)`$ previously observed in the $`\pi \pi `$ final state . The PDG lists the $`X(1910)`$ observed in the $`\omega \omega `$ final state with another $`J^{PC}`$ = $`2^{++}`$ resonance with similar mass and width observed in the $`\eta \eta ^{}`$ final state. In central production the WA102 experiment did not observe the $`f_2(1565)`$ in the $`\pi \pi `$ final state , therefore, the centrally produced $`\omega \omega `$ channel can give information on the validity of the $`f_2(1565)/f_2(1640)`$ assignment. In addition, in the $`\eta \eta ^{}`$ final state of the WA102 experiment a peak was observed at 1.9 GeV which was consistent in mass and width with the $`X(1910)`$. A spin analysis showed that this state was consistent with having $`J^{PC}`$ = $`1^+`$ with spin projection $`J_Z`$ = 1 or $`J^{PC}`$ = $`2^{++}`$ with spin projection $`J_Z`$ = 2. If the latter hypothesis were true then this was the first time that a state had been observed in central production that was produced with spin projection $`J_Z`$ = 2. Hence, if the states observed in $`\omega \omega `$ and $`\eta \eta ^{}`$ are the same and the $`X(1910)`$ has $`J^{PC}`$ = $`2^{++}`$, the $`X(1910)`$ should be observed in the $`J_Z`$ = 2 projection in the $`\omega \omega `$ final state. In central production, the $`\omega \omega `$ final state was previously observed by the WA76 experiment but only 80 events were observed and hence no strong conclusions could be drawn.
In this paper, a study is presented of the $`\omega \omega `$ final state formed in the reaction
$$ppp_f(\omega \omega )p_s$$
(1)
at 450 GeV/c. It represents more than a factor of 60 increase in statistics over previous data on the centrally produced $`\omega \omega `$ final state and, moreover, will present a spin analysis of this channel in central production. The data come from the WA102 experiment which has been performed using the CERN Omega Spectrometer, the layout of which is described in ref. . Reaction (1) has been isolated using the $`\pi ^+\pi ^{}\pi ^0`$ decay mode of both $`\omega `$s. The reaction
$$ppp_f(\pi ^+\pi ^{}\pi ^+\pi ^{}\pi ^0\pi ^0)p_s$$
has been isolated from the sample of events having six outgoing charged tracks and four $`\gamma `$s reconstructed in the GAMS-4000 calorimeter, by first imposing the following cuts on the components of the missing momentum: $`|`$missing $`P_x|<17.0`$ GeV/c, $`|`$missing $`P_y|<0.16`$ GeV/c and $`|`$missing $`P_z|<0.12`$ GeV/c, where the $`x`$ axis is along the beam direction. The two photon mass spectrum, when the mass of the other $`2\gamma `$-pair lies within a band around the $`\pi ^0`$ mass (100–170 MeV), shows a clear $`\pi ^0`$ signal with small background. Events containing a fast $`\mathrm{\Delta }^{++}(1232)`$ were removed if $`M(p_f\pi ^+)<1.3`$ GeV, which left 294 463 centrally produced $`\pi ^+\pi ^{}\pi ^+\pi ^{}\pi ^0\pi ^0`$ events.
Fig. 1a) shows a lego plot of $`M(\pi ^+\pi ^{}\pi ^0)`$ versus $`M(\pi ^+\pi ^{}\pi ^0)`$ (four combinations per event). A clear signal of the $`\omega \omega `$ channel can be observed. Fig. 1b) shows the $`\pi ^+\pi ^{}\pi ^0`$ mass spectrum if the other $`\pi ^+\pi ^{}\pi ^0`$ combination is compatible with being an $`\omega `$ (0.76 $``$ M($`\pi ^+\pi ^{}\pi ^0`$ $``$ 0.81 GeV) where a clear $`\omega `$ signal can be observed. A tight cut has been used around the $`\omega `$ signal to increase the signal to background ratio in the selected sample. In order to decrease the background further the parameter $`\lambda `$ is introduced which describes the $`\omega `$ decay on the Dalitz plot and is defined as:
$$\lambda =\frac{|\stackrel{}{p}_+\times \stackrel{}{p}_{}|^2}{\frac{3}{4}(\frac{1}{9}m^2m_\pi ^2)^2}$$
where $`|\stackrel{}{p}_+\times \stackrel{}{p}_{}|`$ is proportional to the decay matrix element for $`\omega \pi ^+\pi ^{}\pi ^0`$, $`\stackrel{}{p}_\pm `$ is the three momentum of the $`\pi ^\pm `$ in the $`\omega `$ rest frame and $`m^2`$ is the $`\pi ^+\pi ^{}\pi ^0`$ effective mass squared. Superimposed on fig. 1b) as a shaded histogram is the $`\pi ^+\pi ^{}\pi ^0`$ mass distribution for $`\lambda `$ $`>`$ 0.3. As can be seen the signal to background ratio in the $`\omega `$ region has increased.
The $`\omega \omega `$ final state has been selected using the $`\pi ^+\pi ^{}\pi ^0`$ mass cuts described above and by requiring that the $`\lambda `$ $`>`$ 0.3 for each $`\omega `$ candidate. The resulting $`\omega \omega `$ mass spectrum is shown in fig. 1c) and consists of 5067 events. As can be seen there is a peak in the 1.9 GeV region.
The background below the $`\omega `$ signal has several sources including combinatorics and other channels. The combinatorial background is removed, in part, in the selection procedure. The remaining background is approximately 27 %. Four methods have been used to determine the effects of this background; studying the side bands around the $`\omega `$ signal, studying events that do not balance momentum, studying events that do not pass the $`\lambda `$ selection cuts and studying events from the $`\pi ^+\pi ^{}\pi ^+\pi ^{}\pi ^+\pi ^{}`$ channel. Since the majority of the background is due to other physical channels for example $`a_1(1260)a_1(1260)`$ or $`\omega a_1(1260)`$ production, the two methods that best reproduce the background are the one using events that do not pass the $`\lambda `$ cut and the other uses events from the $`\pi ^+\pi ^{}\pi ^+\pi ^{}\pi ^+\pi ^{}`$ channel. These two methods give a very similar representation of the background. In the remainder of this paper the method used to determine the background will be the mean of these two methods. Superimposed on the $`\omega \omega `$ mass spectrum in fig 1c) as a shaded histogram is the estimate of the background.
A spin analysis of the centrally produced $`\omega \omega `$ system has been performed using the method described in ref. for the $`\rho \rho `$ final state modified for the $`\omega \omega `$ channel. The z axis is defined by the momentum vector of the exchanged particle with the greatest four-momentum transferred in the $`\omega \omega `$ centre of mass. Assuming that only angular momenta up to 4 contribute, the amplitudes have been calculated in the spin-orbit (LS) scheme using spherical harmonics.
In order to perform a spin parity analysis the log likelihood function, $`_j=_i\mathrm{log}P_j(i)`$, is defined by combining the probabilities of all events in 50 MeV $`\omega \omega `$ mass bins from 1.5 to 3.0 GeV. The incoherent sum of various event fractions $`a_j`$ is calculated so as to include more than one wave in the fit,
$$=\underset{i}{}\mathrm{log}\left(\underset{j}{}a_jP_j(i)+(1\underset{j}{}a_j)\right)$$
(2)
where the term $`(1_ja_j)`$ represents the phase space background. The negative log likelihood function ($``$) is then minimised using MINUIT . Coherence between different $`J^P`$ states has been neglected in the fit. Different combinations of waves have been tried and insignificant contributions have been removed from the final fit.
It is found necessary to introduce the $`J^{PC}`$ = $`2^{++}`$ wave with both $`J_Z`$ = 0 and 2, the $`J^{PC}`$ = $`0^{++}`$ wave and the $`J^{PC}`$ = $`4^{++}`$ wave with $`J_Z`$ = 1. The results of the best fit are shown in fig. 2.
The $`J^{PC}`$ = $`2^{++}`$ wave with $`J_Z`$ = 2 shows a peak at 1.9 GeV. This wave has been fitted using a spin 2 relativistic Breit-Wigner and a linear background and is shown superimposed. The fit gives M = 1897 $`\pm `$ 11 MeV, $`\mathrm{\Gamma }`$ = 202 $`\pm `$ 32 MeV, parameters consistent with those of the $`X(1910)`$ found from a fit to the $`\eta \eta ^{}`$ final state . Hence this consistent with the fact that the $`X(1910)`$ has $`\omega \omega `$ and $`\eta \eta ^{}`$ decay modes, and we shall refer to it as the $`f_2(1910)`$ hereafter. Correcting for the unseen decay modes and the effects of the detector, the branching ratio $`\omega \omega `$ /$`\eta \eta ^{}`$ of the $`f_2(1910)`$ is 2.6 $`\pm `$ 0.6. There was no evidence for any wave with $`J_Z`$ = 2 in the $`\eta \eta `$ final state of the WA102 experiment and hence an upper limit for the branching ratio $`\eta \eta `$ /$`\eta \eta ^{}`$ of the $`f_2(1910)`$ has been calculated to be $`<`$ 0.2 ($`90\%`$ CL).
The $`J^{PC}`$ = $`2^{++}`$ wave with $`J_Z`$ = 0 shows a broad enhancement. Superimposed on the wave is a shaded histogram representing the $`f_2(1640)`$. As can be seen the $`J^{PC}`$ = $`2^{++}`$ wave with $`J_Z`$ = 0 is not compatible with the $`f_2(1640)`$ observed by other experiments. This non observation does not contradict the claim that the $`f_2(1640)`$ is an $`\omega \omega `$ decay mode of the $`f_2(1565)`$ since this state is also not observed in central production.
The $`J^{PC}`$ = $`0^{++}`$ wave shows some activity near threshold and a broad enhancement around 2 GeV. In a previous analysis of the 4$`\pi `$ channel, the WA102 experiment observed a similar structure in the $`J^{PC}`$ = $`0^{++}`$ $`\rho \rho `$ wave which was identified with the $`f_0(2000)`$. Superimposed on the wave is a shaded histogram representing the $`f_0(2000)`$ assuming that the branching ratio $`\rho \rho `$/$`\omega \omega `$ = 3 as expected for a isoscalar resonance. This well represents the wave in the 2 GeV region.
The $`J^{PC}`$ = $`4^{++}`$ wave with $`J_Z`$ = 1 shows no evidence for the $`f_4(2050)`$ but does show a peak at 2.3 GeV. The change in log likelihood in the three 50 MeV bins around the 2.3 GeV peak produced by introducing the $`J^{PC}`$ = $`4^{++}`$ wave with $`J_Z`$ = 1 is $`\mathrm{\Delta }=28`$. This is the first time that it has been found necessary to introduce any wave with $`J`$ $`>`$ 2 in the WA102 data. This wave has been fitted using a spin 4 relativistic Breit-Wigner and a linear background and is superimposed on the wave. The fit gives M = 2332 $`\pm `$ 15 MeV, $`\mathrm{\Gamma }`$ = 260 $`\pm `$ 57 MeV parameters consistent with those found by the VES experiment . This state is most likely the $`f_4(2300)`$ observed previously in other experiments and we shall refer to it as so hereafter.
States that have a decay to $`\omega \omega `$ might also be expected to have a decay to $`\rho \rho `$ . As was observed above there appears to be evidence for an $`\omega \omega `$ decay of the $`f_0(2000)`$ previously observed in the $`\rho \rho `$ final state. In the previous analysis of the 4$`\pi `$ final state , no evidence was claimed for either a $`J^{PC}`$ = $`2^{++}`$ $`\rho \rho `$ wave with $`J_Z`$ = 2 or a $`J^{PC}`$ = $`4^{++}`$ $`\rho \rho `$ wave with $`J_Z`$ = 1. Because of the large number of possible waves in the 4$`\pi `$ final state ( $``$ 180 for $`J`$ $``$ 2) only waves that changed the log likelihood by more than 100 were considered. The $`J^{PC}`$ = $`2^{++}`$ $`\rho \rho `$ wave with $`J_Z`$ = 2 was rejected because it changed the likelihood by $``$ 60. The $`J^{PC}`$ = $`4^{++}`$ $`\rho \rho `$ wave with $`J_Z`$ = 1 was not considered because only waves with $`J`$ $``$ 2 were included in the fit. If the $`J^{PC}`$ = $`2^{++}`$ $`\rho \rho `$ wave with $`J_Z`$ = 2 and the $`J^{PC}`$ = $`4^{++}`$ $`\rho \rho `$ wave with $`J_Z`$ = 1 are both introduced into the fit of the $`\pi ^+\pi ^{}\pi ^+\pi ^{}`$ channel then the log likelihood increases by 58 units in the region of the $`f_2(1910)`$ and 27 units in the $`f_4(2300)`$ region. If the signal in the $`J^{PC}`$ = $`2^{++}`$ $`\rho \rho `$ wave with $`J_Z`$ = 2 is interpreted as being due to the $`f_2(1910)`$ then after correcting for the unseen decay modes and the effects of the detector the branching ratio $`\rho \rho `$ /$`\omega \omega `$ of the $`f_2(1910)`$ is 2.6 $`\pm `$ 0.4 consistent with it being at isoscalar resonance. Similarly if the signal in the $`J^{PC}`$ = $`4^{++}`$ $`\rho \rho `$ wave with $`J_Z`$ = 1 is interpreted as being due the $`f_4(2300)`$ then after correcting for the unseen decay modes and the effects of the detector the branching ratio $`\rho \rho `$ /$`\omega \omega `$ of the $`f_4(2300)`$ is 2.8 $`\pm `$ 0.5.
In previous analyses a study has been made of how different resonances are produced as a function of the parameter $`dP_T`$, which is the difference in the transverse momentum vectors of the two exchange particles , and as a function of the azimuthal angle $`\varphi `$ which is defined as the angle between the $`p_T`$ vectors of the two outgoing protons. A study of the background subtracted $`\omega \omega `$ system over the whole mass range as a function of $`dP_T`$ has been performed. The fraction of all $`\omega \omega `$ production has been calculated for $`dP_T`$$``$0.2 GeV, 0.2$``$$`dP_T`$$``$0.5 GeV and $`dP_T`$$``$0.5 GeV and gives 0.12 $`\pm `$ 0.02, 0.36 $`\pm `$ 0.02 and 0.52 $`\pm `$ 0.02 respectively. This results in a ratio of production at small $`dP_T`$ to large $`dP_T`$ of 0.23 $`\pm `$ 0.04. This ratio is much lower than has been observed in the $`K^{}(892)\overline{K}^{}(892)`$ and $`\varphi \varphi `$ final states. However, the latter final states have been shown to be dominantly due to the $`f_2(1950)`$ which is produced mainly at small $`dP_T`$ .
The amount of $`f_2(1910)`$ has also been determined in the same $`dP_T`$ intervals and gives 0.20 $`\pm `$ 0.04, 0.62 $`\pm `$ 0.07 and 0.18 $`\pm `$ 0.04 respectively. This results in a ratio of production at small $`dP_T`$ to large $`dP_T`$ of 1.1 $`\pm `$ 0.3. This value is consistent with what has been observed for the glueball candidates the $`f_0(1500)`$, $`f_0(1710)`$ and $`f_2(1950)`$ .
The azimuthal angle ($`\varphi `$) between the $`p_T`$ vectors of the two protons is shown in fig. 3a) for the background subtracted $`\omega \omega `$ channel for the entire mass range and in fig. 3b) for the $`f_2(1910)`$. The distribution for the $`f_2(1910)`$ is similar to that observed for other glueball candidates .
In summary, a spin analysis of the $`\omega \omega `$ channel has been performed for the first time in central production. Evidence is found for the $`f_2(1910)`$ in the $`J^{PC}`$ = $`2^{++}`$ wave with spin projection $`J_Z`$ = 2. This is the only state observed in central production with spin projection $`J_Z`$ = 2. Its $`dP_T`$ and $`\varphi `$ dependencies are similar to those observed for other glueball candidates. In addition, evidence is found for a state with $`J^{PC}`$ = $`4^{++}`$ consistent with the $`f_4(2300)`$. The $`f_0(2000)`$, previously observed in the $`\rho \rho `$ final state, is confirmed.
Acknowledgements
This work is supported, in part, by grants from the British Particle Physics and Astronomy Research Council, the British Royal Society, the Ministry of Education, Science, Sports and Culture of Japan (grants no. 07044098 and 1004100), the French Programme International de Cooperation Scientifique (grant no. 576) and the Russian Foundation for Basic Research (grants 96-15-96633 and 98-02-22032).
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# Geometric chaoticity leads to ordered spectra for randomly interacting fermions
## Abstract
A rotationally invariant random interaction ensemble was realized in a single-$`j`$ fermion model. The dominance of ground states with zero and maximum spin was confirmed and explained with a statistical approach based on the random coupling of individual angular momenta. The interpretation is supported by the structure of the ground state wave functions.
The interplay of regular and chaotic features in many-body quantum dynamics is currently extensively studied both for simple models and for realistic applications to atomic , nuclear , and condensed matter physics , as well as for understanding properties of the QCD vacuum . Typical finite “shell-model” systems such as complex atoms and nuclei are described by the mean field and corresponding residual interaction. The density of the mean field configurations grows exponentially for combinatorial reasons, so that the interaction becomes effectively strong at sufficiently high excitation energy leading to generic chaotic features both in spectral statistics, which rapidly move to the limit of random matrix theory , and in properties of wave functions . Studies of finite many-body systems have to account for the existence of constants of motion such as total angular momentum, isospin and parity. If these conservation laws are exact, one usually deals with the states of each class separately. However, little attention was paid to the problem of correlations between classes of states which are described by the same Hamiltonian but belong to different values of exact integrals of motion.
An obvious and practically important example is angular momentum conservation in a finite Fermi-system. The prediagonalization procedure of projecting the correct value $`J`$ of nuclear spin out of the $`m`$-scheme Slater determinants induces by itself a strong mixing of the states within a shell model configuration . The projected states of various spins acquire a nearly uniform degree of complexity and energy dispersion. For a sufficiently large dimension, the majority of states correspond to a complicated quasi-random coupling of individual spins. This “geometric chaoticity” was used long ago in evaluating the level density for a given $`J`$. It also plays an important role in the response to external fields, large amplitude collective motion, dissipation and so on . The similarity of different $`J`$-classes with respect to mixing was demonstrated in the nuclear shell model by the studies of complexity, occupation numbers, strength functions and pairing properties. This raises also a question of existence of compound rotational bands which would connect complicated states having different $`J`$ but almost the same mixing.
A new angle of looking at the problem was introduced by refs. where the spectrum of a random but rotationally invariant Hamiltonian was obtained for a shell-model Fermi system. In spite of the random character of the two-body interaction, the fraction $`f_0`$ of the ensemble realizations with a ground state spin $`J_0=0`$ was much higher than the total statistical fraction $`f_0^s`$ of $`J=0`$ states in shell-model space. This result was confirmed in refs. , as well as for the interacting boson model . A new feature discovered in was an excess of the probability $`f_{J_{max}}`$ for the ground state to have the maximum possible spin $`J_{max}`$. The emergence of regular features as an output of a random interaction seems to contradict the notion of geometrical chaoticity. Below we show that, vice versa, the geometric chaoticity provides a base for explaining the main features of the pattern.
First we give a couple of trivial examples which point out the possible source of the effects, namely an analog of the Hund rule in atomic physics. Consider a system of $`N`$ pairwise interacting spins with the Hamiltonian
$$H=A\underset{ab}{}𝐬_a𝐬_b=A[𝐒^2Ns(s+1)].$$
(1)
If the interaction strength $`A`$ is a random variable with zero mean, then the ground state of the system will have equal, $`f_0=f_{S_{max}}=1/2`$, probabilities to have spin $`S=0`$ or $`S=S_{max}`$ (antiferromagnetism or ferromagnetism). A similar situation takes place in the degenerate pairing model where the pair creation, $`P_0^{}`$, pair annihilation, $`P_0`$, and particle number, $`N`$, operators form an SU(2) pseudospin algebra. Then the eigenenergy is simply proportional to the pairing constant so that, for a random sign of this constant, the ground state pseudospin will be 0 (unpaired state of maximum seniority) or maximum possible (fully paired state of zero seniority), on average in 50% of cases. In the Elliott SU(3) model , as well as in any model with a rotational spectrum, the normal or inverted bands will happen evenly if the moment of inertia takes positive or negative values randomly.
Let us consider a system of interacting fermions. For simplicity we limit ourselves here to a case of $`N`$ identical particles on a single-$`j`$ shell which provides a generic framework for the extreme limit of strong residual interaction. Rotational invariance is preserved, so that all single-particle $`m`$-states are degenerate in energy. Within this space, the general two-fermion rotationally invariant interaction can be written as
$$H=\underset{L\mathrm{\Lambda }}{}V_LP_{L\mathrm{\Lambda }}^{}P_{L\mathrm{\Lambda }},$$
(2)
where the pair operators with pair spin $`L`$ and its projection $`\mathrm{\Lambda }`$ are defined as
$$P_{L\mathrm{\Lambda }}^{}=\frac{1}{\sqrt{2}}\underset{mn}{}C_{mn}^{L\mathrm{\Lambda }}a_m^{}a_n^{},P_{L\mathrm{\Lambda }}=\frac{1}{\sqrt{2}}\underset{mn}{}C_{mn}^{L\mathrm{\Lambda }}a_na_m;$$
(3)
and $`C`$ are the Clebsch-Gordan coefficients. Because of Fermi statistics, only even $`L`$ values are allowed in the single-$`j`$ space. This fact was ignored in the attempt to construct the quasiparticle ensemble with identical distributions of the parameters $`V_L`$ in the particle-particle channel and the parameters $`\stackrel{~}{V}_K`$ for the same interaction transformed to the particle-hole channel, $`H_{K\kappa }\stackrel{~}{V}_K(a^{}a)_{K\kappa }(a^{}a)_{K\stackrel{~}{\kappa }}`$ (the difference between the interactions in the two channels was discussed long ago by Belyaev , and served as a justification for an interpolating model “pairing plus multipole-multipole forces”). Since $`K`$ can take both even and odd values, the number of parameters is different in the two representations, and $`\stackrel{~}{V}_K`$ cannot be independent if $`V_L`$ are.
Assuming that the coupling constants $`V_L`$ are random, uncorrelated and uniformly distributed between -1 and 1, we get the distribution $`f_J`$ of the ground state spin $`J_0`$ shown in Fig. 1(a–e) for $`N=4`$ and $`N=6`$ at different values of $`j`$. For comparison we show by dotted lines the statistical distributions $`f_J^s`$ based on the fraction of states of given $`J`$ in the entire Hilbert space for given $`N`$. The overwhelming probability $`f_0`$ shows the same phenomenon in the uniform ensemble as observed earlier in Gaussian ensembles of $`V_L`$ . Further evidence of the dominance of $`J_0=0`$ configurations is given by the example, Fig. 1 (e) , for an odd number of particles, where excess of the ground state spin $`J_0=j`$ is evidently related to the ground spin $`J_0=0`$ in the neighboring even system.
First we note that the effect seems to exist already in a crude approximation modeling fermionic pairs by bosons. The commutation relations for the fermion pair operators (3) are ($`L`$ and $`L^{}`$ are even),
$$[P_{L^{}\mathrm{\Lambda }^{}},P_{L\mathrm{\Lambda }}^{}]=\delta _{L^{}L}\delta _{\mathrm{\Lambda }^{}\mathrm{\Lambda }}+2\underset{mm^{}n}{}C_{m^{}n}^{L^{}\mathrm{\Lambda }^{}}C_{nm}^{L\mathrm{\Lambda }}a_m^{}a_m^{}.$$
(4)
The second term in (4) is of the order $`N/\mathrm{\Omega }`$ where $`\mathrm{\Omega }`$ is the capacity ($`=2j+1`$ in our case) of the fermionic orbitals. It is small for a small number of fermions; for a nearly filled shell its effect is also small because of the particle-hole symmetry of states. For intermediate shell occupation this term is not small but can be approximately substituted by its mean value (the monopole part with spin $`K=0`$). Then, after a simple renormalization, $`P_{L\mathrm{\Lambda }}`$ become bosonic operators. This is the assumption used in the original boson expansion techniques and later in the interacting boson models: fermionic pairs $`P_{L\mathrm{\Lambda }}`$ are substituted by bosons $`b_{L\mathrm{\Lambda }}`$, and the Hamiltonian (2) becomes a sum of random bosonic energies $`_{L\mathrm{\Lambda }}\omega _Ln_{L\mathrm{\Lambda }}`$. The ground state in each realization corresponds to the condensation of the bosons into the single-boson states $`|L\mathrm{\Lambda })`$ with the lowest value of $`\omega _L`$. For a given $`L`$, the many-boson states with different $`J`$ allowed for the condensate are degenerate, but the value $`L=0`$ is singled out by the obvious fact that for $`\omega _0=`$min all degenerate states have total spin $`J=0`$ while for the minimum boson energy $`\omega _L`$ at $`L0`$ any specific value of $`J`$, including $`J=0`$, appears only in a small fraction of states. If all $`V_L`$ have the same distribution, we expect $`f_0^b1/k`$ where $`k`$ is a number of (equiprobable) values of $`L`$. All other values $`J0`$ appear with small probabilities $`1/k^2`$. This is demonstrated by Fig. 1(f) where the pattern is qualitatively similar to that in Fig. 1(a–e). The bosonic effect gives only a part (decreasing with increasing $`j`$) of the $`J_0=0`$ dominance observed for the fermions. Another argument against the dominance of the bosonic correlations is given in Fig. 1(c). Here we see that after exact elimination of the monopole term ($`V_{L=0}0`$), the picture does not significantly change although the value $`V_0`$ is now the lowest only in a small fraction, $`2^{(k1)}`$, of all cases (when all $`V_{L0}`$ are positive).
In our opinion, the main effect comes from the statistical correlations of the fermions. They resolve the bosonic degeneracy in favor of the $`J=0`$ and $`J=J_{max}`$ ground states. In the strong mixing among nearly degenerate states, the eigenstates emerge as complicated chaotic superpositions. The only constraints left are the conservation laws for the particle number and total spin. The latter can be taken into account by the standard cranking approach . Thus, we model the system by the Fermi-gas in statistical equilibrium with the occupation numbers $`n_m`$ of individual orbitals characterized by the angular momentum projection $`m`$ onto the cranking axis. The presence of the constraints creates a “body-fixed frame” and splits effective quasiparticle energies, although instead of the collective rotation around a perpendicular axis we have here a random coupling of individual spins with the symmetry (cranking) axis being the only direction which is singled out in the system . Under the constraints
$$N=\underset{m}{}n_m,M=\underset{m}{}mn_m,$$
(5)
equilibrium statistical mechanics leads to the Fermi-Dirac distribution
$$n_m=\frac{1}{\mathrm{exp}(\gamma m\mu )+1}$$
(6)
determined by the Lagrange multipliers of the chemical potential $`\mu `$ and cranking frequency $`\gamma `$; in the end the total projection $`M`$ (equivalent to the $`K`$ quantum number for axially deformed nuclei) is identified with the total spin $`J`$.
The quantities $`\mu (N,M)`$ and $`\gamma (N,M)`$ can be found directly from (5). At $`M=0`$ we have $`\gamma =0`$, so that the expansion in powers of $`\gamma `$ allows one to study the most important region around $`M=0`$; the power expansion is sufficient for all $`M`$ except for the edges. With no cranking, one has the uniform distribution of occupancies $`n_m^0=\overline{n}=N/\mathrm{\Omega }`$. With the perturbational cranking, the occupation numbers are
$$n_m=\overline{n}\left[1\gamma m(1\overline{n})+\frac{\gamma ^2}{2}(m^2m^2)(1\overline{n})(12\overline{n})+\mathrm{}\right].$$
(7)
Here $`m^2=(1/\mathrm{\Omega })_mm^2=𝐣^2/3`$, and terms of higher orders are not shown explicitly. The expectation value of energy in our statistical system can be written as
$$H=\underset{L\mathrm{\Lambda }m_1m_2}{}V_L|C_{m_1m_2}^{L\mathrm{\Lambda }}|^2n_{m_1}n_{m_2}.$$
(8)
Neglecting the correlations between the occupation numbers, $`n_{m_1}n_{m_2}n_{m_1}n_{m_2}`$, we use the statistical result (7) and calculate the geometrical sums with the Clebsch-Gordan coefficients. Expressing the parameter $`\gamma `$ in terms of the total spin $`MJ`$, we come to the result including the terms of the second order in $`J^2`$,
$$H_{N,J}=\underset{L}{}(2L+1)V_L[h_0(L)+h_2(L)J^2+h_4(L)J^4],$$
(9)
where
$$h_0(L)=\overline{n}^2,h_2(L)=\frac{3}{2}\frac{𝐋^22𝐣^2}{𝐣^4\mathrm{\Omega }^2},$$
(10)
$$h_4(L)=\frac{9}{40}\frac{(12\overline{n})^2(3𝐋^4+3𝐋^212𝐣^2𝐋^26𝐣^2+8𝐣^4)}{(1\overline{n})^2N^2\mathrm{\Omega }^2𝐣^8}.$$
(11)
$`J_0`$ is determined by the ensemble distributions of $`h_{2,4}=_L(2L+1)V_Lh_{2,4}(L).`$ For all realizations of the random interaction with non-negative $`h_2`$ and positive $`h_4`$, the ground state has spin $`J_0=0`$. If $`h_2>0`$ but $`h_4<0`$, one has a local minimum of energy at $`J=0`$ although there is a possibility to reach the absolute energy minimum at $`J_{max}=(1/2)N(\mathrm{\Omega }N)`$. This will not happen if at $`J=J_{max}`$ we still have $`h_2+J_{max}^2h_4>0`$. Therefore the probability to have the ground spin state equal to zero turns out to be, in this approximation,
$$f_0=_{S(h_2,h_4)}𝑑h_2𝑑h_4𝒫(h_2h_4),$$
(12)
where the region $`S`$ is defined by the conditions $`h_2>0,h_4>(h_2/J_{max}^2)`$. Since $`h_4`$ is small, the result is close to that for the right semi-plane $`h_2>0`$, and $`f_0`$ should be close to 50%. For a Gaussian distribution of the parameters $`V_L`$ with zero mean and variances $`\sigma _L`$, the distribution of the linear combinations $`h_{2,4}`$ is again Gaussian, and the integral over the region $`S`$ in (12) gives for this case
$$f_0=\frac{1}{4}+\frac{1}{2\pi }\mathrm{arctan}\left[\frac{D+A/J_{max}^2}{\sqrt{ABD^2}}\right],$$
(13)
which is close to 1/2. Here we introduced the combinations of geometric factors weighted with the corresponding variances,
$$A=\underset{L}{}(h_2(L))^2\sigma _L^2,D=\underset{L}{}h_2(L)h_4(L)\sigma _L^2,B=\underset{L}{}(h_4(L))^2\sigma _L^2.$$
(14)
The $`\gamma `$-expansion fails for large momenta. However, the states with high $`M`$ can be constructed exactly. For Fig. 2 we used our statistical approach near $`J=0`$ in conjunction with the exact values in the end region $`J=J_{max}`$ to improve the above result for $`f_0`$ and to get an upper bound for $`f_{J_{max}}`$. Thus the statistical approach provides a good estimate for the dominance of $`J=0`$ and $`J=J_{max}`$ in the ground state; more subtle effects such as odd-even staggering should be considered separately.
Although the energy spectra with random two-body interactions bear clear resemble the ordered spectra of pairing forces, the structure of the eigenstates is close to that expected for chaotic dynamics . Fig. 3(b) shows the distribution $`P(x)`$ of the overlaps $`x=|J=0,\mathrm{g}.\mathrm{s}.|0,\mathrm{p}|^2`$ of ground states with spin 0 obtained in the random ensemble with the ground state $`|0,\mathrm{p}`$ for the degenerate pairing model, the latter corresponding to the case of fixed $`V_0=1,V_{L0}=0`$. In the chaotic limit the wave functions are expected to behave as random superpositions of basis states with uncorrelated components $`C`$ uniformly spread over a unit sphere, $`P(C)\delta (C^21)`$. This is equivalent to the distribution of a single component $`P(C_1)(1C_1^2)^{(n3)/2}`$ where $`n`$ is the space dimension. For $`n1`$, the distribution $`P(C_1)`$ is close to Gaussian whereas the overlaps $`x=C_1^2`$ obey the Porter-Thomas distribution. In the case of Fig. 3 ($`N=6`$ particles, $`j=11/2`$) the dimension of the $`J=0`$ space is small, $`n=3`$, so that $`P(C_1)`$ is constant, and we expect $`P(x)1/\sqrt{x}`$, as in the case of the pion multiplicity for the disordered chiral condensate. Another case considered in Fig. 3(a) corresponds to the overlap of the degenerate pairing model ground state with the ground state in the model with $`V_0=1`$,$`V_{L0}`$ random. Of course, here the completely paired state can appear as the ground state even for random strengths in the channels $`L0`$ which gives the peak at the overlap $`x=1`$. But the character of the distribution changes as well becoming effectively two-dimensional: for $`n=2`$, $`P(x)1/\sqrt{x(1x)}`$.
To conclude, we have shown that statistical correlations of fermions in a finite Fermi system with random interactions drive the ground state spin to its minimum or maximum value. This effect is related to the geometrical chaoticity of the random spin coupling of individual particles. This means, that the dominance of $`0^+`$ ground states in even-even nuclei may at least partly come from incoherent interactions rather than solely from coherent pairing. The structure of ground states with an “antiferromagnetic” type ordering, $`J_0=0`$, is compatible with the predictions for chaotic dynamics. Quantitative relations between the effects of geometric chaoticity and pure dynamic effects in finite many-body systems should be an interesting subject for further detailed studies.
The authors wish to acknowledge P. Cejnar whose expertise in the interacting boson model was very helpful. The authors are grateful to G.F. Bertsch, B.A. Brown, V. Cerovski, V.V. Flambaum, M. Horoi, F.M. Izrailev, and D. Kusnezov for constructive discussions. This work was supported by the NSF grant 96-05207.
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# SYNCHROTRON EMISSION FROM HOT ACCRETION FLOWS AND THE COSMIC MICROWAVE BACKGROUND ANISOTROPY
## 1 INTRODUCTION
The upcoming cosmic microwave background (CMB) experiments, e.g. MAP and the Planck Surveyor, will be able determine the primordial anisotropies to an unprecedented level of accuracy. Because of its high sensitivity, excellent angular resolution and wide range of frequencies, Planck in particular, will be extremely sensitive to extragalactic foreground point sources, which provide the major source of uncertainty in the measurement of the intrinsic fluctuations.
Several studies have therefore been carried out to calculate the contribution of point sources to the CMB anisotropies. Much of this work (see Toffolatti et al. 1999a,b; De Zotti et al. 1999; Gawiser & Smoot 1997; Sokasian, Gawiser & Smoot 1998) has dealt with the contribution from radio sources, the number counts of which are determined down to $`\mu `$Jy but only up to frequencies $``$ 8 GHz. These counts are usually extrapolated to the higher frequencies relevant for the CMB experiments. This implies that the available counts are sensitive enough to include the most significant contribution from the “steep” and “flat” spectrum sources (with $`F_\nu \nu ^\alpha `$, and $`\alpha 0`$, such as compact radio galaxies and radio loud quasars), but are missing, or are strongly under-representing, an important contribution from a class of sources with inverted spectra ($`\alpha 0`$; e.g. De Zotti et al. 1999). This is further emphasized by recent observations at 28.5 GHz, which find up to a factor of 7 more sources than predicted from low-frequency surveys (Cooray et al. 1998). Inverted-spectrum sources, such as those discussed here, may peak in the frequency range of a few tens to a few hundreds GHz, and could therefore provide a considerable contribution in the region where the most sensitive CMB experiments are carried out.
GHz Peaked Spectrum (GPS) sources (see O’Dea et al. 1998, Guerra, Haarsma & Partridge 1998) have been recognized to be an important class of inverted-spectrum sources. Their emission is attributed to synchrotron radiation from compact and high density regions often associated with the early stages of the formation of more classical double radio sources (the so called “young source” scenario; Philips & Mutel 1982). However, as pointed out by Toffolatti et al. (1999), there may be another, distinct, class of strongly inverted spectra due to thermal synchrotron emission in hot or advection dominated accretion flows (ADAFs). Unlike the relatively rare and bright GPS sources (peak fluxes of $`110`$ Jy), usually associated with bright active galaxies or quasars at high redshifts, ADAF sources should be common in nearby galaxies and provide the most significant contribution to the emission in the high radio frequencies of the faint ($``$ a few mJy) radio cores observed in such galaxies.
The reason why we consider hot accretion flows to be common in nearby galactic nuclei is that, in recent years, it has become apparent (e.g. Fabian & Rees 1995; Narayan & Yi 1995; Di Matteo et al. 2000 and references therein) that the nuclei of such galaxies, which host the largest black holes known with masses of $`10^810^{10}M_{}`$ (e.g., Magorrian et al. 1998), are remarkably underluminous for the typically expected accretion rates (determined from measuraments of densities and sound speeds of their hot interstellar medium). In particular, it has been shown (e.g. Di Matteo et al. 2000) that the relative quiescence and spectral characteristics of the early-type galactic nuclei can be well-explained if the central black holes accrete via low radiative-efficiency accretion flows or ADAFs (Rees et al. 1982; for a review see, e.g., Narayan, Mahadevan & Quataert 1998). Moreover, it has been proposed (Di Matteo and Allen 1999) that such flows, which also produce significant emission in the $`X`$-ray band, could provide a significant contribution to the cosmic $`X`$-ray background (XRB). Within the context of these models, a significant fraction of the hard number counts in the X-ray energies should arise from sources at low redshift ($`z1`$). This picture is supported by recent deep Chandra observations, which have resolved about 40 per cent of the hard XRB in point sources in bright early-type galaxies (Mushotzky et al. 2000).
The potential contribution of GPS sources to fluctuations in the CMB anisotropy has been discussed by De Zotti et al. (1999). In this Paper, we examine the specific contribution of inverted spectra ADAF sources in the nuclei of early-type galaxies to the CMB anisotropy. We evaluate their foreground contribution to the small-scale cosmic microwave fluctuations in the low-energy channels foreseen for the Plank surveyor mission. These sources, if indeed common in elliptical galaxies, should be much more numerous albeit fainter than the GPS population, and may therefore provide a stronger noise contribution at the small angular scales.
While it is important to assess the potential contribution of advection-dominated sources to the CMB fluctuations, the forthcoming CMB experiments themselves will, for the first time, provide a large statistical sample of objects with inverted radio spectra. Because most of the ADAF emission occurs in the high radio and in the $`X`$-ray band, Planck observations will possibly provide the most powerful test for the presence of ADAFs around supermassive black holes. In particular, such studies will provide strong constraints on the spectral properties of this class of objects, and will help determine how common they are in the nearby Universe. Confirming the presence of these sources would also support the conjecture that they provide a significant contribution to the hard XRB.
## 2 SYNCHROTRON EMISSION FROM ACCRETION FLOWS IN EARLY-TYPE GALAXIES
Radio continuum surveys (at $`\nu `$ 8 GHz) of elliptical and S0 galaxies have shown that the sources in radio–quiet galaxies tend to be extended but with a compact component with relatively flat or slowly rising radio spectra (with typical spectral indexes of 0.3–0.4). Recent VLA studies at high radio frequencies (up to 43 GHz), although carried out only on a limited sample of objects, have shown that all of the observed compact cores have spectra rising up to $`2030`$ GHz (e.g., Di Matteo et al. 1999).
Although the low-frequency radio emission in these galaxies might still have a significant contribution from the scaled-down radio jets also present in these systems, it has been proposed that the high-frequency emission can be easily accounted for if the supermassive black holes in elliptical galaxies are accreting via ADAFs (Fabian & Rees 1995; Mahadevan 1997; Di Matteo et al. 1999). In an ADAF around a supermassive black hole, the majority of the observable emission is in the high radio and X–ray bands. In the high-frequency radio band, the emission results from synchrotron radiation due to the strong magnetic field in the inner parts of the accretion flow. The X-ray emission is due either to bremsstrahlung or inverse Compton scattering. In the thermal plasma of an ADAF, synchrotron emission rises steeply with decreasing frequency. Under most circumstances the emission becomes self–absorbed and gives rise to a black–body spectrum (in the Rayleigh-Jeans limit) below a turnover frequency $`\nu _\mathrm{c}`$. Above this frequency it decays exponentially as expected from a thermal plasma, due to the superposition of cyclotron harmonics.
The spectral models and self-consistent temperature profile calculations for the ADAF flows in the elliptical galaxy cores are described in detail by Di Matteo et al. (1999; 2000b). Recent studies have also shown that large outflows may be important in such low–radiative efficiency accretion flows (Igumenshchev & Abramowicz 1999; Blandford & Begelman 1999; Stone, Pringle & Begelman 1999). This should lead to a suppression of the synchrotron component with respect to the standard ADAF model with no outflow. Flatter density profiles, as expected from strong mass loss, are usually required to explain the high-resolution VLA observations at high-radio frequencies of a number of ellipticals. As in our previous work, therefore, we model the accretion flows by adopting a density profile which satisfies $`\rho R^{3/2+p}`$, where $`R`$ is the radius of the flow and $`0p1`$. Figure 1 shows the synchrotron emission expected from an ADAF around a supermassive black hole with $`M_{BH}10^9M_{}`$, for different values of $`p`$ (effectively for radially dependent mass accretion rates, $`\dot{M}R^p`$). The uppermost curve is the standard ADAF model calculated with the accretion rates in the flows determined from Bondi accretion theory, $`\dot{M}\dot{M}_{\mathrm{Bondi}}M_{\mathrm{BH}}^2n`$(ISM)$`/c_\mathrm{S}`$(ISM); where the ISM density, $`n`$(ISM), and sound speed, $`c_\mathrm{S}`$(ISM), can be determined from X-ray deprojection analysis, and $`M_{\mathrm{BH}}`$ is given in recent studies by Magorrian et al. (1998).
So far, only a limited sample of sources has been observed at high–resolution radio frequencies (up to 43 GHz), such that both the flux and the position of the synchrotron peak could be determined. The shaded region in Figure 1 identifies the range of models that have been shown to best fit the observed fluxes. However, as the sample is not statistically significant, we will also consider various populations of accretion sources which may have different luminosities due to their different density profiles. Note also, from Figure 1, that because of their sharp spectral cut-offs, these sources are expected to contribute only to the Planck low-frequency channels ($`\nu 100`$ GHz), and in particular to the 30 GHz one.
It is worth pointing out that the synchrotron emission, and the position of its peak in an ADAF, is a strong function of many model variables: the emission at the self-absorbed synchrotron peak arises from the inner regions of an accretion flow and scales as $`L_\nu \nu _c^2T`$, where $`\nu _cT^2BT^2\dot{M}^{1/2}M_{\mathrm{BH}}^{1/4}R^{5/4}`$ (Narayan & Yi 1995); $`T`$ is the electron temperature, and $`B`$ the magnetic field strength. Because of the strong dependences on a number of parameters, we will estimate the contribution to the CMB anisotropy by making the most conservative assumptions for the model source parameters. We stress that our analysis is not intended to explore the full range of parameter space available for the accretion flow models. At present, the theoretical and observational uncertainties involved in such calculations are too large to merit such work. However, although schematic, these models may provide a useful guide for the prediction of the confusion noise due to these sources.
## 3 CONTRIBUTION TO CMB FLUCTUATIONS
The contribution to CMB fluctuations from randomly distributed sources has been extensively discussed in the literature (Scheuer 1957, 1974; Condon 1974; Cavaliere & Setti 1976; Franceschini et al. 1989; Tegmark & Efstathiou 1996).
We let $`x=Sf(\theta )`$ be the telescope response to a source of flux $`S`$ located at a distance $`\theta `$ from the beam axis, and let $`R(x)`$ be the mean number of source responses of intensity $`x`$. The fluctuation level generated by randomly distributed sources with a Poisson distribution is given by the second moment $`\sigma `$ of the $`R(x)`$ distribution. If the angular power pattern of the detector, $`f(\theta )`$, is taken to be a Gaussian with full width half maximum $`\theta _0`$, one obtains:
$$\sigma ^2=_0^{x_c}x^2R(x)𝑑x=\pi \theta _0^2I(x_c)$$
(1)
where
$$I(x_c)=_0^{x_c}𝑑xx^2_0^{\mathrm{}}𝑑\psi N\left(\frac{x}{f(\psi )}\right)\mathrm{exp}(4\psi \mathrm{ln2}).$$
(2)
Here $`\psi (\theta /\theta _0)`$, $`N(S_\nu )`$ are the differential source counts per steradian at a given frequency $`\nu `$, and $`x_c=q\sigma `$ is the threshold flux above which sources are considered to be individually detected. In our calculations we adopt the standard value of $`q=5`$.
The rms brightness temperature fluctuation $`(\mathrm{\Delta }T/T)_{\mathrm{rms}}(\mathrm{\Delta }T/T)^2^{1/2}`$ at the frequency $`\nu `$ is related to the confusion standard deviation $`\sigma `$ by
$$\left(\frac{\mathrm{\Delta }T}{T}\right)_{\mathrm{rms}}=\frac{\sigma }{T\omega _b}\left(\frac{B_\nu }{T}\right)^1,$$
(3)
where $`\omega _b=\pi _0^{\mathrm{}}𝑑\theta ^2f(\theta )=\pi (\theta _0/2)^2\mathrm{log}2`$ is the effective beam area, and
$$\frac{B_\nu }{T}=\frac{2k}{c^2}\left(\frac{kT}{h}\right)^2\frac{x^4e^x}{(e^x1)^2}$$
(4)
is the conversion factor from temperature to flux (per steradian). Here $`B_\nu (T)`$ is the Planck function, $`xh\nu /kT`$, and $`T=2.725`$ K (Mather et al. 1999) is the CMB temperature.
We take the number density of our sources to be that of ellipticals, and use the fit<sup>1</sup><sup>1</sup>1The fit was derived from observations of over 1700 galaxies in various magnitude limited samples (the Autofib Redshift Survey). provided by Heyl et al. (1997) for galaxies with $`LL_{}`$,
$$n(z)=n_0(1+z)^{[\gamma _\varphi \gamma _L(\alpha _0+\gamma _\alpha z)]}\mathrm{exp}[1/(1+z)^{\gamma _L}],$$
(5)
where the set of parameters {$`n_0,\gamma _\varphi ,\gamma _L,\alpha _0,\gamma _\alpha `$} is given by Heyl et al. for red and blue ellipticals. These are {$`1.62\times 10^3`$ Mpc<sup>-3</sup>, -6.15, -1.77, -1.05, 1.3} and $`1.87\times 10^3`$ Mpc<sup>-3</sup>, 1.56, -0.35, -1.06, 1.23} for red and blue respectively. We consider the total $`n(z)`$ and add both components. We take a cutoff at $`z1`$ for their contribution. This is consistent with the normalization of the ADAF sources in ellipticals required to explain the hard XRB (Di Matteo & Allen 2000). Notice, however, that our results are rather insensitive to the precise value of this cut-off, as most of the contribution is produced by sources at very low redshift.
Given a population of low-redshift sources with luminosity per unit frequency $`L_\nu `$, the differential number counts are given by $`N(S_\nu )[4\pi (c/H_0)n(z)d_L^2(z)dz/dS_\nu ]_{z=z(S_\nu )}`$, where $`z(S_\nu )`$ is derived by inverting the relation $`S_\nu =L_\nu (1+z)/4\pi d_L^2(z)`$, and $`d_L`$ is the luminosity distance. We adopt a flat cosmology with $`\mathrm{\Omega }_0=1`$, and take $`H_0=3.2\times 10^{17}\mathrm{s}^1h`$, with $`h=0.65`$. The particular choice of cosmology is not relevant to our calculations, again because the bulk of the contribution from these relatively faint objects comes from nearby sources.
For a Poisson distribution of point sources, it is straightforward to compute their angular power spectrum $`C_l`$. If only the contribution from sources with $`S_\nu S_{\mathrm{lim}}`$ is included, this is given by (Tegmark & Efstathiou 1996; Scott & White 1999)
$$C_l(\nu )=\frac{1}{T^2}\left(\frac{B_\nu }{T}\right)^2_0^{\mathrm{S}_{\mathrm{lim}}}S_\nu ^2N(S_\nu )𝑑S_\nu .$$
(6)
## 4 RESULTS
Figure 2 shows the expected temperature fluctuations (Eq. 3) due to synchrotron emission from accretion flows in the nuclei of ellipticals as a function of the angular scale $`\theta _0`$, at $`\nu =30`$ GHz. The solid line corresponds to the case where the luminosities of all sources fall in the region marked by the shaded region of Figure 1. These are ADAF models with a significant amount of outflow and relatively low luminosities. The other lines show the contribution to the fluctuations due to a mixed population of sources containing a fraction $`f`$ of accretion flows with higher synchrotron luminosities<sup>2</sup><sup>2</sup>2 We cannot attempt to model a proper luminosity function for these sources, as the observed sample in Di Matteo et al. (2000b) is too small to allow such modelling. (e.g. as expected from ADAFs with no outflows corresponding to the uppermost curve of Figure 1). Note that even a small fraction of high-luminosity sources can have large effects on the level of fluctuations. This is due to the fact that the source number counts roughly scale as $`N(S)fL^{1.5}`$, and therefore the dependence on the luminosity is stronger than the dependence on $`f`$. Also note that our derived number counts are comparable (if $`f=0`$) to the number counts extrapolated from low-frequency surveys by Toffolatti et al. (1999a), but can be up to a factor of 10 higher if a considerable fraction of higher-luminosity ($`f0.5`$) sources is considered.
Figure 3 shows, in the case of a fraction $`f=0.5`$ of high-luminosity sources in the sample, the expected level of fluctuations if all sources above a flux $`S_{\mathrm{lim}}`$ were individually identified and substracted out from the sample. This is shown for different choices of $`S_{\mathrm{lim}}`$. Identification and removal of individual sources can in principle be done by means of independent surveys.
Using the same parameters as in Figure 3, Figure 4 shows a comparison between the poissonian power spectrum (Eq. 6) produced by the accretion flows in ellipticals and the predicted power spectrum for a standard CDM model. The heavy solid line shows the contribution to the power spectrum from noise in the 30 GHz channel of the Planck LFI. Note that the “flat” ($`C_l`$ = constant) angular power spectrum of the fluctuations due to ADAF sources differs substantially from the power spectrum of primordial fluctuations. We find that the source signal is generally well below that of the intrinsic fluctuations, and it only becomes comparable to these on the small angular scales, where also the instrumental noise increases to roughly the same level. Removal of source signal should be possible even in cases where it gives a strong contribution. Even if the Poisson component of the sources and the noise due to the instrument have similar power spectra, they are indeed different in their nature (as emphasized by Scott & White 1999). In fact, while the sources on the sky contribute to the flux in every observation of a given pixel, the noise, on the other hand, differs from observation to observation, and, by assumption, it is uncorrelated with the signal in that pixel. Therefore, if a given direction in the sky is observed multiple times (as expected for Planck), the instrumental noise component can be separated from the sky signal.
## 5 DISCUSSION
We have computed the temperature fluctuations and power spectrum produced by inverted radio spectra from hot accretion flows in the nuclei of nearby elliptical galaxies in the Planck 30 GHz channel, where their emission is expected to peak. We have shown that the contribution from this class of sources approaches the intrinsic CMB fluctuation level only at small angular scales. However, because of the different nature of its power spectrum, the source contribution should not affect the most important goal of the Planck mission, that is the accurate measurement of the primary CMB anisotropy.
On the other hand, Planck will provide a large statistical sample of sources characterized by inverted spectra. Therefore, it should be possible to use this study to determine how common this mode of accretion is in the nearby supermassive black holes. In particular, as most of the contribution from this population is expected to peak at high radio frequencies, Planck should allow us to study their spectral characteristics. In turn, because different spectral distributions and luminosities reflect the shape of the density profiles, CMB experiments could allow us to gain important information on the physical conditions in these accretion flows. As already noted by Toffolatti et al. (1999b), the implications of such a study could, more generally, be significant as a way of testing the physical processes in the medium surrounding massive black holes, and the evolution of the interstellar medium in galaxies up to moderate redshifts. Even more, it would provide a test for current ideas according to which a significant fraction of the $`X`$-ray background may due to accretion in this regime in early-type galaxies in the local universe (Di Matteo & Allen 1999). Note that such a significant statistical study would be more difficult to carry out with surveys at other wavelengths because of the rapid decline of the ADAF flux, which makes the emission from this type of accretion flows extremely weak in the far infrared and optical bands.
We need to stress that, in principle, the contribution from ADAF sources should be easily disentangled not only from that due to sources with a flat and steep spectrum, but also from that due to GPS sources which also have strongly inverted spectra. GPS sources are typically much brighter (with fluxes typically ranging from a few to 10 Jy) but rarer (usually associated with QSOs) than the expected ADAF sources. The number of GPS sources rapidly decreases with decreasing flux, whereas ADAFs are expected to be much more numerous at faint flux levels. As a result, GPS are only minor contributors to the fluctuations at small angular scales, whereas ADAFs would be mostly significant at these scales. Therefore it should be possible to study these two populations independently.
Note that we have shown the expected temperature fluctuations due to ADAF sources only for the lowest energy channel of Planck. If most of the sources are indeed in the range of luminosities consistent with those observed so far, then this channel is expected to have the largest (possibly major) contribution, due to the high-frequency cutoff in the spectrum of these sources. However, if a substantial population of high-luminosity sources is present, then some contribution should also be present in the other channels of the Planck LFI. The availability of multifrequency data should allow an efficient identification of pixels contaminated by discrete sources. In order to carry out a substraction of the contaminating flux one should therefore take into account that strongly inverted spectra such as those considered here may not be present in most frequency channels but give rise to a strong contamination up to a certain frequency, and then abruptly drop. It should also be pointed out that, contrary to some of the GPS sources for which variability has been observed (e.g. Stanghellini et al. 1998), the radio sources in the hot accretion flows are usually not very variable. A lack of variability is particularly important for a proper removal of sources from the spectral fitting.
We note that ADAFs around massive black holes could also be found in spiral galaxies such as the Galactic nucleus Sgr A. However, even if ADAFs were indeed common in spiral nuclei, their potential contribution to the CMB anisotropy would still be dominated by that from ellipticals. Inferred black hole masses are found to be proportional to the mass of the bulge component of their host galaxy, implying $`M_{BH}10^610^7M_{}`$ for spiral galaxies. As a consequence, the contribution from spirals should be much lower, as the radio flux scales as $`M_{\mathrm{BH}}^{2.53}`$ (Franceschini et al. 1998), and (see §3) their spectrum would peak at frequencies higher than those of elliptical cores and affect higher energy (e.g. sub-mm, mm) channels of CMB experiments. Because of this, given enough sensitivity, the relevance of ADAFs in quiescent spiral nuclei may also be assesed separately by the forthcoming experiments.
Finally, we note that in our analysis we do not take into account the effects of source clustering. Clustering decreases the effective number of objects in randomly placed cells and, consequently, enhances the cell to cell fluctuations. There is indeed evidence that the positions in the sky of a wide variety of extragalactic sources are correlated (Shaver 1988). However, the specific correlation function of our radio–submm sources in early-type galaxies is not well-constrained. The analyses of Toffolatti et al. (1998) have shown that the contribution due to clustering (using the two-point correlation function from sources selected at 5 GHz; Loan, Wall & Lahav 1997) is generally small in comparison with the Poisson term; however, the relative importance of clustering increases if sources are substracted out from the Planck maps down to faint flux levels.
We thank Ramesh Narayan for motivating this work, and Martin White for useful discussions. We also thank the anonimous referee whose comments greatly improved the presentation of this paper. T. D. M. acknowledges support provided by NASA through Chandra Postdoctoral Fellowship grant number PF8-10005 awarded by the Chandra Science Center, which is operated by the Smithsonian Astrophysical Observatory for NASA under contract NAS8-39073.
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# The gerbe of Higgs bundles
### 0. Introduction
The purpose of this work is to describe the (category of) Higgs bundles on a scheme $`X/`$ having a given cameral cover $`\stackrel{~}{X}`$. We show that this category is a $`T_{\stackrel{~}{X}}`$-gerbe, where $`T_{\stackrel{~}{X}}`$ is a certain sheaf of abelian groups on $`X`$, and we describe the class of this gerbe precisely. In particular, it follows that the set of isomorphism classes of Higgs bundles with a fixed cameral cover $`\stackrel{~}{X}`$ is a torsor over the group $`H^1(X,T_{\stackrel{~}{X}})`$, which itself parametrizes $`T_{\stackrel{~}{X}}`$-torsors on $`X`$. This underlying group $`H^1(X,T_{\stackrel{~}{X}})`$ can be described as a generalized Prym variety, whose connected component is either an abelian variety or a degenerate abelian variety.
The hardest part of our work, though, goes into identifying precisely the $`H^1(X,T_{\stackrel{~}{X}})`$-torsor we get, or in other words, identifying the class of the gerbe. This class is surprisingly complicated. One piece of it can be identified as a twist along the ramification divisors of $`\stackrel{~}{X}`$ over $`X`$, and is present for all groups $`G`$. A second piece is a shift which can be present even for unramified covers. While the twist along the ramification expresses properties of the cameral cover, this shift expresses the non-vanishing of a certain group cohomology element - specifically, the extension class $`[N]`$ of the normalizer $`N=N_G(T)`$, which is an element in the cohomology group $`H^2(W,T)`$ of the Weyl group acting on the maximal torus. It vanishes for some groups, such as $`GL(n),PGL(n),SL(2n+1),SO(n)`$, but not for others such as $`SL(2n)`$. Yet a third piece is present only for the groups $`SO(2n+1)`$ (or groups containing them as direct factors); this piece expresses the existence of non-primitive coroots, which amounts to the non vanishing of an element in another cohomology group. We give several examples to illustrate these individual ingredients as well as their combined effect.
Throughout this work, we let $`G`$ be a connected reductive group, and let $`X`$ be a scheme over the complex numbers. A Higgs bundle over $`X`$ is a principal $`G`$-bundle plus some additional data. We describe this additional data next: first for $`G=GL(n)`$, and then for all $`G`$, in subsection 0.1. In the remainder of this introduction we will outline our results 0.2, discuss some examples and applications 0.3, and review some related results in the literature 0.4. The notation we employ is summarized in 0.5.
#### 0.1. Abelianization: Higgs bundles and cameral covers
It is especially easy to spell out the definition when $`G=GL(n)`$. In this case a $`G`$-bundle is the same as a vector bundle $`E`$ over $`X`$, and a Higgs structure on it is a subbundle of commutative associative algebras $`𝐜_X\mathrm{End}_{𝒪_X}(E)`$, which has rank $`n`$ over $`X`$ and such that $`𝐜_X`$ is locally generated by one section. In this case the relative spectrum of $`𝐜_X`$ over $`X`$ is a flat $`n`$-sheeted cover of $`X`$, called the spectral cover corresponding to our Higgs bundle. We will denote it by $`\overline{X}`$.
How can we classify Higgs bundles with a given spectral cover $`\overline{X}`$? The answer is simple: these are in bijection with line bundles on $`\overline{X}`$. Thus, by asking not just for principal $`G`$-bundles, but rather for $`G`$-bundles endowed with a Higgs structure with a fixed spectral cover, we go from a non-abelian problem to an abelian one.
The natural question now is how to extend the above discussion to other reductive groups. It turns out that the notion of an abstract Higgs bundle is quite easy to generalize. Namely, a Higgs bundle is a pair $`(E_G,𝐜_X)`$, where $`E_G`$ is a principal $`G`$-bundle over $`X`$ and $`𝐜_X`$ is a subbundle of the associated bundle of Lie algebras $`𝔤_{E_G}`$, whose fibers are regular centralizers. The precise definition is given in Section 2. Here we only recall that a regular centralizer in the Lie algebra $`𝔤`$ is an abelian subalgebra $`𝐜𝔤`$ which is the centralizer of some regular (but not necessarily semisimple) element $`g𝔤`$. In particular, taking $`g`$ to be regular semisimple, we see that every Cartan subalgebra (i.e. Lie algebra of a maximal torus) is a regular centralizer. In fact, we will see in Section 1 that the set of regular centralizers in $`𝔤`$ is parametrized by an algebraic variety $`\overline{G/N}`$ which is a partial compactification of the parameter space $`G/N`$ for the maximal tori. The simplest Higgs bundles are the unramified ones, i.e. Higgs bundles $`(E_G,𝐜_X)`$ for which all the fibers of $`𝐜_X`$ are maximal tori.
The situation is less transparent with spectral covers. In fact, we do not know a good definition of a spectral cover that would work for any $`G`$ and reproduce for $`GL(n)`$ the old object.
Instead, we use the notion of a cameral cover introduced in . By definition, the latter is a finite flat map $`p:\stackrel{~}{X}X`$ such that the Weyl group $`W`$ of $`G`$ acts on $`\stackrel{~}{X}`$ and certain restrictions on the ramification behaviour are satisfied (cf. Section 2). When $`G=GL(n)`$, we will note below that this notion is different from that of a spectral cover, though equivalent to it.
It turns out that every Higgs bundle determines in a canonical way a cameral cover, so one is led naturally to the problem of classification of Higgs bundles with a given cameral cover. This is the problem we solve in the present paper. Given a cameral cover $`\stackrel{~}{X}`$, we will describe the corresponding Higgs bundles in terms of the “Abelian” data consisting of the maximal torus $`TG`$, the $`W`$-action on $`T`$, and the ramification pattern of $`\stackrel{~}{X}`$ over $`X`$. The “non-Abelian” data involving the group $`G`$ itself is not needed.
#### 0.2. Outline of the results
We formulate the above classification problem in the categorical framework, in terms of the category $`\mathrm{Higgs}_{\stackrel{~}{X}}(X)`$ of Higgs bundles together with an isomorphism between the induced cameral cover and $`\stackrel{~}{X}`$. Our first result shows that this classification problem is indeed abelian.
Namely, starting from $`\stackrel{~}{X}`$ we define a sheaf of abelian groups $`T_{\stackrel{~}{X}}`$. We assert in Theorem 4.4 that $`\mathrm{Higgs}_{\stackrel{~}{X}}(X)`$ is a gerbe over the Picard category of $`T_{\stackrel{~}{X}}`$-torsors. (These notions are reviewed for the reader’s convenience in Section 3.) This result has two immediate consequences.
First, the set of isomorphism classes of objects in our category $`\mathrm{Higgs}_{\stackrel{~}{X}}(X)`$, i.e. the set of isomorphism classes of Higgs bundles with the given cameral cover $`\stackrel{~}{X}`$, if non-empty, carries a simply-transitive action of the abelian group $`H^1(X,T_{\stackrel{~}{X}})`$ (Corollary 4.6), and is therefore non-canonically isomorphic to it. It is thus a generalized Prym variety, cf. : depending on the circumstances, this may appear as a Jacobian of a spectral curve, or as an ordinary Prym, or as various types of Prym-Tyurin varieties , and so on.
The second consequence allows us to determine when Higgs bundles with the given cameral cover $`\stackrel{~}{X}`$ actually exist. This happens if and only if the gerbe is trivial: the cameral cover $`\stackrel{~}{X}`$ determines an obstruction class in $`H^2(X,T_{\stackrel{~}{X}})`$, and Higgs bundles with the given $`\stackrel{~}{X}`$ exist if and only if this class vanishes (Corollary 4.5).
In the above, the sheaf $`T_{\stackrel{~}{X}}`$ is defined in terms of the slightly larger sheaf $`\overline{T}_{\stackrel{~}{X}}`$ (on $`X`$) of $`W`$-equivariant maps $`\stackrel{~}{X}T`$, i.e. $`\overline{T}_{\stackrel{~}{X}}(U):=\mathrm{Mor}_W(\stackrel{~}{U},T)`$, where $`\stackrel{~}{U}`$ is the induced cameral cover of $`U`$. For each positive root $`\alpha :T𝔾_m`$, let $`s_\alpha `$ be the corresponding reflection acting on $`\stackrel{~}{X}`$, and let $`D_X^\alpha \stackrel{~}{X}`$ be its fixed point scheme. Any section $`t`$ of $`\overline{T}_{\stackrel{~}{X}}(U)`$ determines a function $`\alpha t:\stackrel{~}{U}𝔾_m`$ which goes to its own inverse under the reflection $`s_\alpha `$. In particular, its restriction to the ramification locus $`D_X^\alpha `$ equals its inverse, so it equals $`\pm 1`$. The subsheaf $`T_{\stackrel{~}{X}}\overline{T}_{\stackrel{~}{X}}`$ is given by the positive choice:
$$T_{\stackrel{~}{X}}(U):=\{t\overline{T}_{\stackrel{~}{X}}(U)|(\alpha t)|_{D_U^\alpha }=+1\text{for each root}\alpha \}.$$
Although Theorem 4.4 is quite useful, it is not a completely satisfactory result by itself, as it does not describe which $`T_{\stackrel{~}{X}}`$-gerbe we get. Our main result, Theorem 6.4, gives a complete description of the category $`\mathrm{Higgs}_{\stackrel{~}{X}}(X)`$ as the gerbe parametrizing certain “$`R`$-twisted, $`N`$-shifted $`W`$-equivariant $`T`$-bundles on $`\stackrel{~}{X}`$”. The “twist” here is along the ramification divisors, and the “shift” is by the extension class of the normalizer $`N`$.
Our description of this gerbe is based on an explicit description of the underlying Picard category $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}`$ which appears in the statement of Theorem 4.4. An object in this category, i.e. a $`T_{\stackrel{~}{X}}`$-torsor, consists of:
* A (weakly $`W`$-equivariant) $`T`$-bundle $`_0`$ on $`\stackrel{~}{X}`$,
* A group homomorphism $`\gamma _0:N_0\mathrm{Aut}(_0,\stackrel{~}{X}/X)`$, commuting with the projctions to $`W`$, and
* For every simple root $`\alpha _i`$, a trivialization
$$\beta _{i,0}:\alpha _i(_0)|_{D_X^{\alpha _i}}𝒪_{D_X^{\alpha _i}}.$$
The data of $`\gamma _0`$ and $`\beta _0`$ must satisfy some compatibility conditions, which are described in detail in Section 16. (Roughly, these say that the collection $`\beta _0`$ of isomorphisms $`\beta _{i,0}`$ is $`W`$-equivariant, and $`\beta _0,\gamma _0`$ are related by the compatibility condition: $`\gamma _{0|D_X^\alpha }=\stackrel{ˇ}{\alpha }\beta _0.`$) Morphisms in this category are $`T`$-bundle maps that are compatible with the data of $`\gamma _0`$ and $`\beta _0`$.
Our notation here is as follows. An element of the group $`\mathrm{Aut}(_0,\stackrel{~}{X}/X)`$, for a $`T`$-bundle $`_0`$ on $`\stackrel{~}{X}`$, consists of an element $`wW`$ together with an isomorphism $`w^{}(_0)_0`$. The bundle $`_0`$ is weakly $`W`$-equivariant if $`\mathrm{Aut}(_0,\stackrel{~}{X}/X)`$ surjects onto $`W`$, in which case $`\mathrm{Aut}(_0,\stackrel{~}{X}/X)`$ is an extension of $`W`$ by $`\mathrm{Mor}(\stackrel{~}{X},T)`$. Now the semidirect product $`N_0`$ of $`T`$ and $`W`$ induces one such extension, and $`\gamma _0`$ is supposed to induce an isomorphism of this extension with $`\mathrm{Aut}(_0,\stackrel{~}{X}/X)`$. We think of the root $`\alpha `$ as a homomorphism $`T𝔾_m`$, so $`\alpha (_0)`$ is the line bundle associated to $`_0`$ via this homomorphism. Similarly, the coroot $`\stackrel{ˇ}{\alpha }`$ is a homomorphism $`𝔾_mT`$.
In describing our gerbe, we replace each linear feature in the description of $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}`$ by an affine variant. We start with the equivariance: the $`T`$-bundles $`_0`$ were weakly $`W`$-equivariant (which means that $`w^{}(_0)`$ was isomorphic to $`_0`$, for each $`wW`$), and in fact strongly $`W`$-equivariant (which simply means that $`W`$ itself, and hence also the semidirect product $`N_0`$, acted on them).
Our variant of the weakly $`W`$-equivariant $`T`$-bundles $`_0`$ involves $`T`$-bundles $``$ which are $`R`$-twisted weakly $`W`$-equivariant, meaning that now $`w^{}()_X^w`$ is isomorphic to $``$, for each $`wW`$. Here $`_X^w`$ is a $`T`$-bundle on $`\stackrel{~}{X}`$ which encodes the ramification pattern of $`\stackrel{~}{X}`$ over $`X`$. In the simplest case, when $`\stackrel{~}{X}`$ is integral and $`w`$ is the reflection $`s_\alpha `$ corresponding to a simple root $`\alpha `$, we have $`_X^w=_X^\alpha =\stackrel{ˇ}{\alpha }(R_X^\alpha ),`$ where $`R_X^\alpha `$ is the line bundle $`𝒪_{\stackrel{~}{X}}(D_X^\alpha )`$. The precise definitions are given in Section 6.
Next, we need a substitute for the strong equivariance. We replace $`\mathrm{Aut}(_0,\stackrel{~}{X}/X)`$ by the group $`\mathrm{Aut}_R(,\stackrel{~}{X}/X)`$ of isomorphisms $`w^{}()_X^w`$, and the semidirect product $`N_0`$ by the normalizer $`N`$, so we demand that $`\gamma `$ should map $`N`$ to $`\mathrm{Aut}_R(,\stackrel{~}{X}/X)`$.
Finally, $`\beta _i`$ needs to be twisted by the ramification, so it now sends $`\alpha _i()|_{D_X^{\alpha _i}}^{\alpha _i}|_{D_X^{\alpha _i}}`$. One final complication is that $`\beta _i`$ now depends (linearly) on the choice of a lift of $`w_i`$ to an element $`n_iN`$. (This choice of a lift is unnecessary in the linear version, since $`W`$ is a subgroup of $`N_0`$, so the $`w_i`$’s have a canonical lift.)
We can now give an almost complete statement of our main result, Theorem 6.4. It says that a Higgs bundle with given cameral cover $`\stackrel{~}{X}`$ is equivalent to:
* An $`R`$-twisted, weakly $`W`$-equivariant $`T`$-bundle $``$ on $`\stackrel{~}{X}`$,
* A group homomorphism $`\gamma :N\mathrm{Aut}_R(,\stackrel{~}{X}/X)`$, and
* For every simple root $`\alpha _i`$ and lift $`n_iN`$ of the reflection $`s_iW`$ into $`N`$, the data of an isomorphism
$$\beta _i(n_i):\alpha _i()|_{D_X^{\alpha _i}}^{\alpha _i}|_{D_X^{\alpha _i}}.$$
The data of $`\gamma `$ and $`\beta `$ must satisfy several compatibility conditions, which are described in detail in Section 6. (Roughly, these say that the collection $`\beta `$ of isomorphisms $`\beta _i(n_i)`$ is $`N`$-equivariant, and $`\beta ,\gamma `$ are related by the compatibility condition: $`\gamma _{|D_X^\alpha }=\stackrel{ˇ}{\alpha }\beta .`$) In fact, the category $`\mathrm{Higgs}_{\stackrel{~}{X}}(X)`$ is equivalent to the category $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)`$ whose objects are the triples $`(,\gamma ,\beta )`$ as above. Morphisms in this category are again $`T`$-bundle maps that are compatible with the data of $`\gamma `$ and $`\beta `$.
Note that the possible non-triviality of our gerbe can be attributed to three separate causes: the twist along the ramification $`R`$; the shift resulting from non-triviality of the extension class of $`N`$; or the extra complication involved in choosing the $`\beta _i`$. In subsection 6.5 we give a simplified version of our theorem, which avoids this last complication. It applies in all cases except when the group $`G`$ has $`SO(2n+1)`$ as a direct summand.
#### 0.3. Some examples and applications
##### 0.3.1. The unramified case
The cameral cover $`\stackrel{~}{X}X`$ is unramified if and only if the Higgs bundle $`(E_G,𝐜_X)`$ is unramified, i.e. if and only if the bundle of regular centralizers $`𝐜_X`$ is actually a bundle of Cartan subalgebras. In this case the classification (given in ) is easy: specifying a Higgs bundle $`(E_G,𝐜_X)`$ with the unramified cameral cover $`\stackrel{~}{X}`$ is equivalent to giving an $`N`$-bundle $`E_N`$ over $`X`$ together with an identification of the quotient $`E_N/T`$ with $`\stackrel{~}{X}`$. In this case, our $`T`$-bundle $``$ is just $`E_N`$, considered as a $`T`$-bundle over $`E_N/T=\stackrel{~}{X}`$. Since there is no ramification, there is no $`R`$-twist; similarly, there is no $`\beta `$; and $`\mathrm{Aut}_R(,\stackrel{~}{X}/X)`$ is just $`\mathrm{Aut}(E_N,\stackrel{~}{X}/X)`$, which is induced from the extension $`N`$, so $`\gamma `$ is the tautological map.
##### 0.3.2. $`GL(n)`$
Consider first the case of $`G=GL(n)`$. The spectral cover $`\overline{X}`$ is then of degree $`n`$ over $`X`$, while the cameral cover $`\stackrel{~}{X}`$ is of degree $`n!`$. The $`n`$ points of $`\overline{X}`$ above each point $`x`$ of $`X`$ correspond to the $`n`$ simultaneous eigenvectors (in the standard representation) of the corresponding centralizer $`𝐜_x`$, while the $`n!`$ points of $`\stackrel{~}{X}`$ above $`x`$ correspond to the ways of ordering these eigenvectors. In a generic situation, e.g. when the Higgs bundle is unramified or only simply ramified, it is clear that $`\stackrel{~}{X}`$ is precisely the Galois closure of the spectral cover $`\overline{X}`$. Conversely, $`\overline{X}`$ is recovered as the quotient of $`\stackrel{~}{X}`$ by $`S_{n1}`$, the stabilizer in the permutation group $`W=S_n`$ of one of the $`n`$ eigenvectors. Following , we study the relation between the two types of covers in Section 9. In particular, we show that the above correspondence actually extends to an equivalence between cameral and spectral covers, even when we are very far from the generic situation.
##### 0.3.3. The universal objects
The set of all maximal tori $`TG`$, or equivalently the set of Cartan subalgebras in $`𝔤`$, is parametrized by the quotient $`G/N`$. Over $`X=G/N`$ we have the tautological, unramified Higgs bundle: the underlying $`G`$-bundle is the trivial one, $`X\times G`$, and the regular centralizers are the universal family of Cartans. The corresponding (unramified) cameral cover in this case is $`G/TG/N`$. Note that a point of $`G/T`$ is determined by a Cartan together with a Borel containing it.
The cover $`G/TG/N`$ admits a natural partial compactification $`\overline{G/T}\overline{G/N}`$. Here $`\overline{G/N}`$ paramatrizes regular centralizers in the Lie algebra $`𝔤`$, and $`\overline{G/T}`$ is the ramified $`W`$-cover of $`\overline{G/N}`$ parametrizing pairs consisting of a regular centralizer together with a Borel containing it, cf. Section 1 and Section 10. The map $`\overline{G/T}\overline{G/N}`$ is the cameral cover of the tautologocal Higgs bundle on $`\overline{G/N}`$: the underlying $`G`$-bundle is still $`\overline{G/N}\times G`$, and the regular centralizers form the universal group-scheme $`𝒞`$ of centralizers over $`\overline{G/N}`$. We refer to these as universal objects; every Higgs bundle on $`X`$ is locally the pullback of the tautological one via some map $`X\overline{G/N}`$, and every cameral cover of $`X`$ is locally the pullback of $`\overline{G/T}\overline{G/N}`$ via the same map $`X\overline{G/N}`$.
Although our ultimate results are concerned with Higgs bundles on arbitrary schemes, much of our work boils down to a group-theoretic analysis of these universal objects $`\overline{G/N}`$ and $`\overline{G/T}`$. For instance, we will see that the ramification divisors are indexed by the positive roots $`\alpha `$ of $`G`$. In fact, one of the key points of this paper is that the tautological group-scheme $`𝒞`$ can be completely recovered by looking at the ramification pattern of $`\overline{G/T}`$ over $`\overline{G/N}`$. In a strong sense, this says that a regular centralizer can be recovered from the scheme parametrizing those Borels which contain it. This is our Theorem 11.6. We emphasize that it is the phenomenon described in Theorem 11.6 which is “responsible” for the abelianization.
##### 0.3.4. $`SL(2),PGL(2)`$
We saw that in the general case, the final form of the answer is quite involved. A main source of technical difficulties is the possible presense in $`G`$ of non-primitive coroots (cf. ). ¿From the classification of reductive groups we know that this can occur only when $`G`$ has $`SO(2n+1)`$ as a direct factor. So the simplest case where this extra complication occurs is for $`G=SO(3)=PGL(2)`$. In an attempt to illustrate the effect of these non-primitive coroots, we will, in Section 8, work out explicitly and contrast the examples of $`G=SL(2)`$, for which no $`\beta `$’s are necessary because all coroots are primitive, versus $`G=PGL(2)`$, for which the roots are non-primitive. For these groups, both the spectral cover and the cameral cover are double covers of $`X`$, so the entire analysis can be made much more concrete than for a general group. In particular, there are very explicit descriptions of the universal objects $`\overline{G/T},\overline{G/N},G/T,G/N`$, cf. subsection1.6.
##### 0.3.5. $`K`$-valued Higgs bundles
The point of our abstract notion of a Higgs bundle is that it provides a uniform approach to the analysis of various more concrete objects. In the literature, the most common notion of a Higgs bundle is that of a $`K`$-valued Higgs bundle on $`X`$, where $`K`$ is a fixed line bundle on $`X`$. By definition, this means a pair $`(E_G,s)`$, where $`E_G`$ is a principal $`G`$-bundle on $`X`$ and $`s`$ is a section of $`𝔤_{E_G}K`$. Starting with one of our “abstract” Higgs bundles $`(E_G,𝐜_X)`$, we get a $`K`$-valued Higgs bundle by choosing a section of $`𝐜_XK`$. Conversely, a $`K`$-valued Higgs bundle $`(E_G,s)`$ on $`X`$ determines a unique “abstract” Higgs bundle on the open subset $`X_0X`$ where $`s`$ is regular. We say that a $`K`$-valued Higgs bundle is regular if $`X_0=X`$.
Our philosophy is to think of a regular $`K`$-valued Higgs bundle as involving two separate pieces of data. The first requires specifying the basis of “eigenvectors” of the Higgs field, i.e. it amounts to specifying the underlying abstract Higgs bundle. The other piece of the data corresponds to the “eigenvalues”; in our case this amounts to specifying the section $`s`$ of $`𝐜_xK`$. Our point is that this second part of the data is irrelevant for the abelianization process, so we focus on the “eigenvectors” encoded in the abstract Higgs bundle. One obvious advantage of this approach is that it allows the bundle $`K`$ of “values” to be replaced by various other objects, as we will see below.
A little more generally, we can work with the concept of a regularized $`K`$-valued Higgs bundle on $`X`$, which means a triple $`(E_G,𝐜_X,s)`$, with $`(E_G,𝐜_X)`$ a Higgs bundle in our abstract sense, and $`s`$ a (not necessarily regular!) section of $`𝐜_XK`$. The moduli space of regular $`K`$-valued Higgs bundles is open in the moduli of all $`K`$-valued Higgs bundles (for $`X`$ projective), and is also open inside the moduli space of regularized $`K`$-valued Higgs bundles. For a ”general” Higgs bundle, we can expect the complement of $`X_0`$ to have codimension 3, so if $`X`$ is projective of dimension 1 or 2, we expect the open subset of regular Higgs bundles to be nonempty.
In Section 17, we apply our results to show that the algebraic stack $`\mathrm{𝐇𝐢𝐠𝐠𝐬}(X,K)`$ of regularized $`K`$-valued Higgs bundles on $`X`$ fibers over the affine space $`𝐁(X,K)`$ which parametrizes $`K`$-valued cameral covers, i.e. pairs $`(\stackrel{~}{X},v)`$ where $`v`$ is a $`W`$-equivariant map $`v:\stackrel{~}{X}𝔱K`$ (of schemes over $`X`$). The fibers can be identified with the gerbe $`Higgs_{\stackrel{~}{X}}(X)`$ which we studied in the abstract case. In accordance with our general philosophy, the fiber is independent of the bundle $`K`$ or the way $`\stackrel{~}{X}`$ maps to $`K`$: it depends only on the abstract cameral cover $`\stackrel{~}{X}`$.
In case $`X`$ is a smooth, projective curve and $`K`$ is its canonical bundle, we thus recover a version of Hitchin’s integrable system . (There is of course a difference, in that we work with regularized $`K`$-valued Higgs bundles while Hitchin uses semistable $`K`$-valued Higgs bundles.) As an application, our results can be used to establish a duality between the fibers of the Hitchin map for a group $`G`$ and those corresponding to its Langlands dual group $`\stackrel{ˇ}{G}`$.
##### 0.3.6. Bundles on elliptic fibrations
Essentially no new phenomena are encountered if we allow our Higgs bundle to take its “values” in a vector bundle $`K`$. But we can go further and try to take $`K`$ to be any abelian group scheme over $`X`$, such as the relative Picard scheme of some (projective, integral) family $`f:YX`$. This leads us in Section 18 to define a regularized $`G`$-bundle on $`Y`$ to be the data $`(\stackrel{~}{X},E_G,𝐜_X)`$, with $`\stackrel{~}{X}X`$ a cameral cover of $`X`$, and $`(E_G,𝐜_X)\mathrm{Higgs}_{\stackrel{~}{Y}}(Y)`$ a Higgs bundle on $`Y`$ with cameral cover $`\stackrel{~}{Y}:=f^{}\stackrel{~}{X}`$. This notion is most natural in case $`f`$ is an elliptic fibration, since then we know what it means for a bundle (on $`Y`$) to be regular above a point (of $`X`$). Just as was the situation for $`K`$-valued Higgs bundles, “most” $`G`$-bundles on an elliptic curve are indeed regular, and a regular bundle has a unique regularization.
In Theorem 18.5 we apply our results about abstract Higgs bundles to obtain a complete spectral description of regularized $`G`$-bundles on $`Y`$. In the most interesting case, when $`f`$ is an elliptic fibration, this is the main result of . Letting $`\mathrm{𝐑𝐞𝐠}(X,Y)`$ denote the algebraic stack of regularized $`G`$-bundles on $`Y`$, we obtain a “spectral map” $`h:\mathrm{𝐑𝐞𝐠}(X,Y)𝐁(X,Y)`$, sending a regularized bundle to its $`\mathrm{Pic}(Y/X)`$-valued cameral cover, the fibers now being a slightly twisted version of our gerbe $`\mathrm{Higgs}_{\stackrel{~}{X}}(X)`$.
#### 0.4. Some history
The idea of abelianization has its source in quantum field theory and has been extensively exploited by both physicists and mathematicians. This idea was originally applied not to our notion of an abstract Higgs bundle, but rather to $`K`$-valued Higgs bundles. These were considered by Hitchin in case $`X`$ is a curve and $`K`$ its canonical bundle. Other line bundles, on $`X=P^1`$, were considered by Adams, Harnad and Hurtubise and Beauville . Several aspects of spectral covers of $`P^1`$ and their Prym-Tyurin varieties were considered by Kanev in . The abelianization of $`K`$-valued Higgs bundles on other curves was considered by Beilinson and Kazhdan, Bottacin, Donagi and Markman, Faltings, Markman, and Scognamillo , among others. In the case that the base $`X`$ is a curve, these Higgs bundles are related to representations of the fundamental group of a punctured Riemann surface, as well as to integrable systems arising from loop algebras. The notion of a cameral cover was introduced in , where its relation to the various spectral covers was analyzed.
The main point of many of the works cited above is to show, in various interesting special cases, that the fiber of the Hitchin map, i.e. the family of Higgs bundles with given spectral (or cameral) cover, ”is” generically a Jacobian or a Prym variety, depending on the group. A description of this fiber in the general setting was announced in . In particular, the generalized Prym was described there as a certain quotient of $`H^1(\overline{T}_{\stackrel{~}{X}})`$. (This could be off by a finite isogeny: we have seen that the correct description involves $`H^1(T_{\stackrel{~}{X}})`$.) It was also noted there that the fiber is canonically identified not with the generalized Prym variety itself, but with a certain torsor over it. The class of this torsor was described there in terms of the ”twist” arising from the ramification divisor and the ”shift” by the class of the normalizer $`N`$ in $`H^2(W,T)`$. The additional complication which arises only for $`\mathrm{SO}(2n+1)`$ was first noted in . This is encoded in the present work in our $`\beta `$’s.
Higgs bundles on higher dimensional varieties $`X`$, valued in the cotangent bundle $`K:=T^{}X`$, were introduced by Simpson . Through work of Corlette and Simpson, their moduli spaces are related to those of local systems on $`X`$. The version where $`K`$ is replaced by an elliptic fibration was developed in and . These elliptically valued Higgs bundles are of interest because of their relevance to the construction and parametrization of bundles on elliptic fibrations. These have attracted attention recently because of their importance to understanding the conjectured duality between F-theory and the heterotic string, cf. .
#### 0.5. Notation
We work throughout with a fixed connected reductive group $`G`$ over $``$ and we let $`𝔤`$ denote its Lie algebra. We fix a Borel subgroup $`BG`$ and denote by $`l`$ the flag variety $`G/B`$. By definition, $`l`$ classifies Borel subalgebras in $`𝔤`$.
Let $`U`$ be the unipotent radical of $`B`$ and $`T`$ the Cartan quotient $`B/U`$; we will fix a splitting $`TB`$. We will denote by $`𝔟`$ and $`𝔱`$ the Lie algebras of $`B`$ and $`T`$, respectively. The rank $`r`$ of $`G`$ is by definition the dimension of $`T`$. By $`N`$ we will denote the normalizer of $`T`$ (not the nilpotent subgroup!), and by $`W`$ the Weyl group $`N/T`$.
The set of positive roots will be denoted by $`\mathrm{\Delta }^+`$. For $`\alpha \mathrm{\Delta }^+`$, let $`𝔱^\alpha 𝔱`$ denote the corresponding root hyperplane and $`s_\alpha W`$ the corresponding reflection. The set of simple roots we will denote by $`I`$. For $`iI`$, we will use the notation $`s_i`$ instead of $`s_{\alpha _i}`$.
## Part I Main results on Higgs bundles and cameral covers
### 1. Regular centralizers
#### 1.1.
Recall that an element $`x𝔤`$ is called regular if its centralizer $`Z_𝔤(x)`$ has the smallest possible dimension, namely $`r`$ (the rank of $`𝔤`$). Note that with this definition, a regular element need not be semisimple. The set of all regular elements forms an open subvariety of $`𝔤`$, which we will denote by $`𝔤_{reg}`$.
A Lie subalgebra $`𝔞𝔤`$ is called a regular centralizer if $`𝔞=Z_𝔤(x)`$ for some $`x𝔤_{reg}`$. Note that such $`𝔞`$ is automatically abelian. Our first goal is to introduce a variety which parametrizes all regular centralizers in $`𝔤`$.
#### 1.2.
Let $`\mathrm{Ab}^r`$ be the closed sub-variety in the Grassmannian of $`r`$-planes $`Gr_𝔤^r`$ that classifies abelian subalgebras in $`𝔤`$ of dimension $`r`$. Let $`\mathrm{\Gamma }\mathrm{Ab}^r\times 𝔤`$ be the incidence correspondence, i.e. the closed subvariety defined by the condition:
$$(𝔞,x)\mathrm{\Gamma }\text{ if }x𝔞.$$
Let $`\mathrm{\Gamma }_{reg}`$ be the intersection $`\mathrm{\Gamma }(\mathrm{Ab}^r\times 𝔤_{reg})`$.
###### Proposition 1.3.
There is a smooth morphism $`\varphi :𝔤_{reg}\mathrm{Ab}^r`$ whose graph is $`\mathrm{\Gamma }_{reg}`$.
The proof is postponed until Section 10
Let $`\overline{G/N}`$ denote the image of the map $`\varphi `$. The above proposition implies that $`\overline{G/N}`$ is smooth and irreducible. It is clear that $``$-points of $`\overline{G/N}`$ are exactly the regular centralizers in $`𝔤`$.
By definition, the group $`G`$ acts on both $`\mathrm{Ab}^r`$ and $`𝔤_{reg}`$. Therefore, the variety $`\overline{G/N}`$ acquires a natural $`G`$-action and the map $`\varphi `$ is $`G`$-equivariant.
Consider the quotient $`G/N`$; it classifies Cartan subalgebras in $`𝔤`$. These are the centralizers in $`𝔤`$ of regular semisimple elements. Hence, $`G/N`$ embeds into $`\overline{G/N}`$ as an open subvariety. Obviously, the action of $`G`$ on $`G/N`$ by left multiplication is the restriction of its action on $`\overline{G/N}`$.
#### 1.4.
Consider the closed subvariety of $`\overline{G/N}\times l`$ defined by the condition: for $`𝔞\overline{G/N}`$ and $`𝔟^{}l`$,
$$(𝔞,𝔟^{})\overline{G/T}\text{ if }𝔞𝔟^{}.$$
We will denote this variety by $`\overline{G/T}`$ and the natural projection $`\overline{G/T}\overline{G/N}`$ by $`\pi `$. It follows from the definitions that we have a natural $`G`$-action on $`\overline{G/T}`$.
The quotient $`G/T`$ can clearly be identified with the open sub-scheme $`\pi ^1(G/N)`$ of $`\overline{G/T}`$. We have a natural action of the Weyl group $`W=T\backslash N`$ on $`G/T`$; this action is free and the quotient can be identified with $`G/N`$.
In what follows, by a $`W`$-cover of a scheme $`X`$ we will mean a finite flat scheme $`p:\stackrel{~}{X}X`$, acted on by $`W`$ such that $`p_{}𝒪_{\stackrel{~}{X}}`$ is locally isomorphic, as a coherent sheaf with a $`W`$-action, to $`𝒪_X[W]`$. A basic example is $`𝔱𝔱/W`$: as is well known, it is ramified along the complexified walls of the Cartan subalgebra $`𝔱`$.
The following assertion will be proven in Section 10
###### Proposition 1.5.
The variety $`\overline{G/T}`$ is smooth and connected. The $`W`$-action on $`G/T`$ extends to the whole of $`\overline{G/T}`$ and it makes the latter a $`W`$-cover of $`\overline{G/N}`$. Moreover, the two $`W`$-covers $`\overline{G/T}\overline{G/N}`$ and $`𝔱𝔱/W`$ are étale-locally isomorphic.
#### 1.6.
Here is an explicit description of $`\overline{G/N}`$ and $`\overline{G/T}`$ for $`G=SL(2)`$. In this case $`\overline{G/N}`$ is the space of all lines in $`𝔤`$, i.e. $`\overline{G/N}^2`$. We have a natural map $`\overline{G/T}^1\times ^1`$, where the first projection is the natural map $`\overline{G/T}l^1`$ and the second projection is a composition of the first one with the action of $`1S_2W`$ on $`\overline{G/T}`$.
It is easy to see that this map is an isomorphism. Under the identification, $`\pi :\overline{G/T}\overline{G/N}`$ is the symmetrization map $`^1\times ^1^2`$.
#### 1.7. $`G`$-orbits
For each root $`\alpha `$, let $`D^\alpha \overline{G/T}`$ denote the fixed point set of $`s_\alpha `$ on $`\overline{G/T}`$. This is a smooth codimension $`1`$ subscheme of $`\overline{G/T}`$. Indeed, using the étale-local isomorphism between $`\overline{G/T}\overline{G/N}`$ and $`𝔱𝔱/W`$ given in Proposition 1.5, it is enough to prove this statement on $`𝔱`$. However, $`𝔱^{s_\alpha }`$ is just the corresponding root hyperplane $`𝔱^\alpha 𝔱`$.
###### Proposition 1.8.
The $`G`$-orbits in $`\overline{G/T}`$ are precisely the locally closed subsets
$$D^\mathrm{\Delta }^{}:=\underset{\alpha \mathrm{\Delta }^{}}{}(D^\alpha )\underset{\beta \mathrm{\Delta }^{}}{}(D^\beta )$$
where $`\mathrm{\Delta }^{}\mathrm{\Delta }`$ is a subset of the set of roots, closed under linear combinations. The $`G`$-orbits in $`\overline{G/N}`$ are the images of the $`D^\mathrm{\Delta }^{}`$; they are indexed by the $`\mathrm{\Delta }^{}`$ modulo the action of $`W`$.
The proof will be given in Section 10.
### 2. Higgs bundles and cameral covers
#### 2.1. Higgs bundles
A family of Cartan subalgebras parametrized by a scheme $`X`$ is given by a map from $`X`$ to $`G/N`$. Equivalently, it is given by a $`G`$-equivariant map from the trivial $`G`$-bundle over $`X`$ to $`G/N`$. An advantage of this latter description is that there is a natural way to twist it: given any principal $`G`$-bundle $`E_G`$ over $`X`$, we specify a family of Cartan subalgebras in the adjoint bundle $`𝔤_{E_G}:=E_G\underset{𝐺}{\times }𝔤`$ by a $`G`$-equivariant map from $`E_G`$ to the variety $`G/N`$. By generalizing this, we define:
###### Definition 2.2.
A Higgs bundle over a scheme $`X`$ is a pair $`(E_G,\sigma )`$, where $`E_G`$ is a principal $`G`$-bundle over $`X`$ and $`\sigma `$ is a $`G`$-equivariant map $`\sigma :E_G\overline{G/N}`$.
Therefore, according to Proposition 1.3, a Higgs structure in a given $`G`$-bundle $`E_G`$ is the same as a vector subbundle $`𝐜_X`$ of $`𝔤_{E_G}`$ of rank $`r`$ such that $`[𝐜_X,𝐜_X]=0`$ and such that locally in the étale topology $`𝐜_X`$ is the sheaf of centralizers of a section of $`E_G\underset{𝐺}{\times }𝔤_{reg}`$.
The restriction of a Higgs bundle to an open subset $`UX`$ over which $`E_G`$ is trivialized can be specified more simply by a map $`U\overline{G/N}`$. In particular, the universal Higgs bundle over $`\overline{G/N}`$ corresponds to the identity map $`\overline{G/N}\overline{G/N}`$.
#### 2.3. The Higgs category and stack
Higgs bundles over $`X`$ form a category, denoted $`\mathrm{Higgs}(X)`$. An element of $`\mathrm{Hom}((E_G^1,\sigma ^1),(E_G^2,\sigma ^2))`$ is by definition a $`G`$-bundle map $`s:E_G^1E_G^2`$ such that $`\sigma ^2s=\sigma ^1`$.
One can say that $`\mathrm{Higgs}(X)`$ is the category of maps from $`X`$ to the stack $`G\backslash (\overline{G/N})`$. Additionally, for a fixed $`X`$, we can consider the functor on the category of schemes, which attaches to a scheme $`S`$ the category $`\mathrm{Higgs}(S\times X)`$. When $`X`$ is projective, this functor is representable by an algebraic stack, which we will denote by $`\mathrm{𝐇𝐢𝐠𝐠𝐬}(X)`$. (The representability follows because the stack $`\mathrm{𝐁𝐮𝐧}_G(X)`$ classifying principal $`G`$-bundles on $`X`$ is an algebraic stack. We have: $`G\backslash (\overline{G/N})=\mathrm{𝐇𝐢𝐠𝐠𝐬}(\mathrm{Spec}())`$.
#### 2.4. Cameral covers
We will now introduce our second basic object.
###### Definition 2.5.
A $`W`$-cover of a scheme $`X`$ is a scheme $`\stackrel{~}{X}\stackrel{𝜋}{}X`$ finite and flat over $`X`$ such that as an $`𝒪_X`$-module with a $`W`$-action, $`\pi _{}(𝒪_{\stackrel{~}{X}})`$ is locally isomorphic to $`𝒪_X[W]`$.
###### Definition 2.6.
A cameral cover of $`X`$ is a $`W`$-cover $`\stackrel{~}{X}X`$, such that locally with respect to the étale topology on $`X`$, $`\stackrel{~}{X}`$ is a pull-back of the $`W`$-cover $`𝔱𝔱/W`$.
As an example, we note that any $`W`$-cover is cameral when $`G=\mathrm{SL}(2)`$, i.e. $`W=S_2`$. On the other hand, not every $`W=S_3`$-cover is cameral: the stabilizer of each point must be a Weyl subgroup of $`W`$, so, for example, an $`A_3`$ stabilizer is not allowed.
#### 2.7. Openness
It is easy to see that the condition for a $`W`$-cover $`\stackrel{~}{X}X`$ to be cameral is open on $`X`$. Indeed, $`\pi :\stackrel{~}{X}X`$ is cameral if and only if, locally on $`X`$, we can find a $`W`$-equivariant embedding $`\stackrel{~}{X}X\times 𝔱`$. (Note that the space of $`W`$-equivariant maps of $`X`$-schemes $`\stackrel{~}{X}X\times 𝔱`$ is ismorphic to the space of sections of the sheaf $`\mathrm{Hom}_{𝒪_X}^W(𝔱^{}𝒪_X,\pi _{}(𝒪_{\stackrel{~}{X}}))`$, and the latter sheaf is non-canonically isomorphic to $`𝔱𝒪_X`$, since $`\stackrel{~}{X}X`$ was assumed to be a $`W`$-cover.)
#### 2.8. The cameral category and stack
Cameral covers form a category in a natural way, denoted $`\mathrm{Cam}(X)`$. By definition, $`\mathrm{Hom}(\stackrel{~}{X}^1,\stackrel{~}{X}^2)`$ consists of all $`W`$-equivariant isomorphisms $`\stackrel{~}{X}^1\stackrel{~}{X}^2`$. It is easy to see that there exists an algebraic stack $`\mathrm{𝐂𝐚𝐦}`$, such that $`\mathrm{Cam}(X)`$ is the category $`\mathrm{Hom}(X,\mathrm{𝐂𝐚𝐦})`$.
Indeed, consider the space of commutative $`W`$-equivariant ring structures on the vector space $`V:=[W]`$. This is clearly an affine scheme, and let us denote it by $`\mathrm{𝐂𝐨𝐯}`$. By construction, there exists a universal $`W`$-cover $`\stackrel{~}{\mathrm{𝐂𝐨𝐯}}\mathrm{𝐂𝐨𝐯}`$. Let $`\mathrm{𝐂𝐚𝐦}^{}`$ be the maximal open subscheme of $`\mathrm{𝐂𝐨𝐯}`$, over which $`\stackrel{~}{\mathrm{𝐂𝐨𝐯}}`$ is cameral. Let $`Aut^W(V)`$ be the algebraic group of automorphisms of $`V`$ as a $`W`$-module. Clearly, $`Aut^W(V)`$ acts on $`\mathrm{𝐂𝐚𝐦}^{}`$ and the action lifts on $`\stackrel{~}{\mathrm{𝐂𝐨𝐯}}|_{\mathrm{𝐂𝐚𝐦}^{}}`$. We can now let $`\mathrm{𝐂𝐚𝐦}`$ be the stack-theoretic quotient $`Aut^W(V)\backslash \mathrm{𝐂𝐚𝐦}^{}`$.
As for Higgs bundles, for a fixed $`X`$ we can consider the functor $`S\mathrm{Cam}(S\times X)`$. For $`X`$ projective this functor is representable by an algebraic stack $`\mathrm{𝐂𝐚𝐦}(X)`$.
###### Proposition 2.9.
There is a natural functor $`F:\mathrm{Higgs}(X)\mathrm{Cam}(X)`$. In particular, for a projective scheme $`X`$, we obtain a map between algebraic stacks $`\mathrm{𝐇𝐢𝐠𝐠𝐬}(X)\mathrm{𝐂𝐚𝐦}(X)`$.
###### Proof.
Any map $`\sigma :E_G\overline{G/N}`$ determines a cameral cover $`\stackrel{~}{E}_G`$ of $`E_G`$, namely $`\overline{G/T}\underset{\overline{G/N}}{\times }E_G`$, cf. Proposition 1.5.
For a Higgs bundle, which involves a $`G`$-equivariant map $`\sigma `$, the cameral cover $`\stackrel{~}{E}_GE_G`$ is itself $`G`$-equivariant, so by descent theory, it is pulled back from a unique cameral cover $`\stackrel{~}{X}X`$.
Clearly, the assignment $`(E_G,\sigma )\stackrel{~}{X}`$ constructed above is functorial.
Over an open set $`UX`$ where $`E_G`$ is trivialized, the restriction $`\stackrel{~}{U}U`$ of the cameral cover is given in terms of $`\sigma `$ as: $`\overline{G/T}\times _{\overline{G/N}}U`$. For example, applying this to the universal Higgs bundle over $`\overline{G/N}`$ gives the cameral cover $`\overline{G/T}\overline{G/N}`$. For this reason we refer in this paper to $`\overline{G/T}\overline{G/N}`$ (rather than $`𝔱𝔱/W`$) as the universal cameral cover.
#### 2.10. The fiber
Let us now fix a cameral cover $`\stackrel{~}{X}`$. Let $`\mathrm{Higgs}_{\stackrel{~}{X}}(X)`$ denote the category-fiber of the above functor $`F:\mathrm{Higgs}(X)\mathrm{Cam}(X)`$ over $`\stackrel{~}{X}`$.
In other words, the objects of $`\mathrm{Higgs}_{\stackrel{~}{X}}(X)`$ are pairs
$$((E_G,\sigma )\mathrm{Higgs}(X),t:F(E_G,\sigma )\stackrel{~}{X})$$
and $`\mathrm{Hom}((E_G^1,\sigma ^1,t^1),(E_G^2,\sigma ^2,t^2))`$ is the set of all bundle maps $`s:E_G^1E_G^2`$ with $`\sigma ^2s=\sigma ^1`$ and such that the composition
$$\stackrel{~}{X}\stackrel{(t^1)^1}{}F(E_G^1,\sigma ^1)F(E_G^2,\sigma ^2)\stackrel{t^2}{}\stackrel{~}{X}$$
is the identity automorphism of $`\stackrel{~}{X}`$.
The goal of this paper is to describe explicitly the category $`\mathrm{Higgs}_{\stackrel{~}{X}}(X)`$ in terms of the $`W`$-action on $`\stackrel{~}{X}`$.
### 3. Gerbes
#### 3.1.
Since the objects we study have automorphisms, it is difficult to describe them adequately without the use of some categorical language. Specifically, our description requires the notion of an $`A`$-gerbe, where $`A`$ is a sheaf of abelian groups on $`X`$. This is a particularly useful case of the more general notion of a gerbe over a sheaf of Picard categories. In this section we review the corresponding definitions. For more details, the reader is referred to or
Let $`\mathrm{Sch}_{et}(X)`$ denote the big étale site over $`X`$. (By definition, $`\mathrm{Sch}_{et}(X)`$ is the category of all schemes over $`X`$ and the covering maps are surjective étale morphisms.)
#### 3.2.
Recall that a presheaf $`𝒬`$ of categories on $`\mathrm{Sch}_{et}(X)`$ assigns to every object $`UX`$ in $`\mathrm{Sch}_{et}(X)`$ a category $`𝒬(U)`$ and to every morphism $`f:U_1U_2`$ in $`\mathrm{Sch}_{et}(X)`$ a functor $`f_𝒬^{}:𝒬(U_2)𝒬(U_1)`$. Moreover, for every composition $`U_1\stackrel{𝑓}{}U_2\stackrel{𝑔}{}U_3`$ there should be a natural transformation $`f_𝒬^{}g_𝒬^{}(gf)_𝒬^{}`$, such that an obvious compatibility relation for three-fold compositions holds.
A presheaf $`𝒬`$ of categories on $`\mathrm{Sch}_{et}(X)`$ is said to be a sheaf of categories (or a stack) if the following two axioms hold:
Axiom SC-1. For $`UX`$ in $`\mathrm{Sch}_{et}(X)`$ and a pair of objects $`C_1,C_2𝒬(U)`$, the presheaf of sets on $`\mathrm{Sch}_{et}(U)`$ that assigns to $`f:U^{}U`$ the set
$$\mathrm{Hom}_{𝒬(U^{})}(f_𝒬^{}(C_1),f_𝒬^{}(C_2))$$
is a sheaf.
Axiom SC-2. If $`f:U^{}U`$ is a covering, then the category $`𝒬(U)`$ is equivalent to the category of descent data on $`𝒬(U^{})`$ with respect to $`f`$ (i.e. every descent data on $`𝒬(U^{})`$ with respect to $`f`$ is canonically effective, cf. , p. 221).
#### 3.3.
Here is our main example of a sheaf of categories. Fix a cameral cover $`\stackrel{~}{X}X`$. For every object $`U\mathrm{Sch}_{et}(X)`$ write $`\stackrel{~}{U}:=U\underset{𝑋}{\times }\stackrel{~}{X}`$, which is a cameral cover of $`U`$.
We define the presheaf of categories $`\mathrm{Higgs}_{\stackrel{~}{X}}`$ by $`\mathrm{Higgs}_{\stackrel{~}{X}}(U):=\mathrm{Higgs}_{\stackrel{~}{U}}(U)`$ (the functors $`\mathrm{Higgs}_{\stackrel{~}{X}}(U)\mathrm{Higgs}_{\stackrel{~}{X}}(U^{})`$ for $`U^{}U`$ and the corresponding natural transformations are defined in a natural way).
The following is an easy exercise in descent theory:
###### Lemma 3.4.
$`\mathrm{Higgs}_{\stackrel{~}{X}}`$ satisfies SC-1 and SC-2.
#### 3.5.
Recall that a Picard category is a groupoid endowed with a a structure of a tensor category, in which every object is invertible. A basic example (and the source of the name) is the category of line bundles over a scheme.
A sheaf of categories $`𝒫`$ is said to be a sheaf of Picard categories if for every $`(UX)\mathrm{Sch}_{et}(X)`$, $`𝒫(U)`$ is endowed with a structure of a Picard category such that the pull-back functors $`f_𝒫^{}`$ are compatible with the tensor structure in an appropriate sense. If $`𝒫_1`$ and $`𝒫_2`$ are two sheaves of Picard categories one defines (in a straightforward fashion) a notion of a tensor functor between them.
A typical and the most important example of a sheaf of Picard categories can be constructed as follows:
Let $`𝒜`$ be a sheaf of abelian groups over $`\mathrm{Sch}_{et}(X)`$. For an object $`f:UX`$ of $`\mathrm{Sch}_{et}(X)`$ let $`\mathrm{Tors}_𝒜(U)`$ denote the category of $`𝒜|_U`$-torsors on $`U`$. This is a Picard category and it is easy to see that the assignment $`U\mathrm{Tors}_𝒜(U)`$ defines a sheaf of Picard categories on $`\mathrm{Sch}_{et}(X)`$ which we will denote by $`\mathrm{Tors}_𝒜`$.
#### 3.6.
Just as a torsor is a space on which a group acts simply transitively, a gerbe is a category on which a Picard category acts simply transitively: A category $`𝒬`$ is said to be a gerbe over the Picard category $`𝒫`$, if $`𝒫`$ acts on $`𝒬`$ as a tensor category and for any object $`C𝒬`$ the functor $`𝒫𝒬`$ given by
$$P𝒫\mathrm{Action}(P,C)𝒬$$
is an equivalence.
Now, if $`𝒫`$ is a sheaf of Picard categories and $`𝒬`$ is another sheaf of categories we say that $`𝒬`$ is a gerbe over the sheaf of Picard categories $`𝒫`$, if the following holds:
* For every $`(UX)\mathrm{Sch}_{et}(X)`$, $`𝒬(U)`$ has a structure of a gerbe over $`𝒫(U)`$. This structure is compatible with the pull-back functors $`f_𝒫^{}`$ and $`f_𝒬^{}`$.
* There exists a covering $`UX`$, such that $`𝒬(U)`$ is non-empty.
A basic feature of gerbes is that if $`𝒬_1`$ and $`𝒬_2`$ are gerbes over $`𝒫`$, one can form a new gerbe $`𝒬_1\underset{𝒫}{}𝒬_2`$, called their tensor product, cf. .
#### 3.7.
The basic example of a gerbe over an arbitrary sheaf of Picard categories $`𝒫`$, is $`𝒫`$ itself. Here is a less trivial example:
Fix a short exact sequence $`0𝒜𝒜^{\prime \prime }𝒜^{}0`$ of sheaves of abelian groups on $`X`$ and let $`\tau _𝒜^{}`$ be an $`𝒜^{}`$-torsor over $`X`$. We introduce a sheaf of categories $`𝒬=𝒬_{\tau _𝒜^{}}`$ as follows. For $`U\mathrm{Sch}_{et}(X)`$, $`𝒬(U)`$ is the category of all “liftings” of $`\tau _𝒜^{}|_U`$ to an $`𝒜^{\prime \prime }|_U`$-torsor. It is easy to check that $`𝒬`$ is a gerbe over $`𝒫=\mathrm{Tors}_𝒜`$.
In fact, gerbes over $`\mathrm{Tors}_𝒜`$ can be classified cohomologically:
###### Lemma 3.8.
There is a bijection between the set of equivalence classes of gerbes over $`\mathrm{Tors}_𝒜`$ and $`H^2(X,𝒜)`$. For a given gerbe $`𝒬`$ the corresponding class in $`H^2(X,𝒜)`$ vanishes if and only if the category $`𝒬(X)`$ of ”global sections” is non-empty.
In the above example, the class on $`H^2(X,𝒜)`$ corresponds to the image of the class of $`\tau _𝒜^{}`$ under the boundary map $`H^1(X,𝒜^{})H^2(X,𝒜)`$.
#### 3.9.
The following will be needed in Section 18.
Let $`𝐚:𝒫_1`$ and $`𝒫_2`$ be sheaves of Picard categories, and $`𝒫_1𝒫_2`$ be a functor compatible with the tensor structure. We say that $`𝐚`$ is a monomorphism if for every $`U\mathrm{Sch}_{et}(X)`$ the functor $`𝐚(U):𝒫_1(U)𝒫_2(U)`$ is faithful.
We say that $`𝐚`$ is an epimorphism if for every $`U\mathrm{Sch}_{et}(X)`$ and $`P,\stackrel{~}{P}𝒫_1(U)`$ the map of sheaves on $`\mathrm{Sch}_{et}(U)`$: $`\mathrm{Hom}_{𝒫_1(U^{})}(P|_U^{},\stackrel{~}{P}|_U^{})\mathrm{Hom}_{𝒫_2(U^{})}(𝐚(P|_U^{}),𝐚(\stackrel{~}{P})|_U^{})`$ is an epimorphism (in the sense of sheaves).
Similary, if we have three sheaves of Picard categories and tensor functors $`𝐚:𝒫_1𝒫_2`$, and $`𝐛:𝒫_2𝒫_3`$ we say that the form a short exact sequence if $`𝐛`$ is an epimorphism and $`𝐚`$ induces an equivalence between $`𝒫_1`$ and the category-fiber of $`𝒫_2`$ over the unit object in $`𝒫_3`$.
In this case, for every object $`P_3𝒫_3`$, the category fiber of $`𝒫_2`$ over it is, in a natural way, gerbe over $`𝒫_1`$. This generalizes the above example of $`0𝒜𝒜^{\prime \prime }𝒜^{}0`$.
Let now $`𝒬_1`$ be a gerbe over $`𝒫_1`$, and $`𝐚:𝒫_1𝒫_2`$ a tensor functor. In this case one can construct a canonical induced gerbe $`𝒬_2`$ over $`𝒫_2`$ with the property that there exists a functor $`𝒬_1𝒬_2`$, compatible with the $`𝒫_1`$\- and $`𝒫_2`$-actions via $`𝐚`$.
Suppose now that
$$0𝒫_1\stackrel{𝐚}{}𝒫_2\stackrel{𝐛}{}𝒫_30$$
is a short exact sequence of Picard categories, and $`𝒬_1`$ is a gerbe over $`𝒫_1`$. Let $`𝒬_2`$ be the corresponding induced $`𝒫_2`$-gerbe.
The next lemma follows from the definitions in a straightforward way:
###### Lemma 3.10.
There exists a canonical functor $`𝒬_2𝒫_3`$. The category fiber of $`𝒬_2`$ over a given object $`P_3𝒫_3`$ is naturally a $`𝒫_1`$-gerbe, canonically equivalent to the tensor product $`𝒬_1\underset{𝒫_1}{}𝐛^1(P_3)`$.
### 4. $`\mathrm{Higgs}_{\stackrel{~}{X}}`$ is a gerbe
#### 4.1.
Given a cameral cover $`\stackrel{~}{X}X`$, let $`\overline{T}_{\stackrel{~}{X}}`$ be the sheaf “of $`W`$-equivariant maps $`\stackrel{~}{X}T`$” on the étale site over $`X`$. More precisely, for $`U\mathrm{Sch}_{et}(X)`$, $`\overline{T}_{\stackrel{~}{X}}(U)=\mathrm{Hom}_W(\stackrel{~}{U},T)`$, where $`\stackrel{~}{U}`$ is the induced cameral cover of $`U`$ and the subscript “$`W`$” means maps respecting the $`W`$-action.
However, we need a slightly smaller sheaf.
#### 4.2.
Let $`D_X^\alpha `$ (for each positive root $`\alpha `$) be the fixed point scheme of the reflection $`s_\alpha `$ acting on $`\stackrel{~}{X}`$. Locally, this is the pullback of the universal ramification divisor, i.e. $`D^\alpha \overline{G/T}`$.
Let $`\alpha `$ be a root of $`G`$, considered as a homomorphism $`\alpha :T𝔾_m`$. Then any section $`t`$ of $`\overline{T}_{\stackrel{~}{X}}(U)`$ determines a function $`\alpha t:U𝔾_m`$ which goes to its own inverse under the reflection $`s_\alpha `$. In particular, its restriction to the ramification locus $`D_X^\alpha `$ equals its inverse, so it equals $`\pm 1`$. The subsheaf $`T_{\stackrel{~}{X}}\overline{T}_{\stackrel{~}{X}}`$ is defined by the following condition:
$$T_{\stackrel{~}{X}}(U):=\{t\overline{T}_{\stackrel{~}{X}}(U)|(\alpha t)|_{D_U^\alpha }=+1\text{for each root}\alpha \}.()$$
By construction, $`\overline{T}_{\stackrel{~}{X}}/T_{\stackrel{~}{X}}`$ is a $`_2`$-torsion sheaf. Note, in addition, that it suffices to impose condition (\*) for one representative of each orbit of $`W`$ on the set of roots.
Remark. Recall that a coroot $`\stackrel{ˇ}{\alpha }:𝔾_mT`$ is called primitive if $`\mathrm{ker}(\stackrel{ˇ}{\alpha })=1`$ (this is equivalent to saying that $`\stackrel{ˇ}{\alpha }`$ is a primitive element of the lattice of cocharacters of $`T`$.) It is clear that condition (\*) holds automatically for roots whose corresponding coroots are primitive. For example, when the derived group of $`G`$ is simply-connected, all coroots are primitive, i.e. (\*) is automatic and $`T_{\stackrel{~}{X}}=\overline{T}_{\stackrel{~}{X}}`$. In fact, $`G`$ has non-primitive coroots if and only if it contains $`SO(2n+1)`$ (e.g. $`PGL(2)=SO(3)`$) as a direct factor, as is easily seen from the classification of Dynkin diagrams.
#### 4.3.
Our first result can be stated as:
###### Theorem 4.4.
$`\mathrm{Higgs}_{\stackrel{~}{X}}`$ is a gerbe over $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}`$.
Let us list several corollaries of this theorem:
###### Corollary 4.5.
To a cameral cover $`\stackrel{~}{X}`$ there corresponds a class in $`H^2(X,T_{\stackrel{~}{X}})`$, which vanishes if and only if $`\stackrel{~}{X}`$ is the cameral cover corresponding to some Higgs bundle.
This is immediate from Lemma 3.8.
###### Corollary 4.6.
Suppose $`\mathrm{Higgs}_{\stackrel{~}{X}}(X)`$ is non-empty. The set of isomorphism classes of objects in this category carries a simply-transitive action of $`H^1(X,T_{\stackrel{~}{X}})`$. The group of automorphisms of every object is canonically isomorphic to $`T_{\stackrel{~}{X}}(X)`$.
### 5. Ramification
#### 5.1.
We now proceed to the formulation of our main result, Theorem 6.4, which describes the category $`\mathrm{Higgs}_{\stackrel{~}{X}}(X)`$ completely in terms of $`\stackrel{~}{X}`$. For that purpose, we need to introduce some further notation that has to do with the ramification pattern of $`\stackrel{~}{X}`$ over $`X`$.
#### 5.2.
For each root $`\alpha `$ we will define a line bundle $`R_X^\alpha `$ on $`\stackrel{~}{X}`$. Assume first that $`\stackrel{~}{X}`$ is integral. In this case the subscheme $`D_X^\alpha \stackrel{~}{X}`$ is a Cartier divisor, because locally it is the pull-back of $`D^\alpha \overline{G/T}`$. We set $`R_X^\alpha =𝒪(D_X^\alpha )`$.
When $`\stackrel{~}{X}`$ is arbitrary we proceed as follows. The construction is local, so we may assume that $`X`$, and hence also $`\stackrel{~}{X}`$, is affine. Let $`I_X^\alpha `$ be a coherent sheaf on $`\stackrel{~}{X}`$ generated by symbols $`\{g\}`$, for $`\{g𝒪_{\stackrel{~}{X}}|s_\alpha (g)=g\}`$ that satisfy the relations:
$$f\{g\}=\{fg\}\text{ for all }f\text{ such that }s_\alpha (f)=f.$$
Locally, $`I_X^\alpha `$ is the pull-back of the sheaf of ideals of the subscheme $`D^\alpha \overline{G/T}`$. Hence, $`I_X^\alpha `$ is a line bundle. We have a natural map $`I_X^\alpha 𝒪_{\stackrel{~}{X}}`$ that sends $`\{g\}g`$ and, by construction, its cokernel is $`𝒪_{D_X^\alpha }`$.
We define the line bundle $`R_X^\alpha `$ as the inverse of $`I_X^\alpha `$. We have a canonical section $`𝒪_{\stackrel{~}{X}}R_X^\alpha `$ whose locus of zeroes is the subscheme $`D_X^\alpha `$.
#### 5.3.
Consider the $`T`$-bundle $`_X^\alpha :=\stackrel{ˇ}{\alpha }(R_X^\alpha )`$ (i.e. $`_X^\alpha `$ is induced from $`R_X^\alpha `$ by means of the homomorphism $`\stackrel{ˇ}{\alpha }:𝔾_mT`$).
For an element $`wW`$ we introduce the $`T`$-bundle $`_X^w`$ on $`\stackrel{~}{X}`$ as
$$_X^w:=\underset{𝛼}{}_X^\alpha ,$$
where $`\alpha `$ runs over those positive roots for which $`w(\alpha )`$ is negative. For example, for $`w=s_i`$ (a simple reflection), $`_X^{s_i}_X^{\alpha _i}`$.
Observe that given a $`T`$-bundle $``$ on $`\stackrel{~}{X}`$ and an element $`wW`$, there are two ways to produce a new $`T`$-bundle: we can pull back by $`w`$ acting as an automorphism of $`\stackrel{~}{X}`$, or we can conjugate the $`T`$-action by $`w`$. We will always write $`w^{}()`$ for the combination of both actions. For example, for $`G=SL(2)`$, the $`T`$-bundle $``$ is equivalent to a line bundle $`L`$. The two individual actions on $``$ of the non-trivial element $`1S_2=W`$ send $`L`$ to $`(1)^{}(L)`$ and $`L^1`$, respectively, while $`(1)^{}()`$ corresponds to the line bundle $`(1)^{}(L^1)`$. In particular, we have:
(1)
$$w^{}(_X^\alpha )_X^{w^1(\alpha )}.$$
###### Lemma 5.4.
There is a canonical isomoprhism $`_X^{w_1w_2}\stackrel{\varpi (w_1,w_2)}{}w_2^{}(_X^{w_1})_X^{w_2}`$.
The proof follows imediately from the definition of $`_X^w`$ and (1). The following proposition is necessary for the formulation of Theorem 6.4.
###### Proposition 5.5.
Let $`\alpha _i`$ be a simple root and let $`wW`$ be such that $`w(\alpha _i)=\alpha _j`$ (another positive simple root). Then, the line bundle $`\alpha _i(_X^w)|_{D_X^{\alpha _i}}`$ admits a canonical trivialization.
###### Proof.
Let us observe first that, since we are using only roots rather than arbitrary weights, it is sufficient to consider the case when $`[G,G]`$ is simply-connected.
We have $`ws_i=s_jw`$, hence, by Lemma 5.4
$$s_i^{}(_X^w)_X^{s_i}_X^{ws_i}_X^{s_jw}w^{}(_X^{s_j})_X^w.$$
However, by definition $`w^{}(_X^{s_j})_X^{s_i}`$, so we obtain that $`s_i^{}(_X^w)_X^w`$. By restricting to $`D_X^{\alpha _i}`$, we obtain $`\stackrel{ˇ}{\alpha }_i(\alpha _i(_X^w))\stackrel{ˇ}{\alpha }_i(𝒪_{D_X^{\alpha _i}})`$.
Since $`[G,G]`$ is simply-connected, every coroot is primitive. Therefore, there exists a weight $`\lambda `$, such that $`\lambda \stackrel{ˇ}{\alpha }_i=\mathrm{id}:𝔾_m𝔾_m`$. By applying $`\lambda `$ to the above isomorphism $`\stackrel{ˇ}{\alpha }_i(\alpha _i(_X^w))\stackrel{ˇ}{\alpha }_i(𝒪_{D_X^{\alpha _i}})`$, we obtain an isomorphism $`\alpha _i(_X^w)\stackrel{\mathrm{isom}^\lambda }{}𝒪_{D_X^{\alpha _i}}`$.
Now it only remains to check that this isomorphism is independent of the choice of $`\lambda `$. However, since the $`_X^w`$’s are locally pull-backs of the corresponding $`T`$-bundles on $`\overline{G/T}`$, it suffices to consider the universal situation, namely the case $`X=\overline{G/N}`$.
In the latter case, the $`T`$-bundle $`^w|_{D^{\alpha _i}}`$ itself is trivialized over an open dense part of $`D^{\alpha _i}`$, namely over $`D^{\alpha _i}\underset{\alpha \alpha _i}{}(D^\alpha D^{\alpha _i})`$. This is because $`\alpha _i`$ is not among the set of roots which become negative under the action of $`w`$. In particular, we obtain an isomorphism $`\alpha _i(^w)\stackrel{\mathrm{isom}^{}}{}𝒪_{D^{\alpha _i}}`$ over $`D^{\alpha _i}\underset{\alpha \alpha _i}{}(D^\alpha D^{\alpha _i})`$.
Moreover, it is easy to see that for any $`\lambda `$ as above, the isomorphisms $`\mathrm{isom}^\lambda `$ and $`\mathrm{isom}^{}`$ coincide. In particular, $`\mathrm{isom}^\lambda `$ is independent of $`\lambda `$ over $`D^{\alpha _i}\underset{\alpha \alpha _i}{}(D^\alpha D^{\alpha _i})`$ and hence over the whole of $`D^{\alpha _i}`$, which is what we need.
The following notions will be used in the formulation of Theorem 6.4.
###### Definition 5.6.
Let $`_0`$ be a $`T`$-bundle on $`\stackrel{~}{X}`$. We say that it is weakly $`W`$-equivariant if for every $`w`$ there exists an isomorphism $`w^{}(_0)_0`$.
For a weakly $`W`$-equivariant $`T`$-bundle, let $`\mathrm{Aut}(_0)`$ be the group whose elements are pairs: an element $`wW`$ plus an isomorphism $`w^{}(_0)_0`$. By definition, $`\mathrm{Aut}(_0)`$ fits into a short exact sequence:
$$1\mathrm{Hom}(\stackrel{~}{X},T)\mathrm{Aut}(_0)W1.$$
###### Definition 5.7.
A strongly $`W`$-equivariant $`T`$-bundle is a weakly $`W`$-equivariant $`T`$-bundle $`_0`$ plus a choice of a splitting $`\gamma _0:W\mathrm{Aut}(_0)`$.
###### Definition 5.8.
A $`T`$-bundle on $`\stackrel{~}{X}`$ is called weakly $`R`$-twisted $`W`$-equivariant if for every $`wW`$ there exists an isomorphism $`w^{}()_X^w`$.
For a weakly $`R`$-twisted $`W`$-equivariant $`T`$-bundle $``$ we introduce the group $`\mathrm{Aut}_{}()`$. Its elements are pairs $`wW`$ and an ismorphism $`w^{}()_X^w`$. The group law is defined via the isomorphism $`\varpi (w_1,w_2)`$ of Lemma 5.4. By definition, $`\mathrm{Aut}_{}()`$ is also an extension of $`W`$ by means of $`\mathrm{Hom}(\stackrel{~}{X},T).`$
### 6. The main result
#### 6.1.
We need one more piece of notation. For a simple root $`\alpha _i`$, let $`M_i`$ be the corresponding minimal Levi subgroup. Under the projection $`NW`$, the intersection $`N[M_i,M_i]`$ surjects onto $`s_iS_2`$. Let $`𝒩_i`$ denote the preimage of $`s_i`$ in $`N[M_i,M_i]`$.
By definition, if $`n_i`$ and $`n_i^{}`$ are two elements in $`𝒩_i`$, there exists $`c𝔾_m`$ such that $`n_i^{}=\stackrel{ˇ}{\alpha }_i(c)n_i`$.
#### 6.2.
Given a cameral cover $`\stackrel{~}{X}X`$, we introduce the category $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)`$ of “$`R`$-twisted, $`N`$-shifted $`W`$-equivariant $`T`$-bundles on $`\stackrel{~}{X}`$”. Its objects consist of:
* A weakly $`R`$-twisted $`W`$-equivariant $`T`$-bundle $``$ on $`\stackrel{~}{X}`$.
* A map of short exact sequences:
$$\begin{array}{ccccccccc}1& & T& & N& & W& & 1\\ & & \text{natural map}& & \gamma & & \mathrm{id}& & & & \\ 1& & \mathrm{Hom}(\stackrel{~}{X},T)& & \mathrm{Aut}_{}()& & W& & 1\end{array}$$
* For each simple root $`\alpha _i`$ and element $`n_i𝒩_i`$, an isomorphism of line bundles on $`D_X^{\alpha _i}`$
$$\beta _i(n_i):\alpha _i()|_{D_X^{\alpha _i}}R_X^{\alpha _i}|_{D_X^{\alpha _i}}.$$
These data must satisfy three compatibility conditions:
(1) If $`n_i^{}=\stackrel{ˇ}{\alpha }_i(c)n_i`$ for $`c𝔾_m`$, then $`\beta _i(n_i^{})=c\beta _i(n_i)`$.
(2) Let $`\alpha _i`$ be again a simple root and $`n_i𝒩_i`$. Consider the isomorphism
$$\gamma (n_i):s_i^{}()_X^{s_i}.$$
When we restrict it to $`D_X^{\alpha _i}`$ it induces an isomorphism
$$\stackrel{ˇ}{\alpha }_i(\alpha _i()|_{D_X^{\alpha _i}})\stackrel{ˇ}{\alpha }_i(R_X^{\alpha _i}|_{D_X^{\alpha _i}}),$$
by the definition of $`_X^{s_i}`$. We need that this isomorphism coincides with $`\stackrel{ˇ}{\alpha }_i(\beta _i(n_i))`$.
(3) Let $`\alpha _i`$ and $`\alpha _j`$ be two simple roots and let $`wW`$ be such that $`w(\alpha _i)=\alpha _j`$. Let $`\stackrel{~}{w}N`$ be an element that projects to $`w`$, and $`n_j`$ be an element of $`𝒩_j`$. By pulling back the isomorphism $`\beta _j(n_j)`$ with respect to $`w`$, we obtain an isomorphism $`\alpha _i(w^{}())|_{D_X^{\alpha _i}}R_X^{\alpha _i}|_{D_X^{\alpha _i}}`$. In addition, the isomorphisms induced by $`\gamma (\stackrel{~}{w})`$ and Proposition 5.5 lead to a sequence of isomorphisms:
$$\alpha _i()|_{D_X^{\alpha _i}}\stackrel{\gamma (\stackrel{~}{w})}{}\alpha _i(w^{}())|_{D_X^{\alpha _i}}\alpha _i(_X^w)|_{D_X^{\alpha _i}}\stackrel{\text{Proposition }\text{5.5}}{}\alpha _i(w^{}())|_{D_X^{\alpha _i}}.$$
By composing the two, we obtain an isomorphism $`\alpha _i()|_{D_X^{\alpha _i}}R_X^{\alpha _i}|_{D_X^{\alpha _i}}`$ and our condition is that it coincides with $`\beta _i(n_i)`$, where $`n_i=\stackrel{~}{w}^1n_j\stackrel{~}{w}𝒩_i`$.
This concludes the definition of objects of $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)`$. Morphisms between $`(,\gamma ,\beta _i)`$ and $`(^1,\gamma ^1,\beta _i^1)`$ are $`T`$-bundle isomorphism maps $`^1`$, which intertwine in the obvious sense $`\gamma `$ with $`\gamma ^1`$ and $`\beta _i`$ with $`\beta _i^1`$.
#### 6.3.
It is easy to see that $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)`$ can be naturally sheafified. Namely, we define the presheaf of categories $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}`$ by setting for for $`U\mathrm{Sch}_{et}(X)`$, $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(U):=\mathrm{Higgs}_{\stackrel{~}{U}}^{}(U)`$. The pull-back functors are defined in an evident manner and it is easy to see that $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}`$ satisfies SC-1 and SC-2.
Our main result is:
###### Theorem 6.4.
The sheaves of categories $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}`$ and $`\mathrm{Higgs}_{\stackrel{~}{X}}`$ are naturally equivalent.
In particular, we obtain that $`\mathrm{Higgs}_{\stackrel{~}{X}}(X)`$ is equivalent to $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)`$. In other words, a Higgs bundle on $`X`$ with the given cameral cover $`\stackrel{~}{X}`$ is equivalent to a $`T`$-bundle on $`\stackrel{~}{X}`$ which is $`R`$-twisted, $`N`$-shifted $`W`$-equivariant.
#### 6.5. Variant
Assume that all coroots in $`G`$ are primitive, i.e. for every $`\alpha `$, the corresponding $`1`$-parameter subgroup maps injectively into $`T`$.
We claim that the definition of $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)`$ is equivalent to the following (simplified) one. We introduce the category $`\mathrm{Higgs}_{\stackrel{~}{X}}^{\prime \prime }(X)`$ as follows:
Objects of $`\mathrm{Higgs}_{\stackrel{~}{X}}^{\prime \prime }(X)`$ are pairs
* A weakly $`R`$-twisted $`W`$-equivariant $`T`$-bundle $``$ on $`\stackrel{~}{X}`$.
* A map of short exact sequences:
$$\begin{array}{ccccccccc}1& & T& & N& & W& & 1\\ & & \text{natural map}& & \gamma & & \mathrm{id}& & & & \\ 1& & \mathrm{Hom}(\stackrel{~}{X},T)& & \mathrm{Aut}_{}()& & W& & 1,\end{array}$$
such that the following condition holds:
(1’) Let $`\lambda `$ be a weight of $`T`$ such that $`\lambda ,\stackrel{ˇ}{\alpha }_i=0`$, which implies that $`\lambda ()|_{D_X^{\alpha _i}}\lambda (s_i^{}()_X^{s_i})|_{D_X^{\alpha _i}}`$. Our condition is that for every $`n_i𝒩_i`$ the composition
$$\lambda ()|_{D_X^{\alpha _i}}\lambda (s_i^{}()_X^{s_i})|_{D_X^{\alpha _i}}\stackrel{\gamma (n_i)}{}\lambda ()|_{D_X^{\alpha _i}}$$
is the identity map.
Morphisms between $`(,\gamma )`$ and $`(^1,\gamma ^1)`$ are $`T`$-bundle maps, which intertwine between $`\gamma `$ and $`\gamma ^1`$.
Let us show that $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)`$ and $`\mathrm{Higgs}_{\stackrel{~}{X}}^{\prime \prime }(X)`$ are naturally equivalent. Indeed, if we have an object $`(,\gamma ,\beta _i)\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)`$, the corresponding object of $`\mathrm{Higgs}_{\stackrel{~}{X}}^{\prime \prime }(X)`$ is obtained by just forgetting the $`\beta _i`$’s.
Conversely, if $`(,\gamma )\mathrm{Higgs}_{\stackrel{~}{X}}^{\prime \prime }(X)`$, we reconstruct the $`\beta _i`$’s as follows:
For a simple root $`\alpha _i`$ and $`n_i𝒩_i`$ consider the isomorphism $`\gamma (n_i)`$ restricted to $`D_X^{\alpha _i}`$. It yields an isomorphism
$$\stackrel{ˇ}{\alpha }_i(\alpha _i())|_{D_X^{\alpha _i}}\stackrel{ˇ}{\alpha }_i(R_X^{\alpha _i})|_{D_X^{\alpha _i}}.$$
Since $`\stackrel{ˇ}{\alpha }_i`$ is primitive, there exists a weight $`\lambda ^{}`$ with $`\lambda ^{},\stackrel{ˇ}{\alpha }_i=1`$. By evaluating $`\lambda `$ on the above isomorphism, we obtain the required identification $`\beta _i(n_i):\alpha _i()|_{D_X^{\alpha _i}}R_X^{\alpha _i}|_{D_X^{\alpha _i}}`$. This isomorphism does not depend on the choice of $`\lambda ^{}`$ because of our condition (1’) on $`\gamma `$.
The fact that conditions (1) and (2) hold follows from the construction. Condition (3) follows from the way in which we build the isomorphism of Proposition 5.5.
## Part II Basic examples
### 7. The universal example: $`\overline{G/N}`$
#### 7.1.
In the category $`\mathrm{Higgs}_{\overline{G/T}}(\overline{G/N})`$ there is a canonical tautological object. One of the main steps in the proof of Theorem 6.4 is to exhibit the corresponding canonical object in $`\mathrm{Higgs}_{\overline{G/T}}^{}(\overline{G/N})`$. This is our goal in this section.
#### 7.2.
Consider the canonical $`T`$-bundle $`_l=G/U`$ over $`l=G/B`$ and let us denote by $`_{can}`$ its pull-back to $`\overline{G/T}`$ under the natural projection $`\overline{G/T}l`$. This will be the first piece in the data $`(_{can},\gamma _{can},\beta _{i,can})`$.
When we restrict $`_{can}`$ to $`G/T\overline{G/T}`$, it becomes identified with $`GG/T`$. Hence for every element $`\stackrel{~}{w}N`$ that projects to $`wW`$, we obtain an isomorphism $`\gamma _{can}(n):w^{}(_{can})_{can}`$ over $`G/T`$, given by right multipliction by $`\stackrel{~}{w}^1`$ on $`G`$.
However, when exended to the whole of $`\overline{G/T}`$, the above identification is meromorphic and the configuration of its zeroes and poles is given by a divisor on $`\overline{G/T}`$ with values in the cocharacter lattice of $`T`$.
###### Theorem 7.3.
For a simple reflection $`s_i`$, the divisor of the above meromorphic map $`s_i^{}(_{can})_{can}`$ is given by $`\stackrel{ˇ}{\alpha }_i(D^{\alpha _i})`$.
The proof will be given in Section 15.
Since $`_X^{w_1w_2}w_2^{}(_X^{w_1})_X^{w_2}`$, Theorem 7.3 implies that for any element $`wW`$, the divisor of zeroes/poles of the above meromorphic map $`w^{}(_{can})_{can}`$ coincides with $`_{\overline{G/T}}^w`$. Hence, we obtain the data of $`\gamma _{can}:N\mathrm{Aut}_{}(_{can})`$.
Finally, we have to specify the data of $`\beta _{i,can}`$ and check the compatibility conditions. Let us first consider the case when $`[G,G]`$ is simply connected. As was explained in Section 6.5, in this case the data of $`\beta _{i,can}`$ can be recovered from $`\gamma _{can}`$, once we check that condition (1’) holds.
Thus, let $`\alpha _i`$ be a simple root and let $`\lambda `$ be a weight orthogonal to $`\stackrel{ˇ}{\alpha }_i`$. It suffices to check condition (1’) at the generic point of $`D^{\alpha _i}`$. Let $`M_i`$ be the corresponding minimal Levi subgroup. We have a closed embedding $`\overline{M_i/T}\overline{G/T}`$ (cf. Section 10.5) and its orbit under the $`G`$-action is the open subset of $`\overline{G/T}`$ equal to $`G/T(D^{\alpha _i}\underset{\alpha \alpha _i}{}(D^\alpha D^{\alpha _i}))`$. In particular, it contains a dense subset of $`D^{\alpha _i}`$.
Since all our constructions were $`G`$-equivariant, this implies that condition (1’) for $`\alpha _i`$ is equivalent to the corresponding statement for $`M_i`$. Moreover, we can replace $`M_i`$ by an isogenous group, namely $`[M_i,M_i]\times Z(M_i)`$. However, in the latter case our compatibility condition becomes obvious, as $`\lambda `$ factors through $`Z(M_i)`$.
Now, let $`G`$ be arbitrary. Choose an isogeny $`G^{}G`$ such that $`[G^{},G^{}]`$ is simply-connected. The varieties $`\overline{G/T}`$ and $`\overline{G^{}/T^{}}`$ are canonically identified and the $`T`$-bundle $`_{can}`$ is induced from the $`T^{}`$-bundle $`_{can}^{}`$ under $`T^{}T`$. Therefore, once we know the data of $`\beta _{i,can}^{}`$ for $`_{can}^{}`$ that satisfies the compatibility conditions, it produces the corresponding data for $`_{can}`$.
Thus, we have constructed a canonical $`G`$-equivariant object of $`\mathrm{Higgs}_{\overline{G/T}}^{}`$ over $`\overline{G/N}`$.
### 8. Some simple cases
#### 8.1. The unramified situation
We call a Higgs bundle $`(E_G,\sigma )`$ unramified if $`\sigma `$ maps $`E_G`$ to $`G/N`$. Such a map amounts to a reduction of the structure group from $`G`$ to $`N`$. The category of unramified Higgs bundles is therefore equivalent to the category of principal $`N`$-bundles.
The functor $`F:\mathrm{Higgs}(X)\mathrm{Cam}(X)`$ sends an $`N`$-bundle $`E_N`$ to $`\stackrel{~}{X}:=T\backslash E_N`$, which is a principal $`W`$-bundle over $`X`$ (i.e. an étale $`W`$-cover).
In this case the assertion of Theorem 6.4 is quite evident.
#### 8.2. $`G=SL(2)`$.
Fix an $`S_2`$-cover $`p:\stackrel{~}{X}X`$ and consider the subsheaf of $`p_{}(𝒪_{\stackrel{~}{X}})`$ consisting of $`S_2`$-anti-invariants. We will denote it by $`𝐜_X`$.
It is easy to see that the category $`\mathrm{Higgs}_{\stackrel{~}{X}}(X)`$ is canonically equivalent to the category of pairs $`(L,\gamma ^{})`$, where $`L`$ is a line bundle on $`\stackrel{~}{X}`$ and $`\gamma ^{}`$ is an isomorphism $`\mathrm{det}(p_{}(L))𝒪_X`$.
Let $`D_X\stackrel{~}{X}`$ be the ramification divisor and let $`R_X`$ be the corresponding line bundle (cf. Section 5). It is easy to see that the category $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)`$ (which in our case is equivalent to its simplified version $`\mathrm{Higgs}_{\stackrel{~}{X}}^{\prime \prime }(X)`$) consists of pairs $`(L,\gamma )`$, where $`L`$ is a line bundle on $`\stackrel{~}{X}`$ and $`\gamma `$ is an isomorphism $`(1)^{}(L^1)R_XL`$ such that the composition
$$L(1)^{}(R_X)(1)^{}((1)^{}(L^1)R_X)\stackrel{(1)^{}(\gamma )}{}(1)^{}(L)\stackrel{𝛾}{}L(1)^{}(R_X)$$
is minus the identity map.
Let us visualize the equivalence $`\mathrm{Higgs}_{\stackrel{~}{X}}(X)\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)`$ of Theorem 6.4 in this case. Indeed, for any line bundle $`L`$ on $`\stackrel{~}{X}`$ we have a canonical $`S_2`$-equivariant isomorphism
$$p^{}(\mathrm{det}(p_{}(L)))R_XL(1)^{}(L).$$
Therefore, a data of $`\gamma ^{}`$ defines the data of $`\gamma `$, and it is easy to see that this sets up an equivalence.
#### 8.3. $`G=PGL(2)`$.
In this case the only coroot is non-primitive, so one has to work a little harder.
By definition, objects of $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)`$ are the following data:
* A line bundle $`L`$ on $`\stackrel{~}{X}`$.
* An $`S_2`$-equivariant isomorphism of line bundles $`\gamma :L(1)^{}(L)R_X^2`$.
* An identification $`\beta :L|_{D_X}R_X|_{D_X}`$, which is compatible in the obvious sense with the restriction of $`\gamma `$ to $`D_X`$.
Let us make the statement of Theorem 6.4 explicit in this case too. Starting from an object $`(E_G,\sigma ,t)`$ in $`\mathrm{Higgs}_{\stackrel{~}{X}}(X)`$ we can locally choose a principal $`SL(2)`$-bundle $`E_G^1`$, which induces $`E_G`$. Then $`(E_G^1,\sigma ,t)`$ is an $`SL(2)`$-Higgs bundle. Using the above analysis for $`SL(2)`$, we can attach to it a pair $`(L^1,\gamma ^1)`$, where $`L^1`$ is a line bundle on $`\stackrel{~}{X}`$ and $`\gamma ^1:(1)^{}((L^1)^1)R_XL^1`$.
The corresponding object of $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)`$ is constructed as follows. We define the line bundle $`L`$ as $`(L^1)^2`$ and $`\gamma :=(\gamma ^1)^2`$. The data of $`\beta `$ comes from the sequence of isomorphisms
$$(L^1)^1R_X|_{D^X}(1)^{}((L^1)^1)R_X|_{D^X}\stackrel{\gamma ^1}{}L^1|_{D_X}.$$
If we choose a different lifting of $`E_G`$ to an $`SL(2)`$-bundle, the corresponding $`L^1`$ will be modified by tensoring with $`p^{}(L^0)`$, where $`L^0`$ is a line bundle on $`X`$ with $`(L^0)^2𝒪`$, which will not affect the resulting $`(L,\gamma ,\beta )`$.
It is an easy exercise to check that the above construction defines an equivalence of categories.
### 9. Spectral covers versus cameral covers for $`G=GL(n)`$
#### 9.1.
Observe first that a regular centralizer in $`𝔤𝔩(n)`$ is the same as an $`n`$-dimensional associative and commutative subalgebra in $`\mathrm{Mat}(n,n)`$ generated by one element.
###### Definition 9.2.
An $`n`$-sheeted spectral cover of a scheme $`X`$ is a finite flat scheme $`p:\overline{X}X`$ such that $`p_{}(𝒪_{\overline{X}})`$ has rank $`n`$ and is locally uni-generated as a sheaf of algebras.
Thus, a Higgs bundle for $`𝔤𝔩(n)`$ is the same as a rank-$`n`$ vector bundle $`E`$ and an $`n`$-sheeted spectral cover $`\overline{X}X`$ with an embedding of bundles of algebras $`p_{}(𝒪_{\overline{X}})\mathrm{End}_{𝒪_X}(E)`$. This is equivalent to saying that $`E`$ is a line bundle over $`\overline{X}`$.
In this section we will analyze the connection of this description of Higgs bundles for $`GL(n)`$ with the one given by Theorem 6.4. The starting point is the observation that the category of $`S_n`$-cameral covers of $`X`$ is naturally equivalent to the category of $`n`$-sheeted spectral covers. Let us describe the functors in both directions:
Given an $`S_n`$-cameral cover $`\stackrel{~}{X}X`$, we define the scheme $`\overline{\stackrel{~}{X}}`$ as $`S_{n1}\backslash \stackrel{~}{X}`$. Conversely, given an $`n`$-sheeted spectral cover $`\overline{X}X`$, we define $`\stackrel{~}{\overline{X}}`$ to be the scheme that represents the functor of orderings of the sheets of $`\overline{X}X`$. This functor attaches to a scheme $`S`$ the set of data consisting of
(A map $`SX`$ and $`n`$ sections $`t_i:S\overline{S}:=\overline{X}\underset{𝑋}{\times }S`$),
such that the characteristic polynomial of the multiplication action on $`p_{}(𝒪_{\overline{X}})`$ of any function $`f𝒪_{\overline{S}}`$ equals $`\underset{𝑖}{\mathrm{\Pi }}(Yft_i)`$, where $`Y`$ is an indeterminate.
It is easy to see that this functor is indeed representable by a scheme finite over $`X`$. The group $`S_n`$ acts on $`\stackrel{~}{\overline{X}}`$ by permuting the $`t_i`$’s.
###### Proposition 9.3.
The functors $`\stackrel{~}{X}\overline{\stackrel{~}{X}}`$ and $`\overline{X}\stackrel{~}{\overline{X}}`$ send cameral covers to spectral covers and spectral covers to cameral covers, respectively. Moreover, they are inverses of one another.
###### Proof.
Let us consider first the universal situation: $`X_0=\mathrm{Spec}([a_0,\mathrm{},a_{n1}])`$, $`\stackrel{~}{X}_0=\mathrm{Spec}([x_1,\mathrm{},x_n])`$, where the $`x_i`$’s satisfy
$$\underset{𝑖}{\mathrm{\Pi }}(Yx_i)=Y^n+a_{n1}Y^{n1}+\mathrm{}+a_1Y+a_0,$$
and $`\overline{X}_0=\mathrm{Spec}([x_1,a_0,\mathrm{},a_{n1}])`$, where $`x_1`$ satisfies
$$x_1^n+a_{n1}x_1^{n1}+\mathrm{}+a_1x_1+a_0=0.$$
The natural maps $`\stackrel{~}{X}_0X_0`$ and $`\overline{X}_0X_0`$ are a cameral and a spectral cover, respectively and it is easy to see that in this case $`\overline{\stackrel{~}{X}}_0\overline{X}_0`$ and $`\stackrel{~}{\overline{X}}_0\stackrel{~}{X}_0`$.
This proves the first assertion of the proposition. Indeed, any cameral (resp., spectral) cover is locally induced from $`\stackrel{~}{X}_0`$ (resp., $`\overline{X}_0`$).
For a spectral cover $`\overline{X}`$ there is a natural map $`\overline{\stackrel{~}{\overline{X}}}\overline{X}`$, that attaches to a map $`S\stackrel{~}{\overline{X}}`$ given by an $`n`$-tuple $`\{t_1,\mathrm{},t_n\}`$ of maps $`t_i:S\overline{S}`$ the composition $`S\stackrel{t_n}{}\overline{S}\overline{X}`$. The resulting map $`\overline{\stackrel{~}{\overline{X}}}\overline{X}`$ is an isomorphism, because this is so in the universal situation, i.e. for $`\overline{\stackrel{~}{X}}_0X_0`$.
Similarly, we have $`n`$ maps $`\stackrel{~}{X}\overline{\stackrel{~}{X}}`$ which correspond to the natural map $`S_n/S_{n1}\times \stackrel{~}{X}\overline{X}`$. We claim that they define an isomorphism $`\stackrel{~}{X}\stackrel{~}{\overline{\stackrel{~}{X}}}`$.
Indeed, both the fact that these maps satisfy the condition on the characteristic polynomial and that the resulting map is an isomorphism follow from the corresponding facts for $`\stackrel{~}{X}_0`$.
#### 9.4.
Thus, fixing a spectral cover and fixing an $`S_n`$-cameral cover amounts to the same thing. Now, Theorem 6.4 implies that the category $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)`$ is equivalent to the category of line bundles on the corresponding spectral cover $`\overline{X}`$. We would like to explain how to see this equivalence explicitly.
We start with the following observation:
Let $`\stackrel{~}{X}X`$ be an $`S_n`$-cameral cover and let $`\mathrm{Pic}_{\stackrel{~}{X},n}(X)`$ be the groupoid of $`S_n`$-equivariant line bundles $`L`$ on $`\stackrel{~}{X}`$ for which the following condition holds:
For every reflection $`s_{i,j}S_n`$ the isomorphism
$$s_{i,j}^{}(L)L$$
is the identity map on the fixed-point set of $`s_{i,j}`$ in $`\stackrel{~}{X}`$.
###### Proposition 9.5.
The pull-back functor establishes an equivalence between the category of line bundles on $`X`$ and $`\mathrm{Pic}_{\stackrel{~}{X},n}(X)`$.
Let us see first how this proposition implies what we need:
The natural map $`\stackrel{~}{X}\overline{X}`$ is itself an $`S_{n1}`$-cameral cover. On the one hand, by applying the above proposition to this map we obtain that the category of line bundles on $`\overline{X}`$ is equivalent to $`\mathrm{Pic}_{\stackrel{~}{X},n1}(\overline{X})`$.
On the other hand, we claim that $`\mathrm{Pic}_{\stackrel{~}{X},n1}(\overline{X})`$ is equivalent to $`\mathrm{Higgs}_{\stackrel{~}{X}}^{\prime \prime }(X)`$.
Indeed, let us identify the Cartan group of $`GL(n)`$ with the product of $`n`$ copies of $`𝔾_m`$ and let $`\lambda _n:T𝔾_m`$ be the weight corresponding to the last coordinate. Then a functor $`\mathrm{Higgs}_{\stackrel{~}{X}}^{\prime \prime }(X)\mathrm{Pic}_{\stackrel{~}{X},n1}(\overline{X})`$ is given by $`(,\gamma )L:=\lambda _n()`$. It is easy to see that this is indeed an equivalence.
#### 9.6.
Now let us prove Proposition 9.5. The argument will be a prototype of the one we are going to use to prove Theorem 6.4.
Given an object $`L\mathrm{Pic}_{\stackrel{~}{X},n}(X)`$ and a point $`xX`$ we must find an étale neighbourhood of $`x`$ such that, when restricted to the preimage of this neighbourhood, $`L`$ becomes isomorphic to the unit object in $`\mathrm{Pic}_{\stackrel{~}{X},n}(X)`$ (i.e. the one for which $`L=𝒪_{\stackrel{~}{X}}`$ with the tautological $`S_n`$-structure).
First, it is easy to reduce the statement to the case when the ramification over $`x`$ is the maximal possible, i.e. when $`x`$ has only one geometric preimage $`\stackrel{~}{x}`$ in $`\stackrel{~}{X}`$. Further, we can assume that $`X`$ (and therefore also $`\stackrel{~}{X})`$ is a spectrum of a local ring.
Choose some trivialization of $`L`$. Its discrepancy with the $`S_n`$-equivariant structure is a $`1`$-cocycle $`S_n\mathrm{Hom}(\stackrel{~}{X},𝔾_m)`$. We must show that this cocycle is homologous to $`0`$.
Let $`K`$ denote the kernel of the map $`\mathrm{Hom}(\stackrel{~}{X},𝔾_m)𝔾_m`$ given by the evaluation at $`\stackrel{~}{x}`$. Our condition on $`L`$ implies that the above cocycle $`S_n\mathrm{Hom}(\stackrel{~}{X},𝔾_m)`$ takes values in $`K`$. However, since $`\stackrel{~}{X}`$ is local, $`K`$ is divisible and torsion-free. Hence $`H^1(S_n,K)=0`$, so our cocycle is cohomologically trivial.
## Part III Basic structure results over $`\overline{G/N}`$
### 10. The structure of $`\overline{G/N}`$
#### 10.1.
The rest of the paper is devoted to the proofs of various results announced in the previous sections. We start with the proof of Proposition 1.3.
###### Proof.
First we need to show that the map $`\varphi :𝔤_{reg}\mathrm{Ab}^r`$ is well-defined, which is equivalent to saying that the projection $`\mathrm{\Gamma }_{reg}𝔤_{reg}`$ is an isomorphism.
Since the latter projection is proper and $`𝔤_{reg}`$ is reduced, it is enough to show that the scheme-theoretic preimage in $`\mathrm{\Gamma }_{reg}`$ of every $`x𝔤_{reg}`$ is isomorphic to $`\mathrm{Spec}()`$.
This is clear on the level of $``$ points, since by definition of regular elements, the only abelian $`r`$-dimensional subalgebra in $`𝔤`$ that contains $`x`$ is its own centralizer.
For $`𝔞\mathrm{Ab}^r`$, the tangent space $`T_𝔞(\mathrm{Ab}^r)`$ can be identified with the space of maps $`T:𝔞𝔤/𝔞`$ that satisfy:
$$y_1,y_2𝔞,[T(y_1),y_2]+[y_1,T(y_2)]=0𝔤.$$
We claim that the tangent space to $`\mathrm{\Gamma }_{reg}(\mathrm{Ab}^r\times x)`$ at $`𝔞\times x`$ is zero. Indeed, this is the space of maps $`T:𝔞𝔤/𝔞`$ as above, for which, moreover $`[T(y),x]=0,y𝔞`$. However, since $`𝔞=Z_𝔤(x)`$, any such $`T`$ is identically zero.
This implies that $`T_{𝔞\times x}(\mathrm{\Gamma }_{reg}\mathrm{Ab}^r\times x)=0`$ which means that $`\mathrm{\Gamma }_{reg}(\mathrm{Ab}^r\times x)`$ is reduced, i.e. $`\mathrm{Spec}()`$.
Now let us show that $`\varphi `$ is smooth. Let $`𝔞\mathrm{Ab}^r`$ be equal to $`\varphi (x)`$. Using the above description of the tangent space to $`\mathrm{Ab}^r`$ it is easy to see that $`d\varphi `$ sends an element $`u𝔤T_x(𝔤_{reg})`$ to the unique map $`T:𝔞𝔤/𝔞`$ that satisfies:
$$[T(y),x]+[y,u]=0,y𝔞.$$
Consider now the map $`\mathrm{𝐞𝐯}:T_𝔞(\mathrm{Ab}^r)𝔤/𝔞`$ given by $`TT(x)`$. The above description of $`d\varphi `$ implies that the composition
$$𝔤T_x(𝔤_{reg})\stackrel{d\varphi }{}T_𝔞(\mathrm{Ab}^r)\stackrel{𝐞𝐯}{}𝔤/𝔞$$
coincides with the tautological projection $`𝔤𝔤/𝔞`$.
However, since $`x`$ is regular, the fact that $`[T(x),y]=[x,T(y)],y𝔞`$ implies that $`\mathrm{𝐞𝐯}`$ is an injection. We conclude that $`\mathrm{𝐞𝐯}`$ is an isomorphism, hence $`\mathrm{Im}(\varphi )`$ is contained in the smooth locus of $`\mathrm{Ab}^r`$. Furthermore, $`d\varphi `$ is surjective, so $`\varphi `$ is smooth as claimed.
#### 10.2.
Let $`\stackrel{~}{𝔤}`$ be the closed sub-variety in $`𝔤\times l`$ defined by the condition: $`(x,𝔟^{})\stackrel{~}{𝔤}\text{ if }x𝔟^{}`$. Let $`\stackrel{~}{𝔤}_{reg}`$ denote the intersection $`\stackrel{~}{𝔤}(𝔤_{reg}\times l)`$ and let $`\stackrel{~}{\pi }`$ denote the projection $`\stackrel{~}{𝔤}_{reg}𝔤_{reg}`$. It is clear that as a variety, $`\stackrel{~}{𝔤}_{reg}`$ is smooth and connected, since it is an open subset in a vector bundle over $`l`$.
###### Proposition 10.3.
There exists a natural $`G`$-invariant map $`\stackrel{~}{\varphi }:\stackrel{~}{𝔤}_{\mathrm{reg}}\overline{G/T}`$, such that the following square is Cartesian:
$$\begin{array}{ccc}\stackrel{~}{𝔤}_{reg}& \stackrel{\stackrel{~}{\varphi }}{}& \overline{G/T}\\ \stackrel{~}{\pi }& & \pi & & \\ 𝔤_{reg}& \stackrel{\varphi }{}& \overline{G/N}\end{array}$$
###### Proof.
Consider the fibered product $`\overline{G/T}\underset{\overline{G/N}}{\times }𝔤_{reg}`$. By definition of $`\overline{G/T}`$, there is a closed embedding
$$\overline{G/T}\underset{\overline{G/N}}{\times }𝔤_{reg}\stackrel{~}{𝔤}_{reg}$$
that sends a triple $`(𝔞\overline{G/N},𝔟^{}l,x𝔤_{reg})\overline{G/T}\underset{\overline{G/N}}{\times }𝔤_{reg}`$ to $`(x,𝔟^{})\stackrel{~}{𝔤}_{reg}`$.
We claim that this embedding is in fact an isomorphism. Indeed, the statement is obvious over the preimage in $`\stackrel{~}{𝔤}_{reg}`$ of the regular semisimple locus of $`𝔤`$. Therefore, the two schemes coincide at the generic point of $`\stackrel{~}{𝔤}_{reg}`$. This implies what we need, since $`\stackrel{~}{𝔤}_{reg}`$ is reduced.
Now we are ready to prove Proposition 1.5:
###### Proof.
The map $`\stackrel{~}{\varphi }:\stackrel{~}{𝔤}_{reg}\overline{G/T}`$ is smooth, since it is a base change of a smooth map. Hence, the fact that $`\stackrel{~}{𝔤}_{reg}`$ is smooth and connected implies that $`\overline{G/T}`$ has the same properties.
A well-known theorem of Kostant (cf. or , p. 277) says that the restriction of the Chevalley map $`𝔤𝔱/W`$ to $`𝔤_{reg}`$ is smooth and that it gives rise to a Cartesian square:
$$\begin{array}{ccc}\stackrel{~}{𝔤}_{reg}& & 𝔱\\ & & & & \\ 𝔤_{reg}& & 𝔱/W\end{array}$$
Therefore, the natural action of $`W`$ on the preimage in $`\stackrel{~}{𝔤}`$ of the regular semisimple locus in $`𝔤`$ extends to the whole of $`\stackrel{~}{𝔤}_{reg}`$. The same is true for $`\overline{G/T}`$, because the map $`\stackrel{~}{\varphi }`$ is flat and surjective. The étale local isomorphism follows from comparison of our Cartesian square with that of Proposition 10.3.
#### 10.4.
Now let us prove Proposition 1.8.
###### Proof.
Let $`\mathrm{\Delta }^{}`$ be as in the formulation of the proposition. Consider an element $`t𝔱`$ such that $`\alpha (t)=0`$ for $`\alpha \mathrm{\Delta }^{}`$ and $`\beta (t)0`$ for $`\beta \mathrm{\Delta }^{}`$.
In this case $`𝔪:=Z_𝔤(t)`$ is a Levi subalgebra of $`𝔤`$. Let $`M`$ be the corresponding Levi subgroup. It is well-known that $`𝔪𝔟`$ is a Borel subalgebra in $`𝔪`$. Let $`u`$ be an element in the unipotent radical of $`𝔪𝔟`$, which is regular with respect to $`M`$.
We then see that $`x=t+u`$ is a regular element in $`𝔤`$, since $`Z_𝔤(x)=Z_𝔪(u)`$. It is known that if a Borel subalgebra contains a regular element, then it also contains its centralizer (cf. Lemma 11.5). Therefore, $`(Z_𝔪(u),𝔟)\overline{G/T}`$. Moreover, it is easy to see that every pair $`(𝔞,𝔟^{})\overline{G/T}`$ is $`G`$-conjugate to one of the above form.
To conclude the proof, it remains to show that $`(Z_𝔪(u),𝔟)\underset{\alpha \mathrm{\Delta }^{}}{}(D^\alpha )\underset{\beta \mathrm{\Delta }^{}}{}(D^\beta )`$. For that, it suffices to show that the image of $`(t+u,𝔟)`$ as above under $`\stackrel{~}{𝔤}_{reg}𝔱`$ belongs to the corresponding locus of $`𝔱`$. However, the above image is just $`t`$, which makes the assertion obvious.
#### 10.5. Levi subgroups
Let $`JI`$ be a subset. It defines a root subsystem $`\mathrm{\Delta }_J`$ and let $`M_J`$ (resp., $`P_JG`$, $`W_JW`$) denote the corresponding standard Levi subgroup (resp., standard parabolic, Weyl subgroup). Let $`N_{M_J}`$ be the intersection $`M_JN`$, which is the normalizer of $`T`$ in $`M_J`$.
It is easy to see that the natural map $`M_J/N_{M_J}G/N`$ extends to a map $`i_J:\overline{M_J/N_{M_J}}\overline{G/N}`$. In fact, $`\overline{M_J/N_{M_J}}`$ is a closed sub-variety of $`\overline{G/N}`$ which corresponds to $`\{𝔞\overline{G/N}|𝔞𝔪_J\}`$.
###### Proposition 10.6.
There is a canonical $`W`$-equivariant isomorphism:
$$\stackrel{~}{i}_J:W\stackrel{W_J}{\times }\overline{M_J/T}\overline{M_J/N_{M_J}}\underset{\overline{G/N}}{\times }\overline{G/T}.$$
###### Proof.
First, we have a natural closed embedding
$$\overline{M_J/T}\overline{M_J/N_{M_J}}\underset{\overline{G/N}}{\times }\overline{G/T}\overline{G/T}.$$
Its image consists of pairs $`(𝔞,𝔟^{})\overline{G/T}`$ such that $`𝔞𝔪_J`$ and $`𝔟^{}𝔭_J:=\mathrm{Lie}(P_J)`$.
This map is compatible with the $`W_J`$-action. Hence, it extends to a finite map
$$\stackrel{~}{i}_J:W\stackrel{W_J}{\times }\overline{M_J/T}\overline{M_J/N_{M_J}}\underset{\overline{G/N}}{\times }\overline{G/T}.$$
Since both varieties are smooth, in order to prove that $`\stackrel{~}{i}_J`$ is an isomorphism, it suffices to do so over the open part, i.e. over $`M_J/N_{M_J}`$. However, in the latter case, the assertion becomes obvious.
It is easy to see that the $`G`$-orbit of $`\overline{M_J/N_{M_J}}\underset{\overline{G/N}}{\times }\overline{G/T}\overline{G/T}`$ (resp., $`\overline{M_J/T}\overline{G/T}`$) is the union of those $`D^\mathrm{\Delta }^{}`$ for which $`\mathrm{\Delta }^{}`$ is $`W`$-conjugate to a subset of $`\mathrm{\Delta }_J`$ (resp., $`\mathrm{\Delta }^{}\mathrm{\Delta }_J`$.)
### 11. The group-scheme of centralizers
In this section we will formulate two basic theorems, Theorem 11.6 and Theorem 11.8, which will be used for the proof of our first main result, Theorem 4.4.
#### 11.1. The universal centralizers $`𝒞`$ and $`𝐜`$
Consider the constant group-scheme $`G\times \overline{G/N}`$ over $`\overline{G/N}`$, and let $`𝒞G\times \overline{G/N}`$ be its closed group-subscheme of “centralizers”. In other words, $`𝒞`$ is defined by the condition that $`(gG,𝔞\overline{G/N})𝒞`$ if $`g`$ commutes with $`𝔞`$. Clearly, $`𝒞`$ is equivariant with respect to the $`G`$-action on $`\overline{G/N}`$.
Note that the corresponding bundle $`𝐜`$ of Lie algebras can be identified with the tautological rank $`r`$ vector bundle over $`\overline{G/N}`$ which comes from the embedding $`\overline{G/N}Gr_𝔤^r`$. Another interpretation of this $`𝐜`$, considered as a subbundle of the trivial bundle $`𝔤\times \overline{G/N}`$, is that it is the family $`𝐜_{\overline{G/N}}`$ of centralizers of the universal Higgs bundle on $`\overline{G/N}`$, which was studied in detail in Section 7. (Recall from Section 2 that a Higgs bundle $`(E_G,\sigma )`$ on any $`X`$ determines, and is determined by, a subbundle $`𝐜_X`$ consisting of regular centralizer subalgebras of the adjoint bundle $`𝔤_{E_G}`$.)
###### Proposition 11.2.
The group-scheme $`𝒞`$ is commutative and smooth over $`\overline{G/N}`$ and is irreducible as a variety.
###### Proof.
Let $`𝒞^{}`$ be the group-subscheme of $`G\times 𝔤_{reg}`$ over $`𝔤_{reg}`$ defined by the condition:
$$𝒞^{}:=\{(g,x)G\times 𝔤_{reg}|\mathrm{Ad}_g(x)=x\}.$$
First, let us show that $`𝒞^{}`$ is commutative and smooth over $`𝔤_{reg}`$.
Let $`(g,x)`$ be a $``$-point of $`𝒞^{}`$. The tangent space to $`𝒞^{}`$ at $`(g,x)`$ consists of pairs $`(\xi ,y)𝔤\times 𝔤`$ such that $`\mathrm{Ad}_g([x,\xi ])=\mathrm{Ad}_g(y)y`$. The differential of the map $`𝒞^{}𝔤_{reg}`$ sends $`(\xi ,y)`$ to $`y`$. We claim that it is surjective.
It is known that if $`G`$ is of adjoint type, then the centralizer of every regular element is connected. (In particular, each $`Z_G(x)`$ is commutative; this holds even if $`G`$ is not of adjoint type.). Therefore, $`\underset{gZ_G(x)}{\mathrm{Span}}(\mathrm{Ad}_g(y)y)=\mathrm{Im}(\mathrm{ad}_{Z_𝔤(x)})`$. However, the latter, as we saw in the proof of Proposition 1.3, coincides with $`\mathrm{Im}(\mathrm{ad}_x)`$, since $`x`$ is regular.
To prove that $`𝒞^{}`$ is smooth over $`𝔤_{reg}`$, it remains to observe that the fibers of $`𝒞^{}`$ are smooth (since they are algebraic groups in char.$`0`$) and all have dimension $`r`$, by the definition of $`𝔤_{reg}`$. The fact that $`𝒞^{}`$ is commutative was established in the course of the above argument.
Now let us prove the assertion for $`𝒞`$. We have a natural closed embedding $`𝒞\underset{\overline{G/N}}{\times }𝔤_{reg}𝒞^{}`$, which is an isomorphism over the regular semisimple locus of $`𝔤_{reg}`$. Hence, it is an isomorphism, because $`𝒞^{}`$ is reduced. Therefore, since the map $`\varphi :𝔤_{reg}\overline{G/N}`$ is flat and surjective, this shows that $`𝒞`$ is commutative and smooth over $`\overline{G/N}`$. It is irreducible, because this is obviously true over $`G/N`$.
#### 11.3. The group scheme $`𝒯`$.
Now we will introduce another group-scheme over $`\overline{G/N}`$, seemingly of a different nature. Recall the sheaves $`\overline{T}_{\stackrel{~}{X}},T_{\stackrel{~}{X}}`$ introduced in section 4.
Consider the contravariant functor Schemes $``$ Groups which assigns to a scheme $`S`$ the set of pairs
(A map $`S\overline{G/N}`$, a $`W`$-equivariant map $`\stackrel{~}{S}:=S\underset{\overline{G/N}}{\times }\overline{G/T}T).`$
It is easy to see that this functor is representable by an abelian group-scheme over $`\overline{G/N}`$, which we will denote by $`\overline{𝒯}`$. Therefore, once $`S\overline{G/N}`$ is fixed, $`\mathrm{Hom}_{\overline{G/N}}(S,\overline{𝒯})\mathrm{\Gamma }(S,\overline{T}_{\stackrel{~}{S}})`$. In other words, $`\overline{𝒯}`$ represents the sheaf $`\overline{T}_{\overline{G/T}}`$ on $`\mathrm{Sch}_{et}(\overline{G/N})`$.
Clearly, the $`G`$-action on $`\overline{G/T}`$ gives rise to a $`G`$-action on $`\overline{𝒯}`$.
We define the open group-subscheme $`𝒯`$ of $`\overline{𝒯}`$ by the following condition (\**):
$`\mathrm{Hom}(S,𝒯)`$ consists of those pairs $`(S\overline{G/N},\stackrel{~}{S}T)`$ as above, for which for every root $`\alpha `$ the composition
$$S\underset{\overline{G/N}}{\times }D^\alpha \stackrel{~}{S}T\stackrel{𝛼}{}𝔾_m$$
avoids $`1𝔾_m`$.
Since for any map $`S𝒯`$, the above composition takes values in $`\pm 1𝔾_m`$, condition (\**) is equivalent to condition (\*) in the definition of the sheaf $`T_{\stackrel{~}{S}}`$ (cf. Section 4.2): for a fixed map $`S\overline{G/N}`$, $`\mathrm{Hom}_{\overline{G/N}}(S,𝒯)\mathrm{\Gamma }(S,T_{\stackrel{~}{S}})`$, i.e., the group-scheme $`𝒯`$ represents the sheaf $`T_{\overline{G/T}}`$ on $`\mathrm{Sch}_{et}(\overline{G/N})`$.
#### 11.4.
A remarkable fact is that the group-schemes $`𝒞`$ and $`𝒯`$ are canonically isomorphic. Here we will construct a map between them in one direction.
Let $``$ denote the universal group-scheme of Borel subgroups over $`l`$. Let us denote by $`\stackrel{~}{}`$ its pull-back to $`\overline{G/T}`$. In addition, let us denote by $`\stackrel{~}{𝒞}`$ the pull-back of $`𝒞`$ to $`\overline{G/T}`$.
Both $`\stackrel{~}{}`$ and $`\stackrel{~}{𝒞}`$ are group-subschemes of the constant group-scheme $`G\times \overline{G/T}`$.
###### Lemma 11.5.
$`\stackrel{~}{𝒞}`$ is a closed group-subscheme of $`\stackrel{~}{}`$.
Indeed, since $`\stackrel{~}{𝒞}`$ is reduced and irreducible, it suffices to check that over $`G/N`$, $`\stackrel{~}{𝒞}`$ is contained in $`\stackrel{~}{}`$. However, this is obvious.
We have a natural projection $`T\times l`$. By composing it with the inclusion of Lemma 11.5, we obtain a map
$$𝒞\underset{\overline{G/N}}{\times }\overline{G/T}T.$$
This map respects the group law on $`𝒞`$ and $`T`$ and commutes with the $`W`$-action. (This is because it suffices to check both facts after the restriction to $`G/N`$, where they become obvious.)
Hence, we obtain a homomorphism of group-schemes $`\overline{\chi }:𝒞\overline{𝒯}`$.
###### Theorem 11.6.
The above map $`\overline{\chi }:𝒞\overline{𝒯}`$ defines an isomorphism $`\chi :𝒞𝒯`$.
The proof will be given in the next section.
#### 11.7.
Now we will formulate the second key result which will be used in the proof of Theorem 4.4.
Consider the functor that assigns to a scheme $`S`$ the set of triples $`(\overline{G/N}_S^1,\overline{G/N}_S^2,\nu )`$, where $`\overline{G/N}_S^1`$ and $`\overline{G/N}_S^2`$ are two $`S`$-points of $`\overline{G/N}`$ and $`\nu `$ is a $`W`$-equivariant isomorphism
$$\nu :\stackrel{~}{S}^1\stackrel{~}{S}^2,$$
where $`\stackrel{~}{S}^i`$ is the $`W`$-cover of $`S`$ induced by $`\overline{G/N}_S^i`$ from $`\pi :\overline{G/T}\overline{G/N}`$.
It is easy to see that this functor is representable. Let $``$ denote the representing scheme. Since the $`W`$-cover $`\overline{G/T}\overline{G/N}`$ is $`G`$-equivariant, we obtain a natural map $`\xi :G\times \overline{G/N}`$ which covers the map $`G\times \overline{G/N}\stackrel{\mathrm{Action}\times \mathrm{id}}{}\overline{G/N}\times \overline{G/N}`$.
###### Theorem 11.8.
The above map $`\xi :G\times \overline{G/N}`$ is smooth and surjective.
This theorem will be proven in Section 13.
#### 11.9.
The scheme $``$ lives over $`\overline{G/N}\times \overline{G/N}`$. Let $`_\mathrm{\Delta }`$ denote its restriction to the diagonal. By definition, $`_\mathrm{\Delta }`$ is a group-scheme over $`\overline{G/N}`$ which represents the functor of $`W`$-equivariant automorphisms of $`\overline{G/T}`$ over $`\overline{G/N}`$.
Let $`\mathrm{St}G\times \overline{G/N}`$ be the closed group-subscheme of stabilizers, i.e.
$$(g,𝔞)\mathrm{St}\text{ if }\mathrm{Ad}_g(𝔞)=𝔞.$$
Obviously, $`𝒞`$ is a closed normal group-subscheme of $`\mathrm{St}`$.
The map $`\xi :G\times \overline{G/N}`$ gives rise to a map $`\xi _\mathrm{\Delta }:\mathrm{St}_\mathrm{\Delta }`$.
###### Proposition 11.10.
$`_\mathrm{\Delta }`$ represents the quotient group-scheme $`\mathrm{St}/𝒞`$.
###### Proof.
Theorem 11.8 implies that the map $`\xi _\mathrm{\Delta }:\mathrm{St}_\mathrm{\Delta }`$ is smooth and surjective. Therefore, all we have to show is that if $`S\mathrm{St}`$ is a map such that the induced automorphism of $`\stackrel{~}{S}`$ is trivial, then $`S`$ maps to $`𝒞`$.
Observe that $`_\mathrm{\Delta }`$ acts on $`𝒯`$ via its action on $`\overline{G/T}`$. Since the isomorphism $`\chi :𝒞𝒯`$ is $`G`$-equivariant, we obtain a commutative diagram of actions:
$$\begin{array}{ccc}\mathrm{St}\underset{\overline{G/N}}{\times }𝒞& & 𝒞\\ \xi _\mathrm{\Delta }\times \chi & & \chi & & \\ _\mathrm{\Delta }\underset{\overline{G/N}}{\times }𝒯& & 𝒯,\end{array}$$
where the top horizontal arrow is the adjoint action.
Therefore, if a map $`S\mathrm{St}`$ induces the trivial automorphism of $`\stackrel{~}{S}`$, its adjoint action on $`𝒞`$ is trivial too. But this means that it factors through $`𝒞`$.
Similarly, one shows:
###### Corollary 11.11.
The scheme $``$ represents the quotient group-scheme $`(G\times \overline{G/N})/𝒞`$.
#### 11.12.
Here is one more interpretation of Theorem 11.8:
Clearly, the scheme $``$ with its two projections to $`\overline{G/N}`$ is a groupoid over that latter scheme. According to Theorem 11.8, the above projections are smooth and, therefore, we can consider the algebraic stack $`\backslash (\overline{G/N})`$.
###### Corollary 11.13.
The stack $`\backslash (\overline{G/N})`$ is canonically isomorphic to the stack $`\mathrm{𝐂𝐚𝐦}`$ of Section 2.4.
### 12. Proof of Theorem 11.6
#### 12.1.
We start by establishing a result on compatibility of our objects with restrictions to Levi subgroups. We then verify the Lie-algebraic version of the theorem by restricting to an $`𝔰𝔩(2)`$ subalgebra, and finally we refine this to prove the desired group-theoretic version.
#### 12.2.
Let $`M=M_J`$ be a standard Levi subgroup of $`G`$ (cf. Section 10.5) and let $`𝒞_M`$ be the corresponding sheaf of centralizers over $`\overline{M/N_M}`$.
On the one hand, there is a natural closed embedding
$$𝒞_Mi_J^{}(𝒞):=\overline{M/N_M}\underset{\overline{G/N}}{\times }𝒞.$$
On the other hand, we have the group scheme $`\overline{𝒯}`$ over $`\overline{G/N}`$, as well as the group-scheme $`\overline{𝒯}_M`$ over $`\overline{M/N_M}`$. This time, by Proposition 10.6, we have a canonical isomorphism
$$\overline{𝒯}_Mi_J^{}(\overline{𝒯}):=\overline{M/N_M}\underset{\overline{G/N}}{\times }\overline{𝒯}.$$
Moreover, it induces an isomorphism $`𝒯_Mi_J^{}(𝒯)`$, since if a root $`\alpha `$ is not $`W`$-conjugate to a root in $`M`$, then $`s_\alpha `$ has no fixed points on $`W\stackrel{W_J}{\times }\overline{M/T}`$.
###### Proposition 12.3.
The map $`𝒞_Mi_J^{}(𝒞)`$ is an isomorphism. Moreover, the diagram
$$\begin{array}{ccc}𝒞_M& \stackrel{\chi _M}{}& \overline{𝒯}_M\\ & & & & \\ i_J^{}(𝒞)& \stackrel{i_J^{}(\chi )}{}& i_J^{}(\overline{𝒯})\end{array}$$
is commutative.
###### Proof.
The map $`𝒞_Mi_J^{}(𝒞)`$ is an isomorphism because it is a closed embedding and at the same time an isomorphism over the generic point of $`\overline{M/N_M}`$. Commutativity of the diagram can be checked over the preimage of $`M/N_M`$, in which case it becomes obvious.
#### 12.4.
We will now prove the assertion of Theorem 11.6 on the Lie-algebra level.
Let $`𝐭`$ denote the sheaf of Lie algebras corresponding to $`\overline{𝒯}`$. Obviously, it is isomorphic to $`\mathrm{Lie}(𝒯)`$ as well. By definition, we have: $`𝐭(𝔱\pi _{}(𝒪_{\overline{G/T}}))^W`$. Since $`\pi _{}(𝒪(\overline{G/T}))`$ is locally isomorphic to $`𝒪_{\overline{G/N}}[W]`$, $`𝐭`$ is a vector bundle of rank $`r`$ over $`\overline{G/N}`$.
On the other hand, recall that in subsection 11.1 we defined the sheaf $`𝐜`$ of Lie algebras corresponding to $`𝒞`$. Our map $`\overline{\chi }:𝒞\overline{𝒯}`$ induces a map $`d\overline{\chi }:𝐜𝐭`$ which, for simplicity, we abbreviate as $`d\chi :𝐜𝐭`$.
###### Proposition 12.5.
The map
$$d\chi :𝐜𝐭$$
is an isomorphism.
###### Proof.
The proof will consist of two steps. The first step will be a reduction to the case of $`SL(2)`$ and the second one will be a proof of the assertion for $`SL(2)`$.
Step 1.
Both $`𝐜`$ and $`𝐭`$ are vector bundles of rank $`r`$ over $`\overline{G/N}`$ and the map $`d\chi `$ is clearly an isomorphism over $`G/N`$. Since the variety $`\overline{G/N}`$ is smooth, it remains to show that $`d\chi `$ is an isomorphism on an open subset of $`\overline{G/N}`$ whose complement has codimension at least $`2`$.
It follows from Section 10.5 that such an open subset is formed by the union of the $`G`$-orbits of the images of $`i_J(\overline{M_J/N_{M_J}})`$, where $`J=\{\alpha _j\}`$ for all simple roots $`\alpha _j`$.
Therefore, by $`G`$-equivariance and by Proposition 12.3, it suffices to show that the map
$$d\chi _{M_J}:𝐜_{M_J}𝐭_{M_J}$$
is an isomorphism. This reduces us to the case when $`G`$ is a reductive group of semi-simple rank $`1`$.
Moreover, the statement is clearly invariant under isogenies, so we may replace $`G`$ by $`Z^0(G)\times [G,G]`$. Clearly, the assertion in such a case is equivalent to the one for $`[G,G]`$, which in turn can be replaced by $`SL(2)`$.
Step 2.
For $`G=SL(2)`$, the variety $`\overline{G/N}`$ can be identified with $`^2`$ in such a way that the sheaf $`𝐜`$ goes over to $`𝒪(1)`$. Moreover, $`\overline{G/T}\overline{G/N}`$ can be identified with the $`S_2`$-cover $`\pi :^1\times ^1^2`$.
To prove the assertion, it is enough to show that $`𝐭`$ has degree $`1`$, since any non-zero map between two line bundles of the same degree is automatically an isomorphism.
By definition, $`𝐭`$ is the $`𝒪(^2)`$-module of anti-invariants of $`S_2`$ in $`\pi _{}(𝒪(^1\times ^1))`$. Therefore,
$$𝐭det(\pi _{}(𝒪(^1\times ^1)))𝒪(1).$$
#### 12.6.
Now we will check that the map $`\chi `$ induces an isomorphism between $``$-points of $`𝒞`$ and $`𝒯`$. Evidently, this assertion, combined with Proposition 12.5 and Proposition 11.2, implies Theorem 11.6.
Let $`𝔞\overline{G/N}`$ be the centralizer of a regular element $`x𝔤`$. As we saw in the proof of Proposition 11.2, on the one hand, the fiber of $`𝒞`$ at $`𝔞=Z_𝔤(x)`$ can be identified with $`Z_G(x)`$. On the other hand, the fiber of $`𝒯`$ at $`𝔞`$ can be identified with $`\mathrm{Hom}_W(l^x,T)^{}`$, where $`l^x`$ is the fixed point scheme of the vector field induced by $`x`$ on $`l`$ and the super-script $``$ corresponds to the (\**) condition in the definition of $`𝒯`$.
Let $`x=x^{ss}+x^{nil}`$ be the Jordan decomposition of $`x`$. We can assume that $`Z_𝔤(x^{ss})`$ is a standard Levi subalgebra $`𝔪`$ and $`x^{nil}`$ is a regular nilpotent element in $`𝔪`$. Using Proposition 12.3, we can replace $`G`$ by $`M`$ and hence we can assume that $`x^{ss}`$ is a central element in $`𝔤`$.
There are natural embeddings $`Z(G)\times \overline{G/N}𝒞`$ and $`Z(G)\times \overline{G/N}𝒯`$, which make the diagram
$$\begin{array}{ccc}Z(G)\times \overline{G/N}& & 𝒞\\ \mathrm{id}& & \chi & & \\ Z(G)\times \overline{G/N}& & \overline{𝒯}\end{array}$$
commute.
###### Proposition 12.7.
Let $`x`$ be a regular nilpotent element and let
$`Z_G(x)=Z_G(x)^{ss}\times Z_G(x)^{nil}`$
$`\mathrm{Hom}_W(l^x,T)^{}=\mathrm{Hom}_W(l^x,T)^{ss,}\times \mathrm{Hom}_W(l^x,T)^{nil,}`$
be the Jordan decompositions of the fibers of $`𝒞`$ and $`𝒯`$ at $`Z_𝔤(x)`$. Then the embedding of $`Z(G)`$ induces isomorphisms:
$$Z(G)Z_G(x)^{ss}\text{ and }Z(G)\mathrm{Hom}_W(l^x,T)^{ss,}.$$
It is clear, first of all, that this proposition implies the theorem. Indeed, it is enough to show that $`\chi `$ induces an isomorphism $`Z_G(x)^{nil}\mathrm{Hom}_W(l^x,T)^{nil,}`$. But since these groups are unipotent, our assertion follows from the corresponding assertion on the Lie-algebra level, which has been proven before.
###### Proof.
The fact that $`Z(G)Z_G(x)^{ss}`$ is an immediate consequence of the fact that in a group of adjoint type centralizers of regular elements are connected.
To prove that $`Z(G)\mathrm{Hom}_W(l^x,T)^{ss,}`$, let us observe that if $`x`$ is a regular nilpotent element, $`l^x`$ is a local non-reduced scheme. Its closed point, viewed as a point of $`\overline{G/T}`$, belongs to the intersection of all the $`D^\alpha `$’s.
Let $`\mathrm{Hom}_W(l^x,T)_1`$ be the sub-group of $`\mathrm{Hom}_W(l^x,T)`$ which corresponds to maps $`l^xT`$ that send the closed point of $`l^x`$ to the identity in $`T`$. Clearly, $`\mathrm{Hom}_W(l^x,T)_1`$ is unipotent and $`\mathrm{Hom}_W(l^x,T)\mathrm{Hom}_W(l^x,T)_1\times T^W`$ is the Jordan decomposition of $`\mathrm{Hom}_W(l^x,T)`$.
The proof is concluded by the observation that
$$Z(G)=\{tT^W|\alpha (t)=1,\alpha \mathrm{\Delta }\},$$
which is exactly the (\**) condition.
### 13. Proof of Theorem 11.8
#### 13.1.
We will need an additional property of the isomorphism $`\chi `$.
By definition, we have a canonical $`W`$-equivariant map
$$𝐭\underset{\overline{G/N}}{\times }\overline{G/T}𝔱,$$
hence we obtain a map $`𝐭𝔱/W`$.
###### Lemma 13.2.
The above map coincides with the composition
$$𝐭\stackrel{\chi ^1}{}𝐜𝔤\times \overline{G/N}𝔤𝔱/W,$$
where the last arrow is the Chevalley map.
The proof follows from the fact that the two maps coincide over $`G/N`$.
#### 13.3.
Since $`G\times \overline{G/N}`$ is smooth, to prove the theorem, we need to show that any map $`S`$ can be lifted, locally in the étale topology, to a map $`SG\times \overline{G/N}`$.
Thus, let $`𝔞^1`$ and $`𝔞^2`$ be two $`S`$-points of $`\overline{G/N}`$ and let $`\nu :\stackrel{~}{S}^1\stackrel{~}{S}^2`$ be an isomorphism between the corresponding cameral covers. The maps $`𝔞^i`$ give rise to vector subbundles $`𝐜_S^i𝔤𝒪_S`$, and Theorem 11.6 implies that
$$𝐜_S^i\mathrm{Hom}_{W,𝒪_S}(𝔱^{},𝒪_{\stackrel{~}{S}^i}),i=1,2.$$
Therefore, the data of $`\nu `$ defines an isomorphism of vector bundles $`\nu ^{}:𝐜_S^1𝐜_S^2`$.
By Proposition 1.3 we can find a section $`x_S^1𝐜_S^1`$, such that $`𝐜_S^1=Z_𝔤(x_S^1)`$. Let $`x_S^2𝐜_S^2`$ be the image of $`x_S^1`$ under $`\nu ^{}`$. By making the choice of $`x_S^1`$ sufficiently generic, we can assume that $`x_S^2`$ is regular, i.e. that $`𝐜_S^2=Z_𝔤(x_S^2)`$.
Consider $`x_S^i,i=1,2`$ as maps $`S𝔤_{reg}`$. Lemma 13.2 implies that their compositions with the Chevalley map
$$S\stackrel{x_S^i}{}𝔤_{reg}𝔱/W$$
coincide. Now, we have the following general assertion that follows from smoothness of the Chevalley map restricted to $`𝔤_{reg}`$:
###### Lemma 13.4.
The adjoint action map $`G\times 𝔤_{reg}𝔤_{reg}\underset{𝔱/W}{\times }𝔤_{reg}`$ is smooth and surjective.
Therefore, locally there exists a map $`g_S:SG`$ that conjugates $`x_S^1`$ to $`x_S^2`$. Then this map conjugates $`𝐜_S^1`$ to $`𝐜_S^2`$, which is what we had to prove.
#### 13.5. Complements
We conclude this section by two remarks regarding the assertions of Theorem 11.6 and Theorem 11.8.
First, let us fix a $``$-point $`𝔞\overline{G/N}`$ and let $`\phi :l^𝔞𝔱`$ be a $`W`$-equivariant map, which according to Theorem 11.6, is the same as an element $`x_\phi 𝔞=𝐜_𝔞`$. One may wonder: how can one express the condition that $`x_\phi `$ is a regular element of $`𝔞`$ in terms of $`\phi `$?
###### Lemma 13.6.
The necessary and sufficient condition for $`x_\phi `$ to be a regular element of $`𝔞`$ is that $`\phi :l^𝔞𝔱`$ is a scheme-theoretic embedding.
###### Proof.
First, one easily reduces the assertion to the case when $`𝔞`$ is the centralizer of a regular nilpotent element, which we will assume.
In this case, $`𝔞`$ entirely consists of nilpotents elements. Let $`\mathrm{St}_G(𝔞)`$ be the normalizer of $`𝔞`$. Since the nilpotent locus in $`𝔤_{reg}`$ is a single $`G`$-orbit, we obtain that $`𝔞𝔤_{reg}`$ is a single $`\mathrm{St}_G(𝔞)`$-orbit.
Thus, let $`\varphi `$ be an embedding. To show that $`x_\phi `$ is regular, it is enough to show that its centralizer in $`\mathrm{St}_G(𝔞)`$ coincides with $`𝒞_𝔞`$.
By Proposition 11.10, the quotient $`\mathrm{St}_G(𝔞)/𝒞_𝔞`$ maps isomorphically to the group of $`W`$-equivariant automorphisms of $`l^𝔞`$. If for some $`n\mathrm{St}_G(𝔞)`$ we have $`\mathrm{Ad}_n(x_\phi )=x_\phi `$, then $`n`$ acts trivially on $`l^𝔞`$, since $`\varphi `$ is an embedding. Hence, $`n𝒞_𝔞`$.
To prove the implication in the other direction, let us observe that $`𝔞𝔤_{reg}`$ is the only $`\mathrm{St}_G(𝔞)`$-invariant open subset of $`𝔞`$ consisting of regular elements only. However, the locus of $`\phi `$ that are embeddings is clearly such a subset.
Secondly, let us see how Corollary 11.11 is related to Proposition 1.8:
Let $`𝔞_1`$ and $`𝔞_2`$ be two $``$-points of $`\overline{G/N}`$. Corollary 11.11 says that they are $`G`$-conjugate if and only if $`\pi ^1(𝔞_1)\pi ^1(𝔞_2)`$ as $`W`$-schemes. The condition of Proposition 1.8 is seemingly weaker (but in fact, equivalent): it implies that $`𝔞_1`$ and $`𝔞_2`$ are $`G`$-conjugate if and only if $`(\pi ^1(𝔞_1))_{red}(\pi ^1(𝔞_2))_{red}`$ as $`W`$-schemes.
## Part IV Proofs of the main results
### 14. Proof of Theorem 4.4
#### 14.1.
We are going to deduce our theorem from Theorem 11.6 and Theorem 11.8 combined with the following “abstract nonsense” observation:
###### Lemma 14.2.
Let $`𝒬`$ be a sheaf of categories on $`\mathrm{Sch}_{et}(X)`$, and $`𝒜`$ be a sheaf of abelian groups on $`\mathrm{Sch}_{et}(X)`$. Suppose that for every $`(UX)\mathrm{Sch}_{et}(X)`$ and every $`C𝒬(U)`$, we are given an isomorphism $`\mathrm{Aut}_{𝒬(U)}(C)𝒜(U)`$ such that the following conditions hold:
(0) There exists a covering $`UX`$ such that $`𝒬(U)`$ is non-empty.
(1) If $`C_1C_2`$ is an isomorphism between two objects in $`𝒬(U)`$, then the induced isomorphism $`\mathrm{Aut}_{𝒬(U)}(C_1)\mathrm{Aut}_{𝒬(U)}(C_2)`$ is compatible with the identification of both sides with $`𝒜(U)`$.
(2) If $`f:U^{}U`$ is a morphism in $`\mathrm{Sch}_{et}(X)`$ and $`C𝒬(U)`$, then the map
$$f_𝒬^{}:\mathrm{Aut}_{𝒬(U)}(C)\mathrm{Aut}_{𝒬(U^{})}(f_𝒬^{}(C))$$
is compatible with the restriction map $`𝒜(U)𝒜(U^{})`$.
(3) For any $`U\mathrm{Sch}_{et}(X)`$ and any two $`C_1,C_2𝒬(U)`$, there exist a covering $`f:U^{}U`$ such that the objects $`f_𝒬^{}(C_1)`$ and $`f_𝒬^{}(C_2)`$ of $`𝒬(U^{})`$ are isomorphic.
Then $`𝒬`$ has a canonical structure of a gerbe over $`\mathrm{Tors}_𝒜`$.
#### 14.3.
We claim that $`\mathrm{Higgs}_{\stackrel{~}{X}}`$ satisfies the conditions of this lemma. Condition (0) is a tautology: locally the cameral cover $`\stackrel{~}{X}X`$ is induced from the universal one by means of a map $`X\overline{G/N}`$.
Let $`(E_G,\sigma ,t)`$ be an object of $`\mathrm{Higgs}_{\stackrel{~}{X}}(U)`$. We must construct an isomorphism
$$\mathrm{Aut}_{\mathrm{Higgs}_{\stackrel{~}{X}}(U)}(E_G,\sigma ,t)T_{\stackrel{~}{X}}(U).$$
Let us first assume that $`E_G`$ is trivialized and our Higgs bundle corresponds to a map $`U\overline{G/N}`$ such that $`\stackrel{~}{U}\stackrel{𝑡}{}\overline{G/T}\underset{\overline{G/N}}{\times }U`$.
In this case, an automorphism of $`(E_G,\sigma )`$ as an object of $`\mathrm{Higgs}(U)`$ is the same as a map $`U\mathrm{St}`$ (cf. Section 11.9) that covers the given map $`X\overline{G/N}`$. Now, Proposition 11.10 implies that this automorphism belongs to $`\mathrm{Aut}_{\mathrm{Higgs}_{\stackrel{~}{X}}(U)}(E_G,\sigma ,t)`$ if and only if the above map factors as $`U𝒞\mathrm{St}`$.
Now we apply Theorem 11.6 which says that $`\mathrm{Hom}_{\overline{G/N}}(U,𝒞)=\mathrm{Hom}_{\overline{G/N}}(U,𝒯)=T_{\stackrel{~}{X}}(U)`$.
The fact that the map $`\chi :𝒞𝒯`$ is $`G`$-equivariant implies that our isomorphism between $`\mathrm{Aut}_{\mathrm{Higgs}_{\stackrel{~}{X}}(U)}(E_G,\sigma ,t)`$ and $`T_{\stackrel{~}{X}}(U)`$ is independent of the choice of a trivialization of $`E_G`$. In particular, by SC-1, it defines the required isomorphism for all $`E_G`$. The fact that conditions (1) and (2) are satisfied is automatic from the construction.
Finally, let us check condition (3). Let $`(E_G^1,\sigma ^1,t^1)`$ and $`(E_G^2,\sigma ^2,t^2)`$ be two objects of $`\mathrm{Higgs}_{\stackrel{~}{U}}(U)`$. Without restricting the generality we can assume that both $`E_G^1`$ and $`E_G^2`$ are trivialized.
In this case, the data of $`(\sigma ^1,\sigma ^2,t^1(t^2)^1)`$ defines a $`U`$-point of the scheme $``$. By Theorem 11.8 we can locally find a map $`g_U:UG`$ which conjugates $`(\sigma ^1,t^1)`$ to $`(\sigma ^2,t^2)`$. We can regard $`g_U`$ as a gauge transformation, i.e. a map $`E_G^1E_G^2`$, which defines an isomorphism between $`(E_G^1,\sigma ^1,t^1)`$ and $`(E_G^2,\sigma ^2,t^2)`$.
Thus, Theorem 4.4 is proved.
### 15. Proof of Theorem 7.3
#### 15.1.
It remains to prove Theorem 6.4. In this section we will prove Theorem 7.3 which takes care of the universal situation.
#### 15.2. Step 1
First we show that our map $`s_i^{}(_{can})_{can}`$ is an isomorphism off $`D^{\alpha _i}`$. To do that let us analyze more closely the situation described in Section 10.5.
Let $`\mathrm{\Delta }_J\mathrm{\Delta }`$ be a root subsystem and let $`M=M_J`$ be the corresponding standard Levi subgroup. Let $`l_M`$ denote the flag variety of $`M`$ and $`B_M=BM`$, $`U_M=UM`$.
It is well-known that there exists a canonical closed embedding $`W_M\backslash W\times l_Ml`$:
A point $`𝔟^{}l`$ belongs to $`w\times l_M`$, if and only if $`𝔟^{}`$ is in relative position $`w`$ with respect to $`P=P_J`$ (this makes sense, as $`P`$-orbits in $`l`$ are parametrized exactly by $`W_M\backslash W`$) and $`𝔟^{}𝔪`$ is a Borel subalgebra in $`M`$.
Consider the restriction of the canonical $`T`$-bundle $`_l`$ to $`W_M\backslash W\times l_M`$. It is easy to see that its further restriction to the connected component $`1\times l_M`$ identifies with $`_{l_M}`$.
Let $`wW`$ be a minimal representative of its coset in $`W_M\backslash W`$. The action of $`w`$ defines a map $`1\times l_Mw^1\times l_M`$. Let us consider the pull-back $`w^{}(_l|_{w\times l_M})`$ as a $`T`$-bundle on $`1\times l_M=l_M`$. Let $`\stackrel{~}{w}N`$ be an element that projects to $`wW`$.
###### Lemma 15.3.
We have a canonical $`M`$-equivariant isomorphism
$$w^{}(_l|_{w\times l_M})_{l_M}.$$
###### Proof.
Both $`w^{}(_l|_{w\times l_M})`$ and $`_{l_M}`$ are $`M`$-equivariant $`T`$-bundles on $`l_M`$. To prove that they are isomorphic, we must show that the two homomorphisms $`BMT`$ corresponding to the base point $`𝔟l_M`$ coincide.
However, this follows from the fact that $`w^1\times 𝔟=\mathrm{Ad}_{\stackrel{~}{w}^1}(𝔟)`$, which is true since $`w`$ is minimal.
#### 15.4.
Let $`w`$ be as above. Consider the map
$$\overline{M/T}\stackrel{\text{Proposition }\text{10.6}}{}\overline{G/T}\stackrel{𝑤}{}\overline{G/T}l.$$
The fact that $`\mathrm{Ad}_{\stackrel{~}{w}^1}(B)M=B_M`$ implies that the above map coincides with
$$\overline{M/T}1\times l_M\stackrel{𝑤}{}w\times l_Ml.$$
Therefore, from Lemma 15.3 we obtain an isomorphism
$$\gamma _{can}^{}(\stackrel{~}{w}):w^{}(_{can})|_{\overline{M/T}}_{can}|_{\overline{M/T}}.$$
Moreover, it is easy to see that the above isomorphism is induced by the restriction to $`\overline{M/T}`$ of the (meromorphic) isomorphism $`\gamma _{can}(\stackrel{~}{w})`$. In particular, the a priori meromorphic isomorphism $`\gamma _{can}(\stackrel{~}{w})`$ is regular on $`\overline{M/T}`$.
Let us now go back to the situation of the theorem. We must check that the meromorphic map $`s_i^{}(_{can})_{can}`$ has no poles along $`D^\alpha `$ if $`\alpha \alpha _i`$. Choose a minimal Levi subgroup $`M_j`$ such that $`w(\alpha _j)=\alpha `$ for some $`wW`$.
Then the fact that $`\alpha \alpha _i`$ implies that both $`w`$ and $`s_iw`$ are minimal representatives of the corresponding cosets in $`W/s_j`$. Then the above discussion shows that $`s_i^{}(_{can})_{can}`$ has no poles on $`w\times \overline{M_j/T}`$.
This proves what we need, since the $`G`$-orbit of $`w\times \overline{M_j/T}`$ contains an open part of $`D^\alpha `$ (cf. Proposition 1.8).
#### 15.5. Step 2
Thus, we have shown that the poles of the map $`s_i^{}(_{can})_{can}`$ can occur only on $`D^{\alpha _i}`$. Let $`M_i`$ be the corresponding minimal Levi subgroup. As we have seen before, there is a natural embedding $`\overline{M_i/T}\overline{G/T}`$, and $`_{can}`$ restricts to the corresponding $`T`$-bundle on $`\overline{M_i/T}`$.
Since $`s_i^{}(_{can})_{can}`$ is $`G`$-equivariant, to determine the contribution of the divisor $`D^{\alpha _i}`$, it is enough to perform the corresponding calculation for $`M_i`$. The latter case easily reduces to $`SL(2)`$.
For $`SL(2)`$, $`\overline{G/T}^1\times ^1`$ and $`_{can}𝒪(1)𝒪`$. Moreover, $`1S_2=W`$ acts on $`\overline{G/T}`$ by swapping the two $`^1`$ factors, with the fixed-point locus $`\overline{G/T}^1`$ being the diagonal $`^1`$. Hence, $`(1)^{}(_{can})𝒪𝒪(1)`$.
Therefore, we have a meromorphic map between $`𝒪𝒪(1)`$ and $`𝒪(1)𝒪`$, which is allowed to have zeroes and poles only on the diagonal. Then it must have a zero of order $`1`$, by degree considerations.
### 16. Proof of Theorem 6.4
#### 16.1. The natural functor
Finally, we are ready to complete the proof of the main result. First, we claim that there is a natural functor $`\mathrm{{\rm Y}}:\mathrm{Higgs}_{\stackrel{~}{X}}\mathrm{Higgs}_{\stackrel{~}{X}}^{}`$:
Let $`(E_G,\sigma ,t)`$ be an object of $`\mathrm{Higgs}_{\stackrel{~}{X}}(U)`$, where $`\sigma :E_G\overline{G/N}`$ is a $`G`$-equivariant map. We can pull-back the universal object of $`\mathrm{Higgs}_{\overline{G/T}}^{}(\overline{G/N})`$ (cf. Section 7) and obtain a $`G`$-equivariant object of $`\mathrm{Higgs}_{\stackrel{~}{E}_G}^{}(E_G)`$, where $`\stackrel{~}{E}_G`$ is the induced cameral cover of $`E_G`$.
By descent, it gives rise to an object of $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(U)`$ and this assignment is clearly a functor between sheaves of categories.
The key fact now is that $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}`$ is also a gerbe over $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}`$. Condition (0) of Lemma 14.2 follows from the mere existence of the functor $`\mathrm{{\rm Y}}`$ and the fact that $`\mathrm{Higgs}_{\stackrel{~}{X}}`$ satisfies condition (0).
Let $`(,\gamma ,\beta _i)`$ be an object of $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(U)`$. We must identify the group of its automorphisms with $`T_{\stackrel{~}{X}}(U)`$.
By definition, this group consists of $`T`$-bundle automorphisms, which respect the data of $`\gamma `$ and $`\beta _i`$. However, a $`T`$-bundle map $``$ is the same as a map $`\stackrel{~}{U}T`$ and compatibility with $`\gamma `$ implies that this map is $`W`$-equivariant. Therefore, we obtain a section of $`\overline{T}_{\stackrel{~}{X}}(U)`$. Now, compatibility with $`\beta _i`$ is exactly condition (\*). (Recall that it suffices to impose condition (\*) for one representative in every $`W`$-orbit on the set of roots. In particular, it is sufficient to impose it for simple roots only.)
It is easy to see that conditions (1) and (2) hold for the above identification of $`\mathrm{Aut}_{\mathrm{Higgs}_{\stackrel{~}{X}}^{}(U)}(,\gamma ,\beta _i)T_{\stackrel{~}{X}}(U)`$. In addition, it follows from the construction of $`\chi `$, that $`\mathrm{{\rm Y}}:\mathrm{Higgs}_{\stackrel{~}{X}}\mathrm{Higgs}_{\stackrel{~}{X}}^{}`$ respects the identifications of groups of automorphisms of objects with $`T_{\stackrel{~}{X}}(U)`$.
Assume for a moment that condition (3) of Lemma 14.2 has been checked. We claim, that this already implies Theorem 6.4, because of the following general fact:
###### Lemma 16.2.
Let $`𝒬_1`$ and $`𝒬_2`$ be two gerbes over $`\mathrm{Tors}_𝒜`$ and let $`\mathrm{{\rm Y}}:𝒬_1𝒬_2`$ be a functor between the corresponding sheaves of categories. Assume that for every $`U\mathrm{Sch}_{et}(X)`$ and $`C𝒬_1(U)`$ we have a commutative square:
$$\begin{array}{ccc}𝒜(U)& & \mathrm{Aut}_{𝒬_1(U)}(C)\\ \mathrm{id}& & \mathrm{{\rm Y}}& & \\ 𝒜(U)& & \mathrm{Aut}_{𝒬_2(U)}(\mathrm{{\rm Y}}(C)).\end{array}$$
Then $`\mathrm{{\rm Y}}`$ is an equivalence of $`\mathrm{Tors}_𝒜`$-gerbes.
#### 16.3. The homogeneous version: $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}^{}`$
It remains to prove that every two objects of $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(U)`$ are locally isomorphic. For that purpose we will introduce a sheaf of Picard categories $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}^{}`$, which will be the “homogeneous” version of $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)`$.
Objects of $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}^{}(U)`$ are triples
$$(_0,\gamma _0,\beta _{i,0}),$$
where $`(_0,\gamma _0)`$ is a strongly $`W`$-equivariant $`T`$-bundle on $`\stackrel{~}{U}`$ and each $`\beta _{i,0}`$ is a trivialization of $`\alpha _i(_0)|_{D_U^{\alpha _i}}`$.
The following compatibility conditions must hold:
(1) For a simple root $`\alpha _i`$, the data of $`\gamma _0(s_i):s_i^{}(_0)_0`$ defines, after restriction to $`D_U^{\alpha _i}`$, a trivialization
$$\stackrel{ˇ}{\alpha }_i(\alpha _i(_0)|_{D_U^{\alpha _i}})\stackrel{ˇ}{\alpha }_i(𝒪_{D_U^{\alpha _i}}).$$
We need that this trivialization coincides with $`\stackrel{ˇ}{\alpha }_i(\beta _{i,0})`$.
(2) Assume that $`wW`$ conjugates a simple root $`\alpha _i`$ to another simple root $`\alpha _j`$. The pull-back of $`\beta _{j,0}`$ under $`w`$ is a trivialization of $`\alpha _i(w^{}(_0))|_{D_U^{\alpha _i}}`$, which via $`\gamma _0(w)`$ defines a trivialization of $`\alpha _i(_0)`$. Our condition is that this trivialization coincides with $`\beta _{i,0}`$.
Morphisms in $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}^{}(U)`$ are by definition maps between strongly $`W`$-equivariant $`T`$-bundles, compatible with the data of $`\beta _{i,0}`$.
If $`(_0^1,\gamma _0^1,\beta _{i,0}^1)`$ and $`(_0^2,\gamma _0^2,\beta _{i,0}^2)`$ are two objects of $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}^{}(U)`$ we can form their tensor product $`(_0^1_0^2,\gamma _0^1\gamma _0^2,\beta _{i,0}^1\beta _{i,0}^2)`$ which will be a new object of $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}^{}(U)`$. Moreover, if $`(_0,\gamma _0,\beta _{i,0})`$ is an object of $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}^{}(U)`$ and $`(,\gamma ,\beta _i)`$ is an object of $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(U)`$, we can take their tensor product and obtain another object of $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(U)`$.
It is easy to see that the above constructions define on $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}^{}`$ a structure of a sheaf of Picard categories and on $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}`$ a structure of a gerbe over it.
Therefore, to prove that every two objects of $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(U)`$ are locally isomorphic, it is enough to show that any object of $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}^{}(U)`$ is locally isomorphic to the unit object, i.e. to the one with $`_0`$ being the trivial $`T`$-bundle and $`(\gamma _0,\beta _{i,0})`$ being the tautological maps. The last assertion is equivalent to:
###### Proposition 16.4.
$`\mathrm{Tors}_{T_{\stackrel{~}{X}}}^{}`$ is equivalent as a sheaf of Picard categories to $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}`$.
We proceed to prove this Proposition by showing that any object in $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}^{}(U)`$ is locally isomorphic to the unit object.
#### 16.5. Step 1
Without restricting the generality, we can assume that $`U=X`$ and we must find an étale covering $`X^{}X`$, over which a given object $`(_0,\gamma _0,\beta _{i,0})`$ becomes isomorphic to the trivial one.
Fix a $``$-point $`xX`$. First, we will reduce our situation to the case when the ramification over $`x`$ is maximal possible, i.e. when $`x`$ belongs to the image of $`\underset{𝛼}{}D_X^\alpha `$, where the intersection is taken over all roots of $`G`$.
After an étale localization we can assume that we have a map $`X𝔱/W`$ so that $`\stackrel{~}{X}=X\underset{𝔱/W}{\times }𝔱`$. Let $`t`$ be a point in $`𝔱`$ which has the same image in $`𝔱/W`$ as $`x`$.
By conjugating $`t`$, we can assume that there exists $`JI`$ such that $`\alpha _j(t)=0`$ for $`j\mathrm{\Delta }_J`$ and $`\beta (t)0`$ for $`\beta \mathrm{\Delta }_J`$.
We have a Cartesian square
$$\begin{array}{ccc}W\stackrel{W_J}{\times }(𝔱\underset{\beta \mathrm{\Delta }_J}{}𝔱^\beta )& & 𝔱\\ & & & & \\ (𝔱\underset{\beta \mathrm{\Delta }_J}{}𝔱^\beta )/W_J& & 𝔱/W,\end{array}$$
in particular, the map $`𝔱/W_J𝔱/W`$ is étale in a neighbourhood of the image of $`t`$ in $`𝔱/W`$.
Therefore, the base change $`XX^{}:=X\underset{𝔱/W}{\times }𝔱/W_J`$ is étale in a neighbourhood of $`x`$. This reduces us to the situation, when the $`W`$-cover $`\stackrel{~}{X}`$ is induced from a $`W_J`$-cover $`\stackrel{~}{X}_J`$, i.e. $`\stackrel{~}{X}W\stackrel{W_J}{\times }\stackrel{~}{X}_J`$.
By restricting $`(_0,\gamma _0,\beta _{i,0})`$ to $`\stackrel{~}{X}_J`$ we obtain an object of $`\mathrm{Tors}_{T_{\stackrel{~}{X}_J}}^{}(X)`$ equivariant with respect to $`W_J`$. Moreover, it is easy to see that this establishes an equivalence between $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}^{}`$ and $`\mathrm{Tors}_{T_{\stackrel{~}{X}_J}}^{}`$, thereby reducing us to the situation when $`\mathrm{\Delta }_J=\mathrm{\Delta }`$.
#### 16.6. Step 2
According to Step 1, we may assume that there exists a unique geometric point $`\stackrel{~}{x}\stackrel{~}{X}`$ over $`x`$. To prove the assertion of the proposition, we can replace $`X`$ by the spectrum of the local ring of $`X`$ at $`x`$. In this case all the $`D_X^\alpha `$’s and $`\stackrel{~}{X}`$ are local too.
Let us choose a trivialization of our line bundle $`_0`$, subject only to the condition that it is compatible with the data of $`\beta _{i,0}`$ at $`\stackrel{~}{x}`$ for every simple root $`\alpha _i`$. We must show that this trivialization can be modified so that it will be compatible with the structure on $`_0`$ of a $`W`$-equivariant $`T`$-bundle, i.e. with the data of $`\gamma _0`$. (The argument given below mimics the proof of Proposition 9.5).
The discrepancy between our initial trivialization and $`\gamma _0`$ is given by a $`1`$-cocycle $`\mu :W\mathrm{Hom}(\stackrel{~}{X},T)`$.
The evaluation at $`\stackrel{~}{x}`$ gives rise to a surjection of $`W`$-modules: $`\mathrm{Hom}(\stackrel{~}{X},T)T`$. Thus we obtain a short exact sequence:
$$0K\mathrm{Hom}(\stackrel{~}{X},T)T0,$$
where $`K`$ consists of maps $`\stackrel{~}{X}T`$ which have value $`1`$ at $`\stackrel{~}{x}`$.
Now, our condition on the trivialization (i.e. its compatibility with $`\beta _{i,0}`$) and condition (1) in the definition of $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}^{}`$ imply that $`\mu (s_i)K`$ for every simple reflection $`s_i`$. Hence, $`\mu `$ takes values in $`K`$.
However, since $`\stackrel{~}{X}`$ is local, $`K`$ is torsion-free and divisible! Hence, $`H^1(W,K)=0`$.
Therefore, we can choose a trivialization of $`_0`$ which respects the $`W`$-equivariant structure and the data of $`\beta _{i,0}`$ at $`\stackrel{~}{x}`$. But this implies that it is compatible with the data of $`\beta _{i,0}`$ on the entire $`D_X^{\alpha _i}`$, $`iI`$.
Indeed, a possible discrepancy takes values in $`\pm 1`$, and its value is constant along every connected component of $`D_X^{\alpha _i}`$. However, by construction, each $`D_X^{\alpha _i}`$ is local with $`\stackrel{~}{x}`$ being its unique closed point.
The proof of Proposition 16.4, and hence of Theorem 6.4, is now complete.
#### 16.7. Variant
As was the case for $`\mathrm{Higgs}^{}`$, we can give a much simplified description of $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}^{}`$ in case our group $`G`$ does not have an $`\mathrm{SO}(2n+1)`$ direct factor. In this case the data consists of a strongly equivariant $`T`$-bundle $`(L_0,\gamma _0)`$, such that for a simple root $`\alpha _i`$ and a weight $`\lambda `$ orthogonal to the corresponding coroot, the isomorphism $`\lambda (s_i^{}(L))|_{D^{\alpha _i}}\lambda (L)|_{D^{\alpha _i}}`$ induced by $`\gamma (s_i)`$ coincides with the tautological one.
## Part V Some applications
The point of our abstract notion of a Higgs bundle, as defined in Section 2, is that it provides a uniform approach to the analysis of various more concrete objects. In the final sections we illustrate the applications to Higgs bundles with values in a line bundle or in an elliptic fibration.
### 17. Higgs bundles with values
#### 17.1.
In Section 2 we defined a Higgs bundle over a scheme $`X`$ to be a pair $`(E_G,\sigma )`$, where $`\pi :E_GX`$ is a principal $`G`$-bundle over $`X`$, and $`\sigma `$ is a $`G`$-equivariant map $`\sigma :E_G\overline{G/N}`$. We noted there that on a given $`G`$-bundle $`E_G`$, a Higgs bundle is specified by a vector subbundle $`𝐜_X`$ of $`𝔤_{E_G}`$ whose fibers are regular centralizers. (Recall that $`𝔤_{E_G}:=E_G\underset{𝐺}{\times }𝔤`$ is the adjoint bundle of $`E_G`$.) In subsection 11.1 we defined the universal centralizer $`𝐜𝔤\times \overline{G/N}`$, corresponding to the universal Higgs bundle over $`\overline{G/N}`$. The family of centralizers $`𝐜_X`$ of a general Higgs bundle $`(E_G,\sigma )`$ over $`X`$ is related to the universal $`𝐜`$ by: $`\pi ^{}𝐜_X=\sigma ^{}𝐜`$, an equality of vector subbundles of the trivial bundle $`𝔤\times E_G`$ on the total space of $`E_G`$. We recall also that by Theorem 12.5, $`𝐜`$ is isomorphic to $`𝐭=\mathrm{Lie}(𝒯)`$.
Let $`K`$ be a line bundle on our base $`X`$. In the literature, the most common notion of a Higgs bundle is:
###### Definition 17.2.
A $`K`$-valued Higgs bundle on $`X`$ is a pair $`(E_G,s)`$, where $`E_G`$ is a principal $`G`$-bundle on $`X`$ and $`s`$ is a section of $`𝔤_{E_G}K`$.
The section $`s`$ of $`𝔤_{E_G}K`$ is called regular at a point $`xX`$ if the corresponding local section of $`𝔤_{E_G}`$ determined by some (hence, any) trivialization of $`K`$ at $`x`$ is regular. We work instead with the following more general notion, which is also better adapted to our setup.
###### Definition 17.3.
A regularized $`K`$-valued Higgs bundle on $`X`$ is a triple $`(E_G,\sigma ,s)`$, with $`(E_G,\sigma )`$ a Higgs bundle on $`X`$ and $`s`$ a section of $`𝐜_XK`$, where $`𝐜_X`$ is the regular centralizer subbundle of the adjoint bundle $`𝔤_{E_G}`$ determined by $`\sigma `$.
#### 17.4. Regular vs. regularized
A regularized $`K`$-valued Higgs bundle $`(E_G,\sigma ,s)`$ on $`X`$ clearly determines the unique $`K`$-valued Higgs bundle $`(E_G,s)`$ on $`X`$. Conversely, if the section $`s`$ of $`𝔤_{E_G}K`$ is everywhere regular, then we can recover $`𝐜_X𝔤_{E_G}`$ as the centralizer of $`s`$, which defines a regularized $`K`$–valued Higgs bundle. When $`s`$ is generically regular, the family $`𝐜_X`$ of centralizers is still unique, if it exists. In general, when $`s`$ is not necessarily regular, our definition adds to the pair $`(E_G,s)`$ a choice of a regular centralizer containing $`s`$.
We want to establish the following result:
###### Theorem 17.5.
A regularized $`K`$-valued Higgs bundle on $`X`$ is the same as a triple:
(a) A cameral cover $`\stackrel{~}{X}X`$,
(b) A $`W`$-equivariant map $`v:\stackrel{~}{X}𝔱K`$ (of schemes over $`X`$).
(c) An object of $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)`$,
###### Proof.
Given Theorem 6.4, it remains to show that the data (b) of a $`W`$-equivariant “value” map $`\stackrel{~}{X}𝔱K`$ is the same as the data of a section $`s`$ of $`𝐜_XK`$. And indeed, giving such a section $`s:X𝐜_XK`$ is equivalent to giving a $`G`$-equivariant section $`\stackrel{~}{s}:E_G\sigma ^{}𝐜K`$ of the pullback $`\pi ^{}𝐜_XK=\sigma ^{}𝐜K`$ over $`E_G`$, cf. 17.1 above. By Theorem 12.5, this is the same as a $`G`$-equivariant section $`\stackrel{~}{s}^{}:E_G\sigma ^{}𝐭K`$. Now by the definition of $`𝒯`$ (cf. subsection 11.3), $`\mathrm{Hom}_{\overline{G/N}}(E_G,𝐭)=\mathrm{Hom}_W(\stackrel{~}{E}_G,𝔱)`$. Here $`\stackrel{~}{E}_G:=E_G\times _X\stackrel{~}{X}`$ is the $`G`$-equivariant cameral cover of $`E_G`$ associated to the Higgs bundle on $`E_G`$ which is $`\pi ^{}`$ of our given Higgs bundle $`(E_G,\sigma )`$ on $`X`$. The section $`\stackrel{~}{s}^{}`$, and hence also our original section $`s`$, are therefore equivalent to a $`W`$-equivariant map of $`X`$-schemes $`\overline{s}:\stackrel{~}{E}_G𝔱K`$ which is also $`G`$-invariant. But this is the same as a $`W`$-equivariant map of $`X`$-schemes $`v:\stackrel{~}{X}𝔱K`$, as claimed. ∎
Note that in the data $`(E_G,\sigma ,s)`$, the section $`s:X𝐜_XK`$ is regular if and only if the corresponding map $`v`$ is an embedding. This follows from Lemma 13.6. So we have:
###### Corollary 17.6.
A regular $`K`$-valued Higgs bundle on $`X`$ is the same as a triple:
(a) A cameral cover $`\stackrel{~}{X}X`$,
(b) A $`W`$-equivariant embedding $`v:\stackrel{~}{X}𝔱K`$ (of schemes over $`X`$).
(c) An object of $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)`$,
#### 17.7. The Hitchin map
To conclude our discussion of $`K`$-valued Higgs bundles, let us note that the data (a) and (b) in the above theorem can be assembled into what can be called “a point of the Hitchin base”.
Assume that $`X`$ is proper, and let $`𝐁(X,K)`$ denote the algebraic stack which classifies the data (a) and (b) of Theorem 17.5. I.e., for a scheme $`S`$, $`\mathrm{Hom}(S,𝐁(X,K))`$ is the category of pairs $`(\stackrel{~}{X}_S,v:\stackrel{~}{X}_S𝔱K)`$, where $`\stackrel{~}{X}_S`$ is a cameral cover of $`S\times X`$, and $`v`$ is a $`W`$-equivariant morphism of $`X`$-schemes.
On the other hand, let $`\mathrm{𝐇𝐢𝐠𝐠𝐬}(X,K)`$ denote the algebraic stack of all regularized $`K`$-valued Higgs bundles on $`X`$. The Hitchin map $`h:\mathrm{𝐇𝐢𝐠𝐠𝐬}(X,K)𝐁(X,K)`$ sends a regularized $`K`$-valued Higgs bundle $`(E_G,\sigma ,s)`$ given by data (a),(b) and (c) to the point of the Hitchin base given by data (a) and (b).
###### Corollary 17.8.
The fibers of the Hitchin map $`h:\mathrm{𝐇𝐢𝐠𝐠𝐬}(X,K)𝐁(X,K)`$ can be identified (as categories) with $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)`$. By Corollary 4.6, the set of isomorphism classes of objects of this fiber is a torsor over the abelian group $`H^1(X,T_{\stackrel{~}{X}})`$, and the torsor class is given in Theorem 6.4.
Note that our description of the fiber of the Hitchin map is independent of the line bundle $`K`$.
#### 17.9.
Let now $`\mathrm{𝐇𝐢𝐭𝐜𝐡}(X,K)`$ denote the scheme of sections of the fibration $`(𝔱K)/WX`$. In fact, $`\mathrm{𝐇𝐢𝐭𝐜𝐡}(X,K)`$ is non-canonically isomorphic to an affine space.
The relation between $`\mathrm{𝐇𝐢𝐭𝐜𝐡}(X,K)`$ and $`𝐁(X,K)`$ is similar in some respects to the relation between the vector space $`𝔱/W`$ parametrizing semisimple adjoint orbits in the Lie algebra $`𝔤`$ and the stack $`𝔤/G`$ of all $`G`$-orbits in $`𝔤`$. In both cases, there is an open embedding of the variety into the stack, and there is a retraction of the stack onto the variety which is the identity on the variety.
In our case, the retraction $`r:𝐁(X,K)\mathrm{𝐇𝐢𝐭𝐜𝐡}(X,K)`$ associates to $`v:\stackrel{~}{X}𝔱K`$ the corresponding map $`X(𝔱K)/W`$. As for the open embedding $`i:\mathrm{𝐇𝐢𝐭𝐜𝐡}(X,K)𝐁(X,K)`$: starting with $`X(𝔱K)/W`$, we recover $`\stackrel{~}{X}`$ as
$$\stackrel{~}{X}:=X\underset{(𝔱K)/W}{\times }(𝔱K),$$
and $`v:\stackrel{~}{X}𝔱K`$ is the second projection.
Obviously, the image $`i(\mathrm{𝐇𝐢𝐭𝐜𝐡}(X,K))𝐁(X,K)`$ is the open substack corresponding to the condition that the map $`\stackrel{~}{X}𝔱K`$ is an embedding. By Corollary 17.8, the preimage of $`\mathrm{𝐇𝐢𝐭𝐜𝐡}(X,K)𝐁(X,K)`$ under the Hitchin map is exactly the open substack of regular $`K`$-valued Higgs bundles. Let $`R`$ denote some regularized $`K`$-valued Higgs bundle on $`X`$. Note that the image $`h(R)𝐁(X,K)`$ determines whether $`R`$ is regular. A point in $`\mathrm{𝐇𝐢𝐭𝐜𝐡}(X,K)`$, on the other hand, can be the image of both regular and irregular $`R`$’s.
#### 17.10. Variant
Definition 17.3, Theorem 17.5 and Corollary 17.8, remain unchanged if we allow $`K`$ to be a vector bundle, as is done in where $`K=\mathrm{\Omega }_X^1`$ is the cotangent bundle. In Definition 17.2, on the other hand, commutativity is not built in, so we must impose it by hand: the components of the section $`s`$, with respect to any local decomposition of $`K`$ as a sum of line bundles, must commute with each other. Equivalently, the bracket of $`s`$ with itself, interpreted as a section of $`𝔤_{E_G}^2K`$, must vanish.
### 18. Elliptic fibrations
Let $`f:YX`$ be a projective, flat, dominant morphism with integral (that is, reduced and irreducible) fibers. Eventually we will specialize this to the case of an elliptic fibration, but for now we will work with the general situation. We want to describe an application of our results to the study of regularized $`G`$-bundles on $`Y`$ in terms of data on the base $`X`$ and along the (eventually, elliptic) fibers.
By a regularization of a $`G`$-bundle $`E_G`$ on $`Y`$ we mean a reduction of its structure group along each fiber to some regular centralizer. In other words, we want a Higgs bundle $`(E_G,\sigma )`$ on $`Y`$ whose group scheme of centralizers $`𝒞_Y`$ (equivalently, its cameral cover $`\stackrel{~}{Y}Y`$) is the pullback of some group scheme of centralizers $`𝒞_X`$ on $`X`$ (respectively, of a cameral cover $`\stackrel{~}{X}X`$). More precisely:
###### Definition 18.1.
A regularized $`G`$-bundle on $`Y`$ consists of the data $`(\stackrel{~}{X},E_G,\sigma )`$, with $`\stackrel{~}{X}X`$ a cameral cover of $`X`$, and $`(E_G,\sigma )\mathrm{Higgs}_{\stackrel{~}{Y}}(Y)`$ a Higgs bundle on $`Y`$ with cameral cover $`\stackrel{~}{Y}:=\stackrel{~}{X}\underset{𝑋}{\times }Y`$.
In the case of an elliptic fibration, there is a natural notion of what it means for a bundle (on $`Y`$) to be regular above a point (of $`X`$). In analogy with the situation for $`K`$-valued Higgs bundles considered in Subsection 17.4, “most” $`G`$-bundles on an elliptic curve are indeed regular, and a regular bundle has a unique regularization. We review these well-known facts below.
#### 18.2.
In general, our current situation is the analogue of Higgs bundles with values, in which we replace the bundle $`K`$ of values from Section 17 by the relative Picard scheme $`\mathrm{Pic}(Y/X)`$. The tensor product $`𝔱_{}K`$ can be identified with $`\mathrm{\Lambda }_{}K`$, so we take its analogue to be $`\mathrm{\Lambda }_{}\mathrm{Pic}(Y/X)=:\mathrm{Bun}_T(Y/X).`$ (Here $`\mathrm{\Lambda }`$ is the lattice of coweights.) Similarly, we will need the analogue of $`𝐜_XK`$. This is the sheaf of groups $`\mathrm{Tors}_{Y/X}:=\mathrm{Tors}_{𝒞_Y,Y/X}`$, the sheafification of the presheaf on $`Y`$ given by:
$`U\{𝒞_Y`$-torsors on $`U`$ modulo pullbacks of $`𝒞_X`$-torsors$`\}.`$
(As above, $`𝒞_X`$ is the group scheme of regular centralizer subgroups with Lie algebra $`𝐜_X`$, and $`𝒞_Y:=f^{}(𝒞_X)`$. ) In fancier language, we could think of $`\mathrm{Tors}_{Y/X}`$ as a sheaf of Picard groupoids. But its objects have no automorphisms, so we are dealing in fact with a sheaf of abelian groups. In more detail:
We introduce the sheaf of Picard categories $`\mathrm{Tors}_{𝒞_Y,Y/X}`$ on $`\mathrm{Sch}_{et}(X)`$ as “$`𝒞_Y`$-torsors on $`Y`$ modulo pull-backs of $`𝒞_X`$-torsors”. The definition of $`\mathrm{Tors}_{𝒞_Y,Y/X}`$ is as follows:
First, consider the presheaf of categories $`\mathrm{Tors}_{𝒞_Y,Y/X}^{\mathrm{pre}}`$, whose objects over $`UX`$ are torsors over $`U\underset{𝑋}{\times }Y`$ with respect to the sheaf $`T_{\stackrel{~}{Y}}`$. Morphisms between two such torsors $`\tau ^{}`$ and $`\tau ^{\prime \prime }`$ are pairs $`(\tau _X,\sigma )`$, where $`\tau _X`$ is a $`T_{\stackrel{~}{X}}`$-torsor on $`U`$ and $`\sigma `$ is an isomorphism $`\tau ^{\prime \prime }\tau ^{}f^{}(\tau _X)`$. (Since $`f:YX`$ is dominant, and thus $`\mathrm{\Gamma }(U,T_{\stackrel{~}{X}})\mathrm{\Gamma }(U\underset{𝑋}{\times }Y,T_{\stackrel{~}{Y}})`$ is an injection, it is easy to see that the morphisms defined this way form a set and not just a category.)
The presheaf $`\mathrm{Tors}_{𝒞_Y,Y/X}^{\mathrm{pre}}`$ satisfies the first sheaf axiom, but not the second one, i.e. not every descent data is automatically effective. By applying the standard sheafification procedure, we obtain from $`\mathrm{Tors}_{𝒞_Y,Y/X}^{\mathrm{pre}}`$ a sheaf of Picard categories, which we denote by $`\mathrm{Tors}_{𝒞_Y,Y/X}`$.
Note, however, that since the morphism $`f:YX`$ is projective, objects of $`\mathrm{Tors}_{𝒞_Y,Y/X}`$ have no non-trivial automorphisms, because for every $`U`$ as above, the map $`\mathrm{\Gamma }(U,T_{\stackrel{~}{X}})\mathrm{\Gamma }(U\underset{𝑋}{\times }Y,T_{\stackrel{~}{Y}})`$ is in fact an isomorphism. Hence, $`\mathrm{Tors}_{Y/X}:=\mathrm{Tors}_{𝒞_Y,Y/X}`$ is in fact a sheaf of groups.
We need an explicit description of this sheaf:
###### Lemma 18.3.
There is a canonical identification:
$`\mathrm{Tors}_{Y/X}(X)=\{v\mathrm{Mor}_W(\stackrel{~}{X},\mathrm{Bun}_T(Y/X))|\alpha _iv_{|D_X^{\alpha _i}}=1\mathrm{Pic}(Y/X),`$ $`\alpha _iI\}.`$
(As always, $`I`$ denotes the set of simple roots $`\alpha _i.`$)
###### Proof.
We identify $`\mathrm{Tors}_{Y/X}`$ and $`\mathrm{Tors}_{Y/X}^{}`$ using Proposition 16.4. There is a natural map $`\iota :\mathrm{Tors}_{Y/X}^{}\mathrm{Mor}_W(\stackrel{~}{X},\mathrm{Bun}_T(Y/X))`$, sending a $`T`$-bundle on $`\stackrel{~}{Y}=\stackrel{~}{X}\underset{𝑋}{\times }Y`$ to its classifying morphism $`v`$. This map $`\iota `$ is clearly injective, and its image is contained in the RHS.
We still have to prove the surjectivity of $`\iota `$, i.e. to show that a morphism $`v`$ in the RHS satisfies the two compatibility conditions between $`\beta `$’s and $`\gamma `$’s stated in 16.3. It suffices to do so locally, and then we may assume that $`f:YX`$ has a section. In this case, we can identify $`\mathrm{Tors}_{Y/X}^{}`$ with the sheaf of $`T`$-bundles on $`\stackrel{~}{Y}`$ satisfying the two compatibility conditions between $`\beta `$’s and $`\gamma `$’s and which additionally are trivialized along the section $`XY`$. Similarly, we can identify $`\mathrm{Bun}_T(Y/X))`$ with $`T`$-bundles on $`Y`$ which are trivialized along the section.
Each of the compatibility conditions requires the equality of two given trivializations of some ($`T`$\- or $`𝔾_m`$-) bundle over $`D_X^i\underset{𝑋}{\times }Y`$. Now our assumption, $`\alpha _iv_{|D_X^i}=1`$, together with the assumed trivialization of all objects along the section, guarantees that these equalities hold over the section. The difference between the two trivializations is therefore a global automorphism which equals the identity along the section, so it is the identity everywhere since the fibers of $`f`$ are integral and proper.
#### 18.4.
By construction, we have a short exact sequence of Picard categories:
$$0\mathrm{Tors}_{T_{\stackrel{~}{X}}}f_{}(\mathrm{Tors}_{T_{\stackrel{~}{Y}}})\mathrm{Tors}_{Y/X}0.$$
As in Subsection 3.7, an element $`v\mathrm{Tors}_{Y/X}(X)`$ determines a $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}`$-gerbe which we denote $`𝒬_v`$. In fact, for $`(UX)\mathrm{Sch}_{et}(X)`$, $`𝒬_v(U)`$ is the category of all possible lifts of $`v`$ to a $`T_{\stackrel{~}{Y}}`$–torsor on $`U\underset{𝑋}{\times }Y`$.
The main result of this section is the following analogue of Theorem 17.5:
###### Theorem 18.5.
A regularized $`G`$-bundle on $`Y`$ is the same as a triple:
(a) A cameral cover $`\stackrel{~}{X}X`$,
(b) A $`W`$-equivariant map $`v:\stackrel{~}{X}\mathrm{Bun}_T(Y/X)`$ (of $`X`$-schemes), satisfying:
$`\alpha _iv_{|D^{\alpha _i}}=1\mathrm{Pic}(Y/X),`$ simple root $`\alpha _i`$, and
(c) An object of $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)\underset{\mathrm{Tors}_{T_{\stackrel{~}{X}}}}{}𝒬_v`$.
###### Proof.
Let us fix a cameral cover $`\stackrel{~}{X}X`$ and consider regularized $`G`$-bundles on $`Y`$ corresponding to this fixed $`\stackrel{~}{X}`$ as a sheaf of categories over $`X`$, denoted by $`\mathrm{Reg}_{\stackrel{~}{X}}(Y)`$.
By Theorem 4.4, $`\mathrm{Reg}_{\stackrel{~}{X}}(Y)`$ is a gerbe over the sheaf of Picard categories $`f_{}(\mathrm{Tors}_{T_{\stackrel{~}{Y}}})`$. This gerbe is induced from the $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}`$-gerbe $`\mathrm{Higgs}_{\stackrel{~}{X}}`$ by the homomorphism $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}f_{}(\mathrm{Tors}_{T_{\stackrel{~}{Y}}})`$, cf. Section 3.9.
Thus, according to Lemma 3.10, we have a functor $`\mathrm{Reg}_{\stackrel{~}{X}}(Y)\mathrm{Tors}_{Y/X}`$, and for a given object $`v\mathrm{Tors}_{Y/X}(X)`$ the category-fiber of the above functor is a $`\mathrm{Tors}_{T_{\stackrel{~}{X}}}`$-gerbe, which can be canonically identified with $`\mathrm{Higgs}_{\stackrel{~}{X}}(X)\underset{\mathrm{Tors}_{T_{\stackrel{~}{X}}}}{}𝒬_v`$.
Finally, according to Lemma 18.3, an object $`v\mathrm{Tors}_{Y/X}(X)`$ is equivalent to data (b) above.
#### 18.6.
Now let us assume that $`X`$ is projective as well. As our analogue of $`\mathrm{𝐇𝐢𝐠𝐠𝐬}(X,K)`$, we will consider the algebraic stack $`\mathrm{𝐑𝐞𝐠}(X,Y)`$ which associates to a scheme $`S`$ the category of regularized $`G`$-bundles on $`S\underset{𝑋}{\times }Y`$ (with respect to the projection $`S\underset{𝑋}{\times }YS`$).
We can now describe an analogue of the Hitchin map. Indeed, let $`𝐁(X,Y)`$ be the stack whose $`S`$-points are pairs: $`(\stackrel{~}{X}_S,v)`$ consisting of a cameral cover of $`S\times X`$ and a $`W`$-equivariant map $`v:\stackrel{~}{X}_S\mathrm{Bun}_T(Y/X)`$ of $`X`$-schemes.
We have a natural map of stacks $`h:\mathrm{𝐑𝐞𝐠}(X,Y)𝐁(X,Y)`$.
###### Corollary 18.7.
The fiber of the spectral map $`\mathrm{𝐑𝐞𝐠}(X,Y)𝐁(X,Y)`$ over a cameral point $`(\stackrel{~}{X},v)𝐁(X,Y)`$ can be identified with $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)\underset{\mathrm{Tors}_{T_{\stackrel{~}{X}}}}{}𝒬_v`$. The set of isomorphism classes of objects of this fiber is a torsor over the abelian group $`H^1(X,T_{\stackrel{~}{X}})`$.
In the case of $`K`$-valued Higgs bundles, we saw in Corollary 17.8 that the fiber of the Hitchin map $`\mathrm{𝐇𝐢𝐠𝐠𝐬}(X,K)𝐁(X,K)`$ is independent of the line bundle $`K`$. Note in contrast that the fiber $`\mathrm{Higgs}_{\stackrel{~}{X}}^{}(X)\underset{\mathrm{Tors}_{T_{\stackrel{~}{X}}}}{}𝒬_v`$ of the spectral map could depend on the original map $`f:YX`$. This dependence is mild though: it affects only the second factor, $`𝒬_v`$. A simplification occurs when $`f:YX`$ has a global section: in this case $`𝒬_v`$ is always trivial, because its defining short exact sequence of Picard categories 18.4 is split. It follows that the category $`\mathrm{Reg}_{\stackrel{~}{X}}(Y)`$ of regularized bundles with a specified cameral cover $`\stackrel{~}{X}`$ factors:
$$\mathrm{Reg}_{\stackrel{~}{X}}(Y)=\mathrm{Tors}_{Y/X}\times \mathrm{Higgs}_{\stackrel{~}{X}}^{}(X).$$
#### 18.8.
In addition to the stack $`𝐁(X,Y)`$, one can also define an analogue of the space $`\mathrm{𝐇𝐢𝐭𝐜𝐡}(X,K)`$: we let the space $`\mathrm{𝐇𝐢𝐭𝐜𝐡}(X,Y)`$ denote the scheme of all sections of the fibration $`(\mathrm{Bun}_T(Y/X))/WX`$. As before, we have an obvious retraction $`𝐁(X,Y)\mathrm{𝐇𝐢𝐭𝐜𝐡}(X,Y)`$. The analogue of the embedding $`\mathrm{𝐇𝐢𝐭𝐜𝐡}(X,K)𝐁(X,K)`$ can be described as follows:
Consider the $`W`$-cover $`\mathrm{Bun}_T(Y/X)(\mathrm{Bun}_T(Y/X))/W`$ and let $`(\mathrm{Bun}_T(Y/X))^0/W`$ be the maximal open subscheme over which this cover is cameral; let $`\mathrm{Bun}_T(Y/X)^0`$ denote its preimage in $`\mathrm{Bun}_T(Y/X)`$.
We will have to shrink $`(\mathrm{Bun}_T(Y/X))^0/W`$ to a still smaller open subscheme:
For a simple root $`\alpha _i`$ consider the corresponding ramification divisor
$$D_{(\mathrm{Bun}_T(Y/X))^0/W}^{\alpha _i}\mathrm{Bun}_T(Y/X)^0.$$
Under the map $`\mathrm{Bun}_T(Y/X)\mathrm{Pic}(Y/X)`$ given by $`\alpha _i`$, the image of $`D_{(\mathrm{Bun}_T(Y/X))^0/W}^{\alpha _i}`$ is contained in the set of $`_2`$-torsion points of $`\mathrm{Pic}(Y/X)`$.
We define the open subscheme $`(\mathrm{Bun}_T(Y/X))^{00}/W`$ of $`(\mathrm{Bun}_T(Y/X))^0/W`$ by removing those points, whose preimage in $`\mathrm{Bun}_T(Y/X)^0`$ maps to a non-unit point in $`\mathrm{Pic}(Y/X)`$ by means of the above map. Let $`\mathrm{Bun}_T(Y/X))^{00}(\mathrm{Bun}_T(Y/X))^{00}/W`$ denote the corresponding cameral cover.
Finally, let $`\mathrm{𝐇𝐢𝐭𝐜𝐡}(X,Y)^0`$ be the open subscheme of $`\mathrm{𝐇𝐢𝐭𝐜𝐡}(X,Y)`$, which corresponds to sections whose values belong to $`(\mathrm{Bun}_T(Y/X))^{00}/W`$.
The fiber product construction gives the desired map $`i:\mathrm{𝐇𝐢𝐭𝐜𝐡}(X,Y)^0𝐁(X,Y)`$. Its image is the open substack corresponding to the locus where the map $`v:\stackrel{~}{X}\mathrm{Bun}_T(Y/X)`$ is an embedding.
#### 18.9. The case of an elliptic fiber
The main relevance of the above results is to the case that $`f:YX`$ is an elliptic fibration. This is due to the existence in this case of a good notion of a regular bundle, analogous to the notion of a regular $`K`$-valued Higgs bundle. Take the group $`G`$ to be semisimple, and consider the case of a single elliptic curve $`Z`$.
For any semistable $`G`$-bundle $`E_G`$ on $`Z`$, the dimension of the group $`H:=\mathrm{Aut}_G(E_G)`$ of (global) automorphisms of $`E_G`$ is $`r`$. We say that $`E_G`$ is regular if $`dim(H)=r`$. In this case, $`H`$ is commutative and there exists an embedding $`HG`$ and a principal $`H`$-bundle $`E_H`$ on $`Z`$ such that $`E_GG\stackrel{𝐻}{\times }E_H`$. A regular bundle has a unique regularization.
These results can be found in and elsewhere. In fact, the moduli space $`M_G(Z)`$ of (S-equivalence classes of) semistable, topologically trivial $`G`$-bundles on the elliptic curve $`Z`$ is well understood. As a complex variety, it is isomorphic to $`M_T(Z)/W`$. (This is proved analytically (e.g. ) using Borel’s result that in a simply connected compact group, any two commuting elements are contained in a maximal torus. An algebraic proof was given in .) Each S-equivalence class contains a unique regular representative as well as a unique semisimple representative (i.e. one whose structure group can be reduced to $`T`$). For a generic point of the moduli space, the S-equivalence class consists of a unique isomorphism class, which is both regular and semisimple. A similar but somewhat more complicated description exists for all reductive $`G`$, cf. . Returning to an elliptic family $`f:YX`$, we find ourselves in a situation analogous to that which we had for $`K`$-valued Higgs bundles: a “generic” $`G`$-bundle on $`Y`$ which is semistable along the elliptic fibers should be regular on the generic fiber, and therefore its restriction to a dense open $`X_0X`$ should admit a unique regularization to which we can apply our results.
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# Discrete Flavor Symmetries and Mass Matrix Textures
## I Introduction
The mechanism that generates the fermion masses is not yet understood. In the standard model (SM) the masses and mixings are simply parameters that can be adjusted to agree with experiment. One hope is that the Yukawa couplings in the SM can be understood more fully when the theory is embedded in a more fundamental theory, and relationships between masses and mixings might then be established. Symmetries based on embedding the gauge symmetries of the SM in larger gauge groups (unified theories) have been used for a long time and some reasonable mass patterns can be derived which are consistent with experiment. These symmetries have come to be called vertical symmetries to distinguish them from the horizontal symmetries (or flavor symmetries) that relate fermions from the different generations. In this paper we show how discrete Abelian horizontal symmetries based on $`Z_m`$ can account for some of the successful texture patterns.
There are a number of reasons one might want to extend a $`U(1)`$ flavor symmetry so that it contains an additional discrete $`Z_m`$ component. (A) The additional discrete symmetry offers a solution to the seemingly inconsistent large mixing observed in the atmospheric neutrino data and a hierarchy in the muon and tau neutrino masses. In models with supersymmetric Abelian flavor symmetries, large $`\nu _\mu \nu _\tau `$ mixing is achieved via a light neutrino mass matrix of the form
$`\left(\begin{array}{cc}C& B\\ B& A\end{array}\right){\displaystyle \frac{v^2}{M}},`$ (1)
$`A,B,C𝒪(1),`$ (2)
where $`v`$ is some electroweak scale vacuum expectation value and $`M`$ is the Majorana neutrino mass scale. The eigenvalues are typically the same order of magnitude. It requires a fine-tuning of the order one parameters $`A`$, $`B`$ and $`C`$ to achieve large mixing between neutrinos and widely separated neutrino masses. Grossman, Nir and Shadmi advocated using a discrete symmetry to maintain the large mixing angle while achieving very different neutrino masses without fine-tuning. Discrete symmetries had been discussed previously, but the more recent experimental results indicating large mixing in the atmospheric (and perhaps solar) neutrino oscillations made this technique especially interesting. The use of a discrete flavor symmetry to understand the mass hierarchies and mixing angles for all Standard Model fermions was pursued in Ref. . Other authors have employed non-Abelian discrete flavor symmetries that have both one and two dimensional representations. This approach is particularly well suited for addressing the supersymmetric flavor problem where the first and second generation of superpartners should have very similar flavor properties and thus belong in the same representation of the flavor group. In this paper we restrict our attention to Abelian discrete symmetries. (B) The phenomenological predictions for quark mass ratios and Cabibbo-Kobayashi-Maskawa (CKM) matrix elements can be retained, but the contributions arise from a smaller number of parameters. (C) The process of creating the baryon asymmetry of the Universe by having CP-violating asymmetric decays of heavy neutrinos can be greatly enhanced as a natural consequence of solving issue (A) above. (D) One can potentially solve the supersymmetric flavor problem by suppressing certain entries in the Yukawa matrices. The mechanism works by aligning the quark mass matrices with the squark mass-squared matrices, and one does not need to require that the first and second generation squarks are degenerate. The mixing matrix for the squark-quark-gluino couplings can be made close to a unit matrix, and the undesirable flavor-changing neutral currents (FCNCs) are suppressed. The required suppressions are not possible with a $`U(1)`$ symmetry.
During the last few years, there has been great interest in using new continuous Abelian as well as discrete Abelian and non-Abelian symmetries in the minimal supersymmetric standard model (MSSM) to describe the experimental (phenomenological) data on the fermion masses and mixings,-. Superstring theories appear to have $`U(1)`$ symmetries and symmetries involving its discrete subgroups as a generic feature. If the $`U(1)`$ flavor symmetry is gauged then a general assignment of flavor charges to the fields will be anomalous. One can imagine the anomaly is canceled via the Green-Schwarz mechanism, and one must check whether the correct relations are satisfied. A convenient way to ensure that the flavor charges are amenable to cancellation is to have the flavor symmetry commute with the $`SU(5)`$ grand unified theory. We present in this paper a model with a $`U(1)\times Z_2`$ flavor symmetry with at least four texture zeros (in the up and down quark Yukawa matrices) that commutes with the $`SU(5)`$ gauge group.
The paper is organized as follows. In Section II we briefly discuss flavor symmetries and how they can in principle account for the experimentally observed hierarchies in the quark masses and mixing angles. In Section III we discuss the possible role of a discrete component in the flavor symmetry. Section IV then lists the phenomenological requirements that must be met in the quark sector of the standard model. In Section V a particular model for which the flavor symmetry commutes with an $`SU(5)`$ grand unified symmetry is presented, and the phenomenology is extended to include the leptons. The consequences for neutrino oscillations and the charged lepton masses are discussed. Finally we present our conclusions in Section VI.
## II Flavor Symmetries
The hierarchical structure of the fermion mass matrices hints that there may be a spontaneously broken family symmetry responsible for the suppression of Yukawa couplings. In this paper we employ supersymmetric Abelian horizontal symmetries. These flavor symmetries allow the fermion mass and mixing hierarchies to be naturally generated from nonrenormalizable terms in the effective low-energy theory.
The idea is quite simple and easily implemented. There is some field $`S`$ which is charged under a $`U(1)`$ family symmetry, and without loss of generality, we can assume that its charge is -1. There are terms contributing to effective Yukawa couplings for the quarks,
$`Q_i\overline{d}_jH_d\left({\displaystyle \frac{S}{\mathrm{\Lambda }_L}}\right)^{m_{ij}}+Q_i\overline{u}_jH_u\left({\displaystyle \frac{S}{\mathrm{\Lambda }_L}}\right)^{n_{ij}},`$ (3)
and the integer exponents $`m_{ij}`$ and $`n_{ij}`$ are easily calculated in terms of the horizontal symmetry charges of the quark and Higgs fields. The scale, $`\mathrm{\Lambda }_L`$, where massive states are integrated out of the fundamental theory to produce an effective theory, is assumed to be larger than the vev, $`<S>`$ of the singlet scalar field so the parameter $`<S>/\mathrm{\Lambda }_L`$ is a small one. We henceforth require the Higgs fields to be uncharged under the $`U(1)`$ family symmetry, then the exponent $`m_{ij}`$ is just the sum of the horizontal charge of the fields $`Q_i`$ and $`\overline{d}_j`$. The hierarchy is generated from terms in the superpotential that carry integer charges $`m_{ij},n_{ij}0`$. If we call the small breaking parameter $`<S>/\mathrm{\Lambda }_L\lambda `$, then the generated terms for say the down quark Yukawa matrix will be of order $`\lambda ^{m_{ij}}`$. We will restrict our attention in this paper to flavor charges for the Standard Model fields that are non-negative. Here texture zeros refer to Yukawa matrix elements that can be replaced by an exact zero without affecting the leading order (in the small parameter $`\lambda `$) results for the mass eigenvalues and mixing angles. An analysis of the possible approaches to explaining the neutrino masses and mixings using $`U(1)`$ symmetries only is given in Ref. .
In models whose flavor symmetry contains two distinct components ($`U(1)\times U(1)`$, $`U(1)\times Z_m`$, etc.) one introduces two singlet scalars, $`S_1`$ and $`S_2`$, with horizontal charges
$`S_1(1,0),S_2(0,1),`$ (4)
which in general can have different vacuum expectation values $`<S_1>`$ and $`<S_2>`$. These can be related to a common expansion parameter $`\lambda `$ by setting
$`{\displaystyle \frac{<S_1>}{\mathrm{\Lambda }_L}}\lambda ^\beta ,{\displaystyle \frac{<S_2>}{\mathrm{\Lambda }_L}}\lambda ^\alpha .`$ (5)
In the following we identify $`\lambda `$ as the Cabibbo angle, and take $`\beta =1`$. In general, one can take $`\alpha 1`$, but for our explicit models we assume $`\alpha =1`$. The contributions to the Yukawa matrices arise from flavor invariant terms in
$`Q_i\overline{d}_jH_d\left({\displaystyle \frac{S_1}{\mathrm{\Lambda }_L}}\right)^{m_{ij}}\left({\displaystyle \frac{S_2}{\mathrm{\Lambda }_L}}\right)^{p_{ij}}+Q_i\overline{u}_jH_u\left({\displaystyle \frac{S_1}{\mathrm{\Lambda }_L}}\right)^{n_{ij}}\left({\displaystyle \frac{S_2}{\mathrm{\Lambda }_L}}\right)^{q_{ij}}.`$ (6)
It should be understood that there are undetermined order one coefficients multiplying these terms, and we assume in this paper that these coefficients are sufficiently close to one so as not to influence the hierarchy, i.e. somewhat greater than $`\lambda `$ and somewhat less than $`1/\lambda `$. Formulae for the Yukawa matrices for the quarks and charged leptons as well as mass matrices for the neutrinos that follow from the Froggatt-Nielsen mechanism are given in the Appendix.
If the flavor symmetry is $`U(1)`$ then there is a charge assignment in the quark sector that satisfies all the phenomenological requirements detailed in Section IV below. This solution was obtained by many authors. The up and down quark Yukawa matrices are
$`𝐔\left(\begin{array}{ccc}\lambda ^8& \lambda ^5& \lambda ^3\\ \lambda ^7& \lambda ^4& \lambda ^2\\ \lambda ^5& \lambda ^2& 1\end{array}\right),𝐃\left(\begin{array}{ccc}\lambda ^4& \lambda ^3& \lambda ^3\\ \lambda ^3& \lambda ^2& \lambda ^2\\ \lambda & 1& 1\end{array}\right).`$ (7)
The model we present in this paper will give the same phenomenological predictions as Eq. (7), but the discrete symmetry will suppress certain entries in comparison to the $`U(1)`$ flavor symmetry pattern shown in Eq. (7). After including a discrete component to the flavor symmetry, a different $`SU(5)`$ grand unified model can be constructed (see Section V).
## III Discrete Abelian Groups
In this section we discuss the two possible texture patterns for a $`2\times 2`$ matrix, and then show how to put these $`2\times 2`$ blocks together to form the realistic case of texture patterns for three generations.
### A Suppressing in $`2\times 2`$ blocks
When the flavor symmetry is $`U(1)`$, there is a sum rule among the exponents in any $`2\times 2`$ block. For example, the up quark Yukawa matrix necessarily has the relationship,
$`n_{ii}+n_{jj}n_{ij}n_{ji}=0`$ (8)
between the exponents, $`n_{ij}q_i+u_j`$. The Yukawa matrices in Eq. (7) obey this rule, for example. However, these relationships between elements of the Yukawa matrices can be avoided if the flavor symmetry has a $`Z_m`$ component. We can illustrate this with a simple example with a $`Z_2`$ symmetry: Consider two generations with $`Z_2`$ flavor charges as follows,
$`Q_L:q_i^Z=(0,1),\overline{u}_R:u_i^Z=(0,1),i=1,2,`$ (9)
where the first number for each field gives the charge for the first generation and the second number gives the charge for the second generation. Then performing the $`Z_2`$ arithmetic in constructing the contribution to the Yukawa matrices yields, in general,
$`\left(\begin{array}{cc}\lambda ^{[q_1^Z+u_1^Z]}& \lambda ^{[q_1^Z+u_2^Z]}\\ \lambda ^{[q_2^Z+u_1^Z]}& \lambda ^{[q_2^Z+u_2^Z]}\end{array}\right).`$ (10)
We use brackets around the exponents, $`[]`$, to denote that we are modding out by two according to the $`Z_2`$ addition rules. In the case of the particular choice of charges in Eq. (9) and taking $`<S_2>/\mathrm{\Lambda }_L\lambda `$
$`\left(\begin{array}{cc}\lambda ^0& \lambda ^1\\ \lambda ^1& \lambda ^0\end{array}\right),`$ (11)
So this set of charges yields a Yukawa matrix that does not satisfy the rule in Eq. (8). If one adds in nontrivial contributions from the $`U(1)`$ part of the flavor symmetry, one sees that the off-diagonal entries in Eq. (11) are suppressed relative to the expectation from Eq. (8). For example, assume that the fields have (in addition to the $`Z_2`$ assignments in Eq. (11)) the $`U(1)`$ charge assignments
$`Q_L:q_i=(3,0),\overline{u}_R:u_i=(1,0),i=1,2,`$ (12)
which, in general, give a contribution to the Yukawa matrices
$`\left(\begin{array}{cc}\lambda ^{q_1+u_1}& \lambda ^{q_1+u_2}\\ \lambda ^{q_2+u_1}& \lambda ^{q_2+u_2}\end{array}\right).`$ (13)
The particular choice of charges in Eq. (12) together with taking $`<S_1>/\mathrm{\Lambda }_L\lambda `$ yields the contribution to the Yukawa matrices from the $`U(1)`$ charges of the form
$`\left(\begin{array}{cc}\lambda ^4& \lambda ^3\\ \lambda ^1& \lambda ^0\end{array}\right).`$ (14)
It should be clear from Eq. (6) that the overall contribution to the Yukawa matrix is the product of the contribution from the $`Z_2`$ charges in Eq. (11) and the contribution from the $`U(1)`$ charges in Eq. (14) for each element of the Yukawa matrix. For the example here we get the following result for the $`U(1)\times Z_2`$ flavor symmetry,
$`𝐔\left(\begin{array}{cc}\lambda ^4& \lambda ^4\\ \lambda ^2& \lambda ^0\end{array}\right).`$ (15)
So one sees that the off-diagonal entries are suppressed relative to the expectation in Eq. (14), and it is not difficult to convince oneself that this suppression comes entirely from the $`Z_2`$ part of the flavor symmetry.
The relevance of the above example to the present paper is the following. We are interested in determining the phenomenological predictions of the Yukawa matrices and then comparing them to the experimental data. This requires that we diagonalize the Yukawa matrices to determine the eigenvalues (masses) and the mixing angles (CKM elements) as explained in the Appendix. In the example we arrived at the Yukawa matrix in Eq. (15), for which it is immediately clear that the eigenvalues are of order $`\lambda ^0`$ and $`\lambda ^4`$, while the mixing angles in $`V_u^L`$ and $`V_u^R`$ (see Eq. (64)) are of order $`\lambda ^4`$ and $`\lambda ^2`$, respectively. So one notes that the left-handed mixing angle is suppressed by $`\lambda ^2`$ in comparison to the expectation from a $`U(1)`$ symmetry alone that gives the same mass eigenvalues, namely
$`\left(\begin{array}{cc}\lambda ^4& \lambda ^2\\ \lambda ^2& \lambda ^0\end{array}\right).`$ (16)
One can compare this simple $`2\times 2`$ example with the second and third generations of Eq. (7). We denote this suppression in the following way,
$`\left(\begin{array}{cc}X& 0\\ 0& X\end{array}\right).`$ (17)
We say that the Yukawa matrix has texture zeros in the off-diagonal positions. A texture zero defined in this way is not a true zero, but is negligible to the leading order in the small parameter $`\lambda `$ as far as the mass eigenvalues and the left-handed mixing (diagonalization) angles are concerned. The right-handed mixing angle is not suppressed but this affects neither the CKM mixing angles nor the mass eigenvalues. The physical observables are the elements of the CKM matrix, Eq. (65), and involve contributions from the diagonalization of the down-Yukawa matrix as well. So if there is a contribution from the diagonalization of down-Yukawa matrix of order $`\lambda `$, the contribution from the up-Yukawa matrix will be negligible in comparison. In this case we promote the texture zero to a true phenomenological zero: the contribution from the off-diagonal elements of the up-quark Yukawa matrix does not contribute to the determination of any physical observable (quark mass or CKM element) to leading order in the expansion in the small parameter $`\lambda `$.
We can also engineer a suppression along the diagonal elements of a Yukawa matrix. This case is somewhat trickier than the previous case, so we proceed now to present another example in the case of just two generations: For the $`Z_2`$ charges, consider the assignment
$`Q_L:q_i^Z=(1,0),\overline{d}_R:d_i^Z=(0,1),i=1,2.`$ (18)
and, for the $`U(1)`$ charges, make the assignment
$`Q_L:q_i=(3,0),\overline{d}_R:d_i=(1,0),i=1,2.`$ (19)
Then one obtains
$`\left(\begin{array}{cc}\lambda ^1& \lambda ^0\\ \lambda ^0& \lambda ^1\end{array}\right).`$ (20)
for the contribution from the $`Z_2`$ charges and
$`\left(\begin{array}{cc}\lambda ^4& \lambda ^3\\ \lambda ^1& \lambda ^0\end{array}\right),`$ (21)
from the $`U(1)`$ sector. The contributions from the full $`U(1)\times Z_2`$ symmetry give the Yukawa matrix,
$`𝐃\left(\begin{array}{cc}\lambda ^5& \lambda ^3\\ \lambda ^1& \lambda ^1\end{array}\right).`$ (22)
The eigenvalues of this matrix are $`\lambda ^0`$ and $`\lambda ^2`$ and the mixing angle for the left-handed diagonalization matrix is of order $`\lambda ^2`$. So one can interpret the entry of order $`\lambda ^5`$ as being phenomenologically irrelevant to leading order in powers of $`\lambda `$ and the texture zero pattern is
$`\left(\begin{array}{cc}0& X\\ X& X\end{array}\right).`$ (23)
Note that the entry is not phenomenologically irrelevant, and is still denoted by $`X`$.
Having generated the Yukawa matrices $`𝐔`$ and $`𝐃`$ in Eqs. (15) and (22), we can account for phenomenological requirements (the full set of experimental data for fermion masses and mixings is given in the next section),
$`{\displaystyle \frac{m_c}{m_t}}\lambda ^4,{\displaystyle \frac{m_s}{m_b}}\lambda ^2,|V_{cb}|\lambda ^2.`$ (24)
The mixing angles for $`𝐔`$ are $`\mathrm{sin}\theta _L^u\lambda ^4`$ and $`\mathrm{sin}\theta _R^u\lambda ^2`$, while for $`𝐃`$ they are $`\mathrm{sin}\theta _L^d\lambda ^2`$ and $`\mathrm{sin}\theta _R^d\lambda ^0`$. The leading order contribution to $`|V_{cb}|`$ according to Eq. (65) is then given entirely by $`\mathrm{sin}\theta _L^d\lambda ^2`$ since $`\mathrm{sin}\theta _L^u`$ is suppressed by a relative factor of $`\lambda ^2`$. The mass eigenvalue ratios in Eq. (24) are properly accounted for. So phenomenologically viable Yukawa matrices can be found with texture zeros, and these zeros reduce the number of unknown order one coefficients that contribute to masses and mixing angles at the leading order in $`\lambda `$.
An important feature of the $`U(1)\times Z_2`$ flavor symmetry is that one can achieve different leading order contributions to the left-handed mixing angles in $`V_u^L`$ and $`V_d^L`$, as shown in the above example. In a model with a $`U(1)`$ symmetry, these mixing angles are determined entirely by the charges $`Q_L`$ (and not $`\overline{u}_R`$ and $`\overline{d}_R`$). So the presence of the $`Z_2`$ symmetry allows one to suppress the contribution to the CKM mixings from either the $`𝐔`$ or the $`𝐃`$ matrix.
It is not difficult to generalize the discussion to an arbitrary $`Z_m`$. The exponents of $`\lambda `$ are given by Eq. (10), and the conditions satisfied by the $`Z_m`$ charges that lead to texture suppressions are
$`[q_i^Z+u_i^Z]+[q_j^Z+u_j^Z][q_i^Z+u_j^Z][q_j^Z+u_i^Z]=\pm m,m2,`$ (25)
where the case $`+m`$ results in a suppression of the diagonal entries of the $`2\times 2`$ matrix, and the case $`m`$ results in a suppression of the off-diagonal entries. We remind the reader that the square brackets in Eq. (25) indicate a modding by the integer $`m`$.
### B Extending $`Z_m`$-induced suppressions to $`3\times 3`$ matrices
In the last subsection examples of how to suppress entries in a $`2\times 2`$ Yukawa matrix were presented. One can extend this result to the three generation case by considering $`2\times 2`$ blocks. One has three such blocks in the case of three generations: namely the \[2-3\], \[1-3\] and the \[1-2\] blocks. One can build up a texture pattern for a $`3\times 3`$ matrix by placing zeros in the desired positions of these $`2\times 2`$ blocks. As demonstrated in the last subsection, in each $`2\times 2`$ block, one can have either a texture zero in the off-diagonal or in a diagonal position, but not both at the same time. As an example consider the matrix:
$$\left(\begin{array}{ccc}0& 0& X\\ 0& X& X\\ X& X& X\end{array}\right)$$
(26)
All the zeros cannot be obtained by a assigning charges in the \[1-2\] block alone, since this would require the zeros to be in both the diagonal and off-diagonal positions. However this texture pattern can be obtained by assigning the zero on the diagonal to the \[1-3\] block, while off-diagonal zeros can be assigned to the \[1-2\] block. As in the case of only two generations, one can obtain the texture pattern by considering only the $`Z_m`$ component of the flavor symmetry. One can obtain the required texture in Eq. (26) when $`m=3`$ by the following assignment of $`Z_3`$ charges
$`Q_L:q_i^Z=(2,0,1),\overline{u}_R:u_i^Z=(2,0,1).i=1,2,3,`$ (27)
The contribution to the $`3\times 3`$ matrix Yukawa matrix from this $`Z_3`$ charge assignment is
$`\left(\begin{array}{ccc}\lambda ^1& \lambda ^2& \lambda ^0\\ \lambda ^2& \lambda ^0& \lambda ^1\\ \lambda ^0& \lambda ^1& \lambda ^2\end{array}\right),`$ (28)
The $`U(1)`$ contributions have not been included yet in Eq. (28). While the \[2-3\] block does not have a suppressing pattern, the suppression in the \[1-3\] block suppresses the diagonal element . Finally, the \[1-2\] block suppresses the off-diagonal and elements.
Continuing with our example, if we assign the $`U(1)`$ flavor charges
$`Q_L:q_i=(6,3,0),\overline{u}_R:u_i=(6,3,0).i=1,2,3,`$ (29)
to the quark fields, we obtain the following contribution to the up-type Yukawa matrix,
$`\left(\begin{array}{ccc}\lambda ^{12}& \lambda ^9& \lambda ^6\\ \lambda ^9& \lambda ^6& \lambda ^3\\ \lambda ^6& \lambda ^3& 1\end{array}\right).`$ (30)
Putting the contributions from both components of the $`U(1)\times Z_3`$ flavor symmetry together gives the following up-type Yukawa matrix (after dropping an overall factor of $`\lambda ^2`$ which is irrelevant as far as the hierarchy is concerned),
$`𝐔`$ $`\left(\begin{array}{ccc}\lambda ^{11}& \lambda ^9& \lambda ^4\\ \lambda ^9& \lambda ^4& \lambda ^2\\ \lambda ^4& \lambda ^2& 1\end{array}\right).`$ (31)
One can always diagonalize matrices arises from Abelian flavor symmetries of the type described here in stages, by diagonalizing the \[2-3\] block, followed by the \[1-3\] block, and finally diagonalizing the \[1-2\] block. The diagonalization in the \[2-3\] block does not produce any texture zero because $`(\lambda ^2)(\lambda ^2)/(1)(\lambda ^4)`$ as in Eq. (8). Each order one coefficient in the \[2-3\] block plays a role in determining the leading order diagonalization of that block. However in the diagonalization of the \[1-3\] block, one notices that element ($`\lambda ^{11}`$) is suppressed by a factor of order $`\lambda ^3`$ compared to the product of the and elements. So to leading order in an expansion in $`\lambda `$, the diagonalization of the matrix in Eq. (31) is the same as a matrix where the $`\lambda ^{11}`$ element is replaced with zero (and we call such an entry a texture zero). So we have the following matrix whose diagonalization is equivalent to leading order to the original matrix
$`𝐔`$ $`\left(\begin{array}{ccc}0& \lambda ^9& \lambda ^4\\ \lambda ^9& \lambda ^4& \lambda ^2\\ \lambda ^4& \lambda ^2& 1\end{array}\right).`$ (32)
Finally we must determine if any of the elements in the \[1-2\] block are suppressed. Suppose the diagonalization has been performed in the \[2-3\] and \[1-3\] blocks. Then the matrix has the form
$`\left(\begin{array}{ccc}\lambda ^8& \lambda ^9& 0\\ \lambda ^9& \lambda ^4& 0\\ 0& 0& 1\end{array}\right).`$ (33)
The $`\lambda ^8`$ entry would be generated in the element (and the $`\lambda ^{11}`$ element can be neglected in comparison, as described above). So a subsequent diagonalization of the \[1-2\] block indicates that texture zeros occur in the off-diagonal elements as
$`𝐔`$ $`\left(\begin{array}{ccc}0& 0& \lambda ^4\\ 0& \lambda ^4& \lambda ^2\\ \lambda ^4& \lambda ^2& 1\end{array}\right).`$ (34)
In other words, to leading order in $`\lambda `$ the diagonalization of the first matrix in Eq. (31) is the same as the diagonalization of the matrix in Eq. (34).
By proceeding in this way, one can systematically construct all possible matrices with texture zeros in the desired positions. The task then is to combine a texture pattern for the up-type Yukawa matrix with another texture pattern for the down-type Yukawa matrix, and check whether all the phenomenological requirements can be satisfied. We now turn to the experimental data for the quark and lepton masses and mixing angles.
## IV Phenomenological Requirements in the Quark Sector
If one must satisfy the phenomenological constraints with positive flavor charges, then the Eq. (7) is the solution that results from a $`U(1)`$ flavor symmetry. Using a $`U(1)\times Z_m`$ flavor symmetry instead will change the exponents by adding $`m`$ in certain elements. The relevant equations for the CKM matrix elements that are valid for this category of matrices are,
$`|V_{us}|`$ $`=`$ $`\left({\displaystyle \frac{d_{12}}{\stackrel{~}{d}_{22}}}{\displaystyle \frac{d_{13}d_{32}}{\stackrel{~}{d}_{22}}}\right)\left({\displaystyle \frac{u_{12}}{\stackrel{~}{u}_{22}}}{\displaystyle \frac{u_{13}u_{32}}{\stackrel{~}{u}_{22}}}\right),`$ (35)
$`|V_{cb}|`$ $`=`$ $`d_{23}+d_{22}d_{32}^{}u_{23},`$ (36)
$`|V_{ub}|`$ $`=`$ $`(d_{13}+d_{12}d_{32}^{}u_{13})\left({\displaystyle \frac{u_{12}}{\stackrel{~}{u}_{22}}}{\displaystyle \frac{u_{13}u_{32}}{\stackrel{~}{u}_{22}}}\right)(d_{23}+d_{22}d_{32}^{}u_{23}),`$ (37)
$`|V_{td}|`$ $`=`$ $`(d_{13}+d_{12}d_{32}^{}u_{13})+\left({\displaystyle \frac{d_{12}}{\stackrel{~}{d}_{22}}}{\displaystyle \frac{d_{13}d_{32}}{\stackrel{~}{d}_{22}}}\right)(d_{23}+d_{22}d_{32}^{}u_{23}),`$ (38)
where $`d_{ij}=𝐃_{ij}/𝐃_{33}`$ and $`\stackrel{~}{d}_{22}=d_{22}d_{23}d_{32}`$ and $`\stackrel{~}{u}_{22}=u_{22}u_{23}u_{32}`$. It is understood that there will in general be relative phases between the terms on the right hand sides of Eqs. (35)-(38), which are the correct forms to evaluate the leading orders for Yukawa matrices of the form considered in this paper.
Taking the expansion parameter to be the Cabibbo angle, $`\lambda =|V_{us}|`$, then the experimental constraints
$`|V_{us}|=0.2196\pm 0.0023,|V_{cb}|=0.0395\pm 0.0017,\left|{\displaystyle \frac{V_{ub}}{V_{cb}}}\right|=0.08\pm 0.02,`$ (39)
on the CKM matrix can be identified in terms of powers of $`\lambda `$ by the followingThere are renormalization scaling factors that relate the experimental data at the electroweak scale, Eq. (39), to the relationships at the high scale.,
$`|V_{us}|\lambda ,|V_{cb}|\lambda ^2,|V_{ub}|\lambda ^3\lambda ^4,\left|{\displaystyle \frac{V_{ub}}{V_{cb}}}\right|\lambda \lambda ^2.`$ (40)
We consider a model of Yukawa matrices to describe the experimental data satisfactorily if the leading order contribution to the CKM elements agrees with Eq. (40). For $`|V_{ub}|`$ and $`|V_{ub}/V_{cb}|`$ we accept two values for the exponent of the leading contribution. The constraint on $`|V_{ub}/V_{cb}|`$ can be expressed in a stronger way at 90% confidence level as $`0.25\lambda 0.5\lambda `$. One also has a constraint on the CKM elements from $`B_d^0\overline{B}_d^0`$ mixing,
$`|V_{tb}^{}V_{td}|=0.0084\pm 0.0018,`$ (41)
which implies that
$`|V_{td}|\lambda ^3.`$ (42)
The eigenvalues of the Yukawa matrices are constrained by the following requirements from experimental observations
$`{\displaystyle \frac{m_c}{m_t}}\lambda ^4,{\displaystyle \frac{m_u}{m_c}}\lambda ^4,{\displaystyle \frac{m_s}{m_b}}\lambda ^2,{\displaystyle \frac{m_d}{m_s}}\lambda ^2.`$ (43)
These phenomenological requirements will be used in the next section to constrain the Yukawa matrix patterns that can successfully reproduce the experimental data.
## V Grand Unified Model
In this section we derive an assignment of charges in $`U(1)\times Z_2`$ that has the maximum number of texture suppressions (four) that is consistent with a $`SU(5)`$ grand unified symmetry. Since the flavor symmetry is required to commute with $`SU(5)`$, this means that there must be a common flavor charge assignment for all particles in each multiplet of $`SU(5)`$. We restrict our attention to the case of a $`Z_2`$ symmetry, since (as described earlier) it is the only possible $`Z_m`$ symmetry that can reproduce a hierarchy in neutrino masses of order $`\lambda ^2`$.
Firstly, we have found that all the solutions from the $`U(1)\times Z_2`$ flavor symmetry that satisfy the quark sector phenomenology have the following property: the entry and the \[3-3\] entry of the down quark Yukawa matrix, $`𝐃`$, are the same order of magnitude. If the flavor symmetry is embedded in a grand unified model, the charged lepton Yukawa matrix will be given by the transpose of $`𝐃`$. Then the feature of Eq. (7), that the right-handed mixing matrix that diagonalizes the down-quark Yukawa matrix, D, is of order one in the \[2-3\] block, is retained. This has the important consequence that, if the lepton charges are related to the down quark charges by a grand unified theory, then the charged lepton matrix will require a large mixing between in the \[2-3\] block to diagonalize it. This results in a large mixing between the second and third generation of neutrinos, and can naturally explain how the atmospheric neutrino mixing can be large (order one) while the quark mixing between the second and third generations, $`|V_{cb}|`$, can be small (order $`\lambda ^2`$). This has been called the “lopsided” solution to the producing the required atmospheric neutrino mixing in grand unified models. This occurs in all the models necessarily after applying the phenomenological requirements $`|V_{cb}|\lambda ^2`$ and $`m_s/m_b\lambda ^2`$.
In models in which the $`U(1)`$ flavor symmetry is gauged and anomalous, one can imagine the anomaly is canceled via the Green-Schwarz mechanism. A convenient way to ensure that the flavor charges are amenable to cancellation is to have the flavor symmetry commute with the $`SU(5)`$ grand unified theory<sup>§</sup><sup>§</sup>§The mixed Standard Model-$`U(1)`$ anomalies can be canceled entirely by the Green-Schwarz mechanism if the $`U(1)`$ charges $`X`$ satisfy the relations $`tr(XT_aT_b)tr(T_aT_b)`$ and $`tr(X^2Y)=0`$ where $`T_a`$ are the Standard Model generators. These relations are satisfied automatically if the $`U(1)`$ charges respect the $`SU(5)`$ symmetry.. In the traditional $`SU(5)`$ grand unified theory, the fields $`Q_L`$ and $`\overline{u}_R`$ are assigned to the $`\mathrm{𝟏𝟎}`$ representation, and the $`\overline{d}_R`$ is assigned to the $`\mathrm{𝟓}^{}`$ representation. We have found a texture pattern for the up and down quark Yukawa matrices with four texture zeros for which the flavor symmetry quantum number assignment commutes with an $`SU(5)`$ grand unified gauge group. This texture pattern yields
$`𝐔\left(\begin{array}{ccc}\lambda ^8& 0& 0\\ 0& \lambda ^4& \lambda ^2\\ 0& \lambda ^2& 1\end{array}\right)𝐃\left(\begin{array}{ccc}0& 0& \lambda ^3\\ 0& \lambda ^2& \lambda ^2\\ \lambda & 1& 1\end{array}\right)`$ $`\begin{array}{cccc}Q_L:& (4,1)& (2,0)& (0,0)\\ \overline{u}_R:& (4,1)& (2,0)& (0,0)\\ \overline{d}_R:& (2,0)& (1,0)& (0,1)\end{array}.`$ (44)
This assignment has common $`U(1)\times Z_2`$ flavor symmetry quantum numbers for the $`Q_L`$ and $`\overline{u}_R`$ fields in the $`\mathrm{𝟏𝟎}`$, and a systematic search reveals that no other texture pattern with four or more texture zeros satisfies this property. Finding an assignment for which the flavor symmetry commutes with SU(5) allows us to assign flavor charges to the rest of the $`SU(5)`$ multiplets, namely the charged leptons and neutrinos.
The texture pattern given by Eq. (44) has the following feature: The CKM mixing $`|V_{cb}|`$ arises from contributions of order $`\lambda ^2`$ from the diagonalizations of both the $`𝐔`$ and $`𝐃`$ Yukawa matrices. All other CKM mixing angles ($`|V_{us}|`$, $`|V_{ub}|`$ and $`|V_{td}|`$) arise solely from the $`𝐃`$ Yukawa matrix.
Given the quantum number assignment in Eq. (44), we can extend the model to encompass the leptons. The field $`\overline{e}_R`$ fills out the $`\mathrm{𝟏𝟎}`$ representation, and the left-handed lepton doublet, $`L_L`$, fills out the $`\mathrm{𝟓}^{}`$ representation, so they should have the quantum number assignments
$`\begin{array}{cccc}i=& 1& 2& 3\\ \overline{e}_R:& (4,1)& (2,0)& (0,0)\\ L_L:& (2,0)& (1,0)& (0,1)\end{array},`$ (45)
These assignments dictated by the Eq. (44) can be compared against constraints obtained from experiment for masses and mixings in the lepton sector. The first phenomenological constraints we consider involve the charged leptons. Using the $`U(1)\times Z_2`$ quantum numbers in Eq. (45), one immediately obtains the charged lepton Yukawa matrix (see the Appendix for formulae),
$`m_\mathrm{}^\pm \left(\begin{array}{ccc}\lambda ^7& \lambda ^4& \lambda ^2\\ \lambda ^6& \lambda ^3& \lambda \\ \lambda ^4& \lambda ^3& \lambda \end{array}\right)v_1.`$ (46)
As desired the and elements are the same order of magnitude. This yields the mass ratios
$`{\displaystyle \frac{m_\mu }{m_\tau }}\lambda ^2,{\displaystyle \frac{m_e}{m_\tau }}\lambda ^4,`$ (47)
which are consistent with the experimental constraints after including renormalization group scalingThe largest scaling effect results from the additional running necessary to reach the muon and the electron mass scales so that one can relate the Yukawa couplings to the physical masses of the charged leptons. The scaling of the Yukawa coupling ratios themselves is negligibly small..
Next consider the light neutrino mass matrix. There are two possibilities that were discussed previously in Ref. . First the light neutrino mass matrix might not have suppressed entries arising from the $`Z_2`$ component of the flavor symmetry, in which case the light neutrino mass matrix is simply given by Eq. (76), where $`L_i`$ in this case is simply the sum of the $`U(1)`$ and $`Z_2`$ quantum numbers of the relevant lepton doublet field, $`L_L`$. For the charge assignments in Eq. (45), this gives
$`m_\nu \left(\begin{array}{ccc}\lambda ^4& \lambda ^3& \lambda ^3\\ \lambda ^3& \lambda ^2& \lambda ^2\\ \lambda ^3& \lambda ^2& \lambda ^2\end{array}\right){\displaystyle \frac{v_2^2}{\mathrm{\Lambda }_L}}.`$ (48)
The remaining constraints on leptons involve the neutrino masses and mixings. The most interesting aspect of the neutrino data is that the atmospheric neutrino mixing appears to be large, perhaps even maximal. As mentioned earlier, it is difficult to understand a hierarchical pattern for the neutrino masses, since large mixing should result when the neutrino masses are of roughly the same order of magnitude. The Super-Kamiokande data suggest that
$`\mathrm{\Delta }m_{23}^22.2\times 10^3\mathrm{eV}^2,\mathrm{sin}^22\theta _{23}^\nu 1,`$ (49)
where the subscripts indicate the generations of neutrinos involved in the mixing (we assume the mixing is between $`\nu _\mu `$ and $`\nu _\tau `$, and not some sterile neutrino).
The solar neutrino flux can be explained by one of three distinct solutions. Two of these involve matter-enhanced oscillation (MSW), while the third involves vacuum oscillations (VO). The two MSW solutions are differentiated by the size of the mixing angle, so one is usually called the small mixing angle (SMA) solution, and the other is called the large mixing angle (LMA) solution. The values required for the mixing parameters in each of these three cases are shown in the table below.
$`\begin{array}{ccc}& \mathrm{\Delta }m_{1x}^2[eV^2]& \mathrm{sin}^22\theta _{1x}\\ \mathrm{MSW}(\mathrm{SMA})& 5\times 10^6& 6\times 10^3\\ \mathrm{MSW}(\mathrm{LMA})& 2\times 10^5& 0.8\\ \mathrm{VO}& 8\times 10^{11}& 0.8\end{array}`$ (50)
The MSW solutions can be obtained with a $`Z_2`$ horizontal symmetry. If the neutrino masses are arranged in a hierarchy, then the best fit to the data is
$`{\displaystyle \frac{\mathrm{\Delta }m_{12}^2}{\mathrm{\Delta }m_{23}^2}}\lambda ^4,\mathrm{sin}\theta _{23}^\nu \lambda ^0,`$ (51)
and either
$`\mathrm{sin}\theta _{12}^\nu \lambda ^2,`$ (52)
for the SMA solution, or
$`\mathrm{sin}\theta _{12}^\nu \lambda ^0,`$ (53)
for the LMA solution. If no $`Z_2`$ symmetry is operative, one gets a light neutrino mass as in Eq. (76), and if $`L_2=L_3`$, there is no natural explanation for the hierarchy in the masses of $`m_{\nu _\mu }`$ and $`m_{\nu _\tau }`$. As explained in Refs. , this can be remedied by assigning the right-handed neutrino fields $`\overline{\nu }_{Ri}`$ (singlets of $`SU(5)`$) the following $`Z_2`$ charges: (0,0,1). The particular $`U(1)`$ assignment for the fields $`\overline{\nu }_{Ri}`$ does not affect the light neutrino mass matrix. In this case, the element of the $`m_\nu `$ matrix is enhanced by a factor $`\lambda ^2`$, giving
$`m_\nu \left(\begin{array}{ccc}\lambda ^4& \lambda ^3& \lambda ^3\\ \lambda ^3& \lambda ^2& \lambda ^2\\ \lambda ^3& \lambda ^2& 1\end{array}\right){\displaystyle \frac{v_2^2}{\mathrm{\Lambda }_L}},`$ (54)
for the charge assignment in Eq. (45).
The neutrino mixing matrix is
$`\left(\begin{array}{ccc}1& \lambda & \lambda \\ \lambda & 1& 1\\ \lambda & 1& 1\end{array}\right).`$ (55)
The solar mixing angle is predicted to be of order $`\lambda `$, falling in between the optimal value for the LMA solution ($`\lambda ^0`$) and the SMA solution ($`\lambda ^2`$). Equation (7) yields a solar mixing angle of order $`\lambda ^3`$, so the presence of the $`Z_2`$ symmetry has the effect in the neutrino sector of enhance the mixing of the first generation to the second and third generations by a factor $`\lambda ^2`$. Several unknown order-one coefficients combine to produce the matrix in Eq. (54), so it is not necessarily inconsistent with the MSW solutions.
In the models described here, one can achieve alignment of the quark mass matrices and the squark mass-squared matrices by certain positioning of texture zeros in the quark Yukawa matrices. This alignment can solve the SUSY flavor problem by making it possible to simultaneously diagonalize the quark mass matrices and the quark-squark-gluino coupling, thereby avoiding the dangerous flavor-changing couplings. In particular, in the models we are discussing here, one can achieve this alignment if there are texture zeros in the down quark Yukawa matrix, $`𝐃`$, in the and elements, and in either the or elements, and in either the or elements. This is easily seen to be the case after a quick inspection of Eqs. (35)-(38): in this case the Cabibbo angle, $`|V_{us}|`$, arises to leading order solely in the up quark Yukawa matrix, $`𝐔`$. The texture patterns that achieve this alignment occur when the down-quark Yukawa matrix has texture zeros in the positions given by the patterns
$`\left(\begin{array}{ccc}X& 0& 0\\ 0& 0& X\\ 0& X& X\end{array}\right)`$ (56)
(57)
$`\left(\begin{array}{ccc}0& 0& X\\ 0& X& 0\\ X& 0& X\end{array}\right)`$ (58)
which have the off-diagonal elements in the \[1-2\] block doubly suppressed. The off-diagonal suppression in the \[1-3\] block in the case of pattern Eq. (56) or the \[2-3\] block in the case of the pattern Eq. (58) need to be doubly suppressed, which is impossible. So one cannot achieve the quark-squark alignment in the context of a $`U(1)\times Z_2`$ flavor symmetry. On the other hand, one can employ the idea of supersymmetry breaking through an anomalous flavor symmetry to the grand unified model presented in this section. One can obtain reasonable suppression of the flavor-changing effects provided the first and second generation sparticles are in the multi-TeV range.
## VI Conclusion
We have shown that if the fermion mass matrices are dictated by an Abelian family symmetry, one can obtain a phenomenologically successful texture pattern by employing additional $`Z_m`$ horizontal symmetries. This four-texture zero model has a flavor symmetry that commutes with an SU(5) grand unified theory with the usual assignment of particles to the $`\mathrm{𝟓}^{}`$ and $`\mathrm{𝟏𝟎}`$ representations. When the quantum numbers are extended to the lepton sector, the charged lepton mass ratios were correctly predicted and a large mixing angle naturally arises to explain the atmospheric neutrino data. A mixing angle of order $`\lambda `$ arises to explain the solar neutrino oscillation data.
Discrete flavor symmetries can suppress entries in the Yukawa matrices and offer the potential of a solution to the supersymmetric flavor problem. A judicious choice of texture zeros can render the quark mass matrices and the squark mass-squared matrices simultaneously diagonalizable, thereby eliminating some strongly constrained flavor-changing couplings. However, we find that this solution cannot be obtained in a model with a single $`Z_2`$ symmetry and satisfy all the other (masses and mixings) phenomenological requirements. However the quantum number assignments can be compatible through suppression of flavor-changing effects when supersymmetry breaking is mediated by the anomalous flavor symmetry.
## Appendix
In this appendix we review the formulae for Yukawa and mass matrices that result from Abelian horizontal symmetries with and without a discrete component. Let the $`U(1)`$ quark charges be given by
$`\begin{array}{ccccccccc}Q_{L1}& Q_{L2}& Q_{L3}& \overline{u}_{R1}& \overline{u}_{R2}& \overline{u}_{R3}& \overline{d}_{R1}& \overline{d}_{R2}& \overline{d}_{R3}\\ q_1& q_2& q_3& u_1& u_2& u_3& d_1& d_2& d_3\end{array}.`$ (61)
Then the up and down quark Yukawa matrices, $`𝐔`$ and $`𝐃`$ are given byWe use the notation $`ab`$ to indicate that $`a`$ and $`b`$ are the same order in $`\lambda `$.
$`𝐔\left(\begin{array}{ccc}\lambda ^{q_1+u_1}& \lambda ^{q_1+u_2}& \lambda ^{q_1+u_3}\\ \lambda ^{q_2+u_1}& \lambda ^{q_2+u_2}& \lambda ^{q_2+u_3}\\ \lambda ^{q_3+u_1}& \lambda ^{q_3+u_2}& \lambda ^{q_3+u_3}\end{array}\right),𝐃\left(\begin{array}{ccc}\lambda ^{q_1+d_1}& \lambda ^{q_1+d_2}& \lambda ^{q_1+d_3}\\ \lambda ^{q_2+d_1}& \lambda ^{q_2+d_2}& \lambda ^{q_2+d_3}\\ \lambda ^{q_3+d_1}& \lambda ^{q_3+d_2}& \lambda ^{q_3+d_3}\end{array}\right).`$ (62)
It is understood that these matrices have unknown coefficients multiplying each element. The contributions to each element arise from a different operator in Eq. (6), so they are in general independent of each other. Since these coefficients are not correlated there is no reason to expect the Yukawa matrices to have a zero eigenvalue.
To compare the predictions of flavor symmetries to these phenomenological constraints, one has to relate the CKM elements to the entries in the Yukawa matrices. The Yukawa matrices $`𝐔`$ and $`𝐃`$ can be diagonalized by biunitary transformations
$`𝐔^{\mathrm{𝐝𝐢𝐚𝐠}}`$ $`=`$ $`V_u^L𝐔V_u^R,`$ (63)
$`𝐃^{\mathrm{𝐝𝐢𝐚𝐠}}`$ $`=`$ $`V_d^L𝐃V_d^R.`$ (64)
The CKM matrix is then given by
$$VV_u^LV_d^L.$$
(65)
The left-handed transformation matrices $`V_u^L`$ and $`V_d^L`$ can be defined in terms of three successive rotations in the (2,3), (1,3) and (1,2) sectors. These rotation angles of the transformation matrices can be expressed in terms of the elements of the Yukawa matrices as follows
$`s_{12}^u`$ $`=`$ $`{\displaystyle \frac{u_{12}}{\stackrel{~}{u}_{22}}}+{\displaystyle \frac{u_{11}u_{21}^{}}{|\stackrel{~}{u}_{22}^{}|^2}}{\displaystyle \frac{u_{13}(u_{32}+u_{23}^{}u_{22})}{\stackrel{~}{u}_{22}}}{\displaystyle \frac{u_{11}u_{31}^{}(u_{23}^{}+u_{32}u_{22}^{})}{|\stackrel{~}{u}_{22}^{}|^2}},`$ (66)
$`s_{13}^u`$ $`=`$ $`u_{13}+u_{11}u_{31}^{}+u_{12}(u_{32}^{}+u_{22}^{}u_{23})+u_{11}u_{21}^{}(u_{23}+u_{22}u_{32}^{}),`$ (67)
$`s_{23}^u`$ $`=`$ $`u_{23}+u_{22}u_{32}^{},`$ (68)
where $`u_{ij}`$ is the $`i,j`$th component of the up quark Yukawa matrix, $`𝐔/(𝐔)_{33}`$, and $`\stackrel{~}{u}_{22}=u_{22}u_{33}u_{23}u_{32}`$. There are corresponding expressions for the $`s_{ij}^d`$ in terms of the components of the down quark Yukawa matrix, $`𝐃`$ (which are slightly more complicated due to the fact that the (2,3) sector mixing in $`V_d^R`$ might be of order one). Clearly contributions to the CKM matrix elements can come from a number of terms. In this paper we are interested in determining only the leading order contribution(s) to the CKM angles and the fermion masses.
Assume now that the lepton fields have charges under a $`U(1)`$ family symmetry
$`\begin{array}{ccccccccc}\overline{e}_{R1}& \overline{e}_{R2}& \overline{e}_{R3}& \mathrm{}_{L1}& \mathrm{}_{L2}& \mathrm{}_{L3}& \overline{\nu }_{R1}& \overline{\nu }_{R2}& \overline{\nu }_{R3}\\ E_1& E_2& E_3& L_1& L_2& L_3& 𝒩_1& 𝒩_2& 𝒩_3\end{array}.`$ (71)
All the flavor charges are non-negative so holomorphic zeros do not play a role. The only suppressed entries will arise because of a discrete component in the flavor symmetry via a mechanism described below.
Given lepton doublet charges $`L_i`$ and right-handed neutrino charges $`𝒩_i`$ one has the following pattern for the charged lepton matrix
$`m_\mathrm{}^\pm \left(\begin{array}{ccc}\lambda ^{L_1+E_1}& \lambda ^{L_1+E_2}& \lambda ^{L_1+E_3}\\ \lambda ^{L_2+E_1}& \lambda ^{L_2+E_2}& \lambda ^{L_2+E_3}\\ \lambda ^{L_3+E_1}& \lambda ^{L_3+E_2}& \lambda ^{L_3+E_3}\end{array}\right)v_1,`$ (72)
and for the neutrino Dirac mass matrix
$`m_D\left(\begin{array}{ccc}\lambda ^{L_1+𝒩_1}& \lambda ^{L_1+𝒩_2}& \lambda ^{L_1+𝒩_3}\\ \lambda ^{L_2+𝒩_1}& \lambda ^{L_2+𝒩_2}& \lambda ^{L_2+𝒩_3}\\ \lambda ^{L_3+𝒩_1}& \lambda ^{L_3+𝒩_2}& \lambda ^{L_3+𝒩_3}\end{array}\right)v_2,`$ (73)
We have defined there the VEVs of the Higgs coupling to the down- and up-type quarks to be $`v_1`$ and $`v_2`$, and one usually defines $`\mathrm{tan}\beta =v_2/v_1`$. To determine the neutrino mixing angles one rotates to a basis where the charged lepton matrix is diagonal. This will give a contribution to the mixing in the light neutrino species. The relevant mixing contributing to atmospheric neutrino oscillations comes from the right hand side of the charge lepton matrix, $`\lambda ^{L_2+E_3}/\lambda ^{L_3+E_3}`$.
The Majarona mass matrix is obtained from the charges of the right-handed neutrino flavor charges $`𝒩_i`$ and a heavy scale we lable as $`\mathrm{\Lambda }_L`$,
$`M_N\left(\begin{array}{ccc}\lambda ^{2𝒩_1}& \lambda ^{𝒩_1+𝒩_2}& \lambda ^{𝒩_1+𝒩_3}\\ \lambda ^{𝒩_1+𝒩_2}& \lambda ^{2𝒩_2}& \lambda ^{𝒩_2+𝒩_3}\\ \lambda ^{𝒩_1+𝒩_3}& \lambda ^{𝒩_2+𝒩_3}& \lambda ^{2𝒩_3}\end{array}\right)\mathrm{\Lambda }_L.`$ (74)
Then one obtains the following form for the light neutrino mass matrix via the see-saw formula
$$m_\nu =m_D\frac{1}{M_N}m_D^T,$$
(75)
where $`m_D`$ is the neutrino Dirac mass matrix. Then,
$`m_\nu \left(\begin{array}{ccc}\lambda ^{2L_1}& \lambda ^{L_1+L_2}& \lambda ^{L_1+L_3}\\ \lambda ^{L_1+L_2}& \lambda ^{2L_2}& \lambda ^{L_2+L_3}\\ \lambda ^{L_1+L_3}& \lambda ^{L_2+L_3}& \lambda ^{2L_3}\end{array}\right){\displaystyle \frac{v_2^2}{\mathrm{\Lambda }_L}}.`$ (76)
If $`L_2=L_3`$ one can obtain $`𝒪(1)`$ mixing in the 2-3 sector. On the other hand, one fails to get a mass hierarchy between the second and third generation, since the two mass eigenvalues for the second and third generations are both of order $`\lambda ^{2L_3}`$.
A discrete Abelian family symmetry can be employed to enhance or suppress masses and mixing angle relative to the predictions obtained when the family symmetry is the continuous $`U(1)`$ symmetry, and this idea was pursued further in specific models. The discrete $`Z_m`$ symmetry can result in the enhancement of entries in the light neutrino mass matrix, and this enhancement is compatible with the neutrino seesaw mechanism. For example, if the $`U(1)`$ quantum numbers in Eq. (71) are replaced by $`U(1)\times Z_2`$ quantum numbers, $`L_3(L_31,1)`$ and $`𝒩_3(𝒩_31,1)`$ so that the charges for the lepton fields are
$`\begin{array}{ccccccccc}\overline{e}_{R1}& \overline{e}_{R2}& \overline{e}_{R3}& \mathrm{}_{L1}& \mathrm{}_{L2}& \mathrm{}_{L3}& \overline{\nu }_{R1}& \overline{\nu }_{R2}& \overline{\nu }_{R3}\\ (E_1,0)& (E_2,0)& (E_3,0)& (L_1,0)& (L_2,0)& (L_31,1)& (𝒩_1,0)& (𝒩_2,0)& (𝒩_31,1)\end{array}.`$ (79)
then one finds that
$`M_N\left(\begin{array}{ccc}\lambda ^{2𝒩_1}& \lambda ^{𝒩_1+𝒩_2}& \lambda ^{𝒩_1+𝒩_3}\\ \lambda ^{𝒩_1+𝒩_2}& \lambda ^{2𝒩_2}& \lambda ^{𝒩_2+𝒩_3}\\ \lambda ^{𝒩_1+𝒩_3}& \lambda ^{𝒩_2+𝒩_3}& \lambda ^{2𝒩_32}\end{array}\right)\mathrm{\Lambda }_L,`$ (80)
so that
$`(M_N)^1\left(\begin{array}{ccc}\lambda ^{2𝒩_1}& \lambda ^{𝒩_1𝒩_2}& \lambda ^{𝒩_1𝒩_3+2}\\ \lambda ^{𝒩_1𝒩_2}& \lambda ^{2𝒩_2}& \lambda ^{𝒩_2𝒩_3+2}\\ \lambda ^{𝒩_1𝒩_3+2}& \lambda ^{𝒩_2𝒩_3+2}& \lambda ^{2𝒩_3+2}\end{array}\right)\mathrm{\Lambda }_L^1.`$ (81)
So the effect of the discrete symmetry in our case is to enhance the 3-3 entry of the $`M_N`$ matrix, and thereby alter the results for the third row and the third column on the inverse matrix, $`(M_N)^1`$. The 3-3 component of the neutrino Dirac mass matrix is also enhanced by the discrete symmetry, so that Eq. (73) is modified to be
$`m_D\left(\begin{array}{ccc}\lambda ^{L_1+𝒩_1}& \lambda ^{L_1+𝒩_2}& \lambda ^{L_1+𝒩_3}\\ \lambda ^{L_2+𝒩_1}& \lambda ^{L_2+𝒩_2}& \lambda ^{L_2+𝒩_3}\\ \lambda ^{L_3+𝒩_1}& \lambda ^{L_3+𝒩_2}& \lambda ^{L_3+𝒩_32}\end{array}\right)v_2,`$ (82)
The light neutrino mass matrix in Eq. (76) is modified so that only the 3-3 entry is enhanced,
$`m_\nu \left(\begin{array}{ccc}\lambda ^{2L_1}& \lambda ^{L_1+L_2}& \lambda ^{L_1+L_3}\\ \lambda ^{L_1+L_2}& \lambda ^{2L_2}& \lambda ^{L_2+L_3}\\ \lambda ^{L_1+L_3}& \lambda ^{L_2+L_3}& \lambda ^{2L_32}\end{array}\right){\displaystyle \frac{v_2^2}{\mathrm{\Lambda }_L}}.`$ (83)
The charged lepton mass matrix, Eq. (72), and hence a large mixing angle is needed to diagonalize the \[2-3\] block. So the large mixing observed in the atmospheric neutrino experiments is accounted for, while the hierarchy of order $`\lambda ^2`$ in the second and third generation neutrino masses is obtained. Generalizing to a discrete symmetry $`Z_m`$ rather than $`Z_2`$ can preserve the large neutrino mixing while enhance the heaviest neutrino mass by a factor $`\lambda ^m`$.
## Acknowledgments
This work was supported in part by the U.S. Department of Energy under Grant No. No. DE-FG02-91ER40661.
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# Redshifts of 10 quasar candidates in the field of the rich absorption line quasar Q0122-38011footnote 1Based on optical spectroscopy obtained at the European Southern Observatory, La Silla, Chile
## 1 Introduction
The existence of large scale structure at high redshift ($`z1`$) provides an important constraint on theories for the formation of structure and evolution of the Universe. One approach to probing for such structure is through the study of intervening metal line absorption systems in quasar spectra. Such systems are believed to trace galaxies through their extended gaseous halos.
Statistical analysis of the redshift distribution of available samples of quasar absorption systems suggest that large scale clustering on comoving scales up to $``$100 Mpc may have been in place already at $`z`$$``$2–3 (e.g. Quashnock et al. QVY96 (1996)).
The complementary technique of probing for large scale structure in the plane of the sky by searching for correlated absorption in adjacent lines of sight is hampered by the relatively low density of high redshift quasars bright enough for detailed absorption line work. Nonetheless, several potential high redshift ’absorption superclusters’ spanning tens of Mpc on the sky have been identified in this manner. These include the $`z1.65`$ absorption systems seen toward PKS 0273$``$233 (Foltz et al. Fetal93 (1993)); the two pairs of damped Ly$`\alpha `$ systems seen at $`z2.38`$ and $`z2.85`$ toward Q2138$``$4427 and Q2139$``$4434 (Francis & Hewitt FH93 (1993); Francis et al. Fetal96 (1996)); the apparent structures at $`z2.3`$ and $`z2.5`$ detected in a dense quasar field near the south Galactic pole by Williger et al. (Wetal96 (1996)); and the well-studied case of the strong absorption spanning $`1.8z2.2`$ in the field of the quasar pair Tol 1037$``$2703/1038$``$2712 (Jakobsen et al. Jetal86 (1986); Dinshaw & Impey DI96 (1996); Lespine & Petitjean LP97 (1997); and references therein).
With the aim of searching for further such cases of intervening high-redshift superclusters Romani et al. (RFS91 (1991)) searched the quasar catalogs for suitable background objects near quasars known to display rich metal line absorption systems. One of the most promising fields identified by Romani et al. is that of the $`z=2.181`$ quasar Q0122$``$380, an object whose absorption spectrum contains at least seven C iv systems between $`1.81z1.98`$ (Carswell et al. Cetal82 (1982)) and happens to lie within a field in which Savage et al. (Setal84 (1984)) have carried out a deep objective prism quasar search. Q0122$``$380 is therefore surrounded by 11 quasar candidates of brightness $`V`$19–20 within a $`1^{}`$ radius, corresponding to comoving distances $`D_{}44h^1`$ Mpc at $`z1.95`$.
In this paper we present exploratory slit spectra and redshifts of these quasar candidates. As it turns out, several of the objects of interest are either not confirmed as quasars, or lie at significantly lower redshift than indicated by their preliminary catalog entries, thereby rendering the field toward Q0122$``$380 rather less promising for the purpose of searching for high redshift superclusters than originally thought.
## 2 Observations
Table 1 lists the coordinates, magnitudes and preliminary redshifts $`z^{}`$ of all objects listed in the Hewitt & Burbidge (1993) catalog and located within a circle of $`1^{}`$ radius centered on Q0122$``$380. These 11 objects along with Q0122$``$380 itself were observed with the ESO 3.6 m telescope and ESO Faint Object Camera and Spectrograph (EFOSC1) on 1995 September 28. EFOSC1 was equipped with a thinned back-side illuminated TEK CCD with 512$`\times `$512, 27$`\mu `$m pixels. We used a 230 Å mm<sup>-1</sup> (B300) grism in combination with a 2″ wide slit to obtain spectra covering the wavelength range 3750–6950 Å at $``$20 Å (FWHM) resolution. The spatial resolution is 0.61″ pixel<sup>-1</sup>.
The seeing conditions on both nights were poor. Nevertheless, the high throughput of EFOSC1 ensured good signal-to-noise (S/N) spectra of objects as faint as $`m_\mathrm{V}=20`$ (S/N$``$8 in the continuum near 5500 Å). For most objects we took a single 1200 s exposure. For 0117$``$379, being the faintest target at $`m_\mathrm{V}=20.0`$, we exposed for 1800 s. Prior to each spectroscopic observation, EFOSC1 was used in direct (filterless) imaging mode for target verification and automated slit acquisition. One quasar candidate, 0117$``$378, was not found at or near its catalog coordinates.
The data were reduced within the IRAF environment, following standard techniques. The spectra were put on a relative flux scale based on the standard stars LDS 235/EG 63 and LTT 2415 (Baldwin & Stone 1984; Stone & Baldwin 1983). The absolute calibration is ill determined, due to variations in seeing and atmospheric extinction over the night.
The resulting spectra of the 8 confirmed quasars are shown in Fig. 1 along with the relevant emission line identifications. The measured redshifts given in Table 1 were determined by averaging the redshifts measured for individual emission lines (e.g. Ly$`\alpha `$, N v $`\lambda `$1240, C ii $`\lambda `$1335, Si iv/O iv$`\lambda `$1400, C iv $`\lambda `$1549, He ii $`\lambda `$1640, Al iii $`\lambda `$1857, C iii$`\lambda `$1909, and Mg ii $`\lambda `$2799). For the wavelength of an emission line we adopted the average of the wavelengths of the peak (maximum signal) and of the center of a gaussian fit to the line, both determined after subtraction of the quasar continuum.
Two of the quasar candidates, 0121$``$379 and 0123$``$372, turn out to be Galactic stars. Their spectra are shown in Fig. 2. Unfortunately, of the observed quasar candidates 0121$``$379 would have been the quasar closest to Q0122$``$380, at an angular separation of $`10\stackrel{}{.}4`$.
## 3 Discussion
Romani et al. (1991) originally drew attention to the rich absorption line quasar Q0122$``$380 on the basis that it was surrounded by a number of quasar candidates closer than $`1^{}`$. Of the 11 objects listed in Table 1, seven had preliminary redshifts $`z>1.8`$, placing them near or behind the absorption seen between $`1.81z1.98`$ toward Q0122$``$380.
Of these seven high redshift candidates, only three (Q0117$``$380, Q0120$``$3781 and 0125$``$376) are confirmed as $`z>1.8`$ quasars. Of the remainder, one (0117$``$378) could not be located, two (0121$``$379 and 0123$``$372) are identified as stars, and another (Q0120$``$3785) turns out to be at a lower redshift of $`z=1.52`$. The four quasar candidates with lower preliminary redshifts $`z<1.8`$ are all confirmed as such.
As is evident from Fig. 1, our low resolution spectra of two of the three confirmed $`z>1.8`$ quasars (Q0120$``$3781 and Q0117$``$380) reveal no obvious C iv absorption features in the wavelength range 4350-4600 Å that could potentially be associated with the absorption spanning $`1.81z1.97`$ toward Q0122$``$380. However, the line detection limit of our spectra is only $`W_\lambda 8`$ Å at these wavelengths.
A more promising case is that of the final object, Q0125$``$376, whose redshift of $`z=1.868`$ lies close to that of the absorption complex at $`z1.91`$ in Q0122$``$380. Moreover, our spectrum of Q0125$``$376 (which has a better S/N ratio than those of Q0120$``$3781 and Q0117$``$380) shows three strong ($`W_\lambda 9`$ Å) absorption features at $`\lambda 4395`$ Å, $`\lambda 4435`$ Å and $`\lambda 4275`$ Å, respectively. The former two features are almost certainly due to C iv absorption from two $`z_{\mathrm{abs}}z_{\mathrm{em}}`$ systems at $`z1.837`$ and $`z1.864`$, an interpretation that is further strengthened by plausible detections of matching Si iv lines at shorter wavelengths. Based on a matching weaker feature seen at the anticipated position of Al iii, we tentatively identify the third feature as the C iv trough of a weak and possibly detached BAL-like complex at $`z1.76`$.
The two $`z_{\mathrm{abs}}z_{\mathrm{em}}`$ systems seen in Q0125$``$376 could conceivably be associated with the absorption seen toward Q0122$``$380, falling squarely between the $`z1.91`$ complex and the weaker $`z1.814`$ system detected in that object by Carswell et al. (1982). In a standard cosmological model with $`q_0=0.5`$ the angular distance between Q0125$``$376 and Q0122$``$380 of $`50\stackrel{}{.}2`$ at $`z=1.9`$, corresponds to a projected comoving separation of D$`{}_{}{}^{}=36h^1`$ Mpc $`(\text{H}\text{0}=h\mathrm{\hspace{0.17em}100}\text{km s}\text{-1}\text{ Mpc}\text{-1})`$, which is comparable to the extent of local superclusters of galaxies.
On the other hand, the presence of possibly BAL-like absorption at lower redshift would argue that the two $`z_{\mathrm{abs}}z_{\mathrm{em}}`$ systems seen in Q0125$``$376 are intrinsic in nature.
While higher resolution observations would be required to further delineate these possibilities and properly map the absorption toward the three confirmed $`z>1.8`$ quasars above, the exploratory observations presented here already make it clear that the field surrounding Q0122$``$380 is not as promising for searching for large scale structure at high redshift as had initially been hoped.
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# 1 Problem
## 1 Problem
Deduce the emissive power of radiation of frequency $`\nu `$ into vacuum at angle $`\theta `$ to the normal to the surface of a good conductor at temperature $`T`$, for polarization both parallel and perpendicular to the plane of emission.
## 2 Solution
The solution is adapted from ref. (see also ), and finds application in the calibration of the polarization dependence of detectors for cosmic microwave background radiation .
Recall Kirchhoff’s law of heat radiation (as clarified by Planck ) that
$$\frac{P_\nu }{A_\nu }=K(\nu ,T)=\frac{h\nu ^3/c^2}{e^{h\nu /kT}1},$$
(1)
where $`P_\nu `$ is the emissive power per unit area per unit frequency interval (emissivity) and
$$A_\nu =1=1\left|\frac{E_{0r}}{E_{0i}}\right|^2$$
(2)
is the absorption coefficient $`(0A_\nu 1)`$, $`c`$ is the speed of light, $`h`$ is Plank’s constant and $`k`$ is Boltzmann’s constant. Also recall the Fresnel equations of reflection that
$$\frac{E_{0r}}{E_{0i}}|_{}=\frac{\mathrm{sin}(\theta _t\theta _i)}{\mathrm{sin}(\theta _t+\theta _i)},\frac{E_{0r}}{E_{0i}}|_{}=\frac{\mathrm{tan}(\theta _t\theta _i)}{\mathrm{tan}(\theta _t+\theta _i)},$$
(3)
where $`i`$, $`r`$, and $`t`$ label the incident, reflected, and transmitted waves, respectively.
The solution is based on the fact that eq. (1) holds separately for each polarization of the emitted radiation, and is also independent of the angle of the radiation. This result is implicit in Planck’s derivation of Kirchhoff’s law of radiation, and is stated explicitly in . That law describes the thermodynamic equilibrium of radiation emitted and absorbed throughout a volume. The emissivity $`P_v`$ and the absorption coefficient $`A_\nu `$ can depend on the polarization of the radiation and on the angle of the radiation, but the definitions of polarization parallel and perpendicular to a plane of emission, and of angle relative to the normal to a surface element, are local, while the energy conservation relation $`P_\nu =A_\nu K(\nu ,T)`$ is global. A “ray” of radiation whose polarization can be described as parallel to the plane of emission is, in general, a mixture of parallel and perpendicular polarization from the point of view of the absorption process. Similarly, the angles of emission and absorption of a ray are different in general. Thus, the concepts of parallel and perpendicular polarization and of the angle of the radiation are not well defined after integrating over the entire volume. Thermodynamic equilibrium can exist only if a single spectral intensity function $`K(\nu ,T)`$ holds independent of polarization and of angle.
All that remains is to evaluate the reflection coefficients $`_{}`$ and $`_{}`$ for the two polarizations at a vacuum-metal interface. These are well known , but we derive them for completeness.
To use the Fresnel equations (3), we need expressions for $`\mathrm{sin}\theta _t`$ and $`\mathrm{cos}\theta _t`$. The boundary condition that the phase of the wave be continuous across the vacuum-metal interface leads, as is well known, to the general form of Snell’s law:
$$k_i\mathrm{sin}\theta _i=k_t\mathrm{sin}\theta _t,$$
(4)
where $`k=2\pi /\lambda `$ is the wave number. Then,
$$\mathrm{cos}\theta _t=\sqrt{1\frac{k_i^2}{k_t^2}\mathrm{sin}^2\theta _i}.$$
(5)
To determine the relation between wave numbers $`k_i`$ and $`k_t`$ in vacuum and in the conductor, we consider a plane wave of angular frequency $`\omega =2\pi \nu `$ and complex wave vector k,
$$𝐄=𝐄_0e^{i(𝐤_𝐭𝐫\omega t)},$$
(6)
which propagates in a conducting medium with dielectric constant $`ϵ`$, permeability $`\mu `$, and conductivity $`\sigma `$. The wave equation for the electric field in such a medium is (in Gaussian units)
$$^2𝐄\frac{ϵ\mu }{c^2}\frac{^2𝐄}{t^2}=\frac{4\pi \mu \sigma }{c^2}\frac{𝐄}{t},$$
(7)
where $`c`$ is the speed of light. We find the dispersion relation for the wave vector $`k_t`$ on inserting eq. (6) in eq. (7):
$$k_t^2=ϵ\mu \frac{\omega ^2}{c^2}+i\frac{4\pi \sigma \mu \omega }{c^2}.$$
(8)
For a good conductor, the second term of eq. (8) is much larger than the first, so we write
$$k_t\frac{\sqrt{2\pi \sigma \mu \omega }}{c}(1+i)=\frac{1+i}{d}=\frac{2}{d(1i)},$$
(9)
where
$$d=\frac{c}{\sqrt{2\pi \sigma \mu \omega }}\lambda $$
(10)
is the frequency-dependent skin depth. Of course, on setting $`ϵ=1=\mu `$ and $`\sigma =0`$ we obtain expressions that hold in vacuum, where $`k_i=\omega /c`$.
We see that for a good conductor $`\left|k_t\right|k_i`$, so according to eq. (5) we may take $`\mathrm{cos}\theta _t1`$ to first order of accuracy in the small ratio $`d/\lambda `$. Then the first of the Fresnel equations becomes
$$\frac{E_{0r}}{E_{0i}}|_{}=\frac{\mathrm{cos}\theta _i\mathrm{sin}\theta _t/\mathrm{sin}\theta _i1}{\mathrm{cos}\theta _i\mathrm{sin}\theta _t/\mathrm{sin}\theta _i+1}=\frac{(k_i/k_t)\mathrm{cos}\theta _i1}{(k_i/k_t)\mathrm{cos}\theta _i+1}\frac{(\pi d/\lambda )(1i)\mathrm{cos}\theta _i1}{(\pi d/\lambda )(1i)\mathrm{cos}\theta _i+1},$$
(11)
and the reflection coefficient is approximated by
$$_{}=\left|\frac{E_{0r}}{E_{0i}}\right|_{}^21\frac{4\pi d}{\lambda }\mathrm{cos}\theta _i=12\mathrm{cos}\theta _i\sqrt{\frac{\nu }{\sigma }}.$$
(12)
For the other polarization, we see that
$$\frac{E_{0r}}{E_{0i}}|_{}=\frac{E_{0r}}{E_{0i}}|_{}\frac{\mathrm{cos}(\theta _i+\theta _t)}{\mathrm{cos}(\theta _i\theta _t)}\frac{E_{0r}}{E_{0i}}|_{}\frac{\mathrm{cos}\theta _i(\pi d/\lambda )(1i)\mathrm{sin}^2\theta _i}{\mathrm{cos}\theta _i+(\pi d/\lambda )(1i)\mathrm{sin}^2\theta _i},$$
(13)
so that
$$_{}_{}\left(1\frac{4\pi d}{\lambda }\frac{\mathrm{sin}^2\theta _i}{\mathrm{cos}\theta _i}\right)1\frac{4\pi d}{\lambda \mathrm{cos}\theta _i}=1\frac{2}{\mathrm{cos}\theta _i}\sqrt{\frac{\nu }{\sigma }}.$$
(14)
An expression for $`_{}`$ valid to second order in $`d/\lambda `$ has been given in ref. . For $`\theta _i`$ near $`90^{}`$, $`_{}1`$, but eq. (14) for $`_{}`$ is not accurate. Writing $`\theta _i=\pi /2\vartheta _i`$ with $`\vartheta _i1`$, eq. (13) becomes
$$\frac{E_{0r}}{E_{0i}}|_{}\frac{\vartheta _i(\pi d/\lambda )(1i)}{\vartheta _i+(\pi d/\lambda )(1i)},$$
(15)
For $`\theta _i=\pi /2`$, $`_{}=1`$, and $`_{,\mathrm{min}}=(5\sqrt{2})/(5+\sqrt{2})=0.58`$ for $`\vartheta _i=2\sqrt{2}\pi d/\lambda `$.
Finally, combining eqs. (1), (2), (12) and (14) we have
$$P_\nu \frac{4\pi d\mathrm{cos}\theta }{\lambda ^3}\frac{h\nu }{e^{h\nu /kT}1},P_\nu \frac{4\pi d}{\lambda ^3\mathrm{cos}\theta }\frac{h\nu }{e^{h\nu /kT}1},$$
(16)
and
$$\frac{P_\nu }{P_\nu }=\mathrm{cos}^2\theta $$
(17)
for the emissivities at angle $`\theta `$ such that $`\mathrm{cos}\theta d/\lambda `$.
The conductivity $`\sigma `$ that appears in eq. (16) can be taken as the dc conductivity so long as the wavelength exceeds 10 $`\mu `$m . If in addition $`h\nu kT`$, then eq. (16) can be written
$$P_\nu \frac{4\pi dkT\mathrm{cos}\theta }{\lambda ^3},P_\nu \frac{4\pi dkT}{\lambda ^3\mathrm{cos}\theta },$$
(18)
in terms of the skin depth $`d`$.
We would like to thank Matt Hedman, Chris Herzog and Suzanne Staggs for conversations about this problem.
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# An Elementary Introduction to Groups and Representations
### 0.1. Preface
These notes are the outgrowth of a graduate course on Lie groups I taught at the University of Virginia in 1994. In trying to find a text for the course I discovered that books on Lie groups either presuppose a knowledge of differentiable manifolds or provide a mini-course on them at the beginning. Since my students did not have the necessary background on manifolds, I faced a dilemma: either use manifold techniques that my students were not familiar with, or else spend much of the course teaching those techniques instead of teaching Lie theory. To resolve this dilemma I chose to write my own notes using the notion of a matrix Lie group. A matrix Lie group is simply a closed subgroup of $`\mathrm{𝖦𝖫}(n;).`$ Although these are often called simply “matrix groups,” my terminology emphasizes that every matrix group is a Lie group.
This approach to the subject allows me to get started quickly on Lie group theory proper, with a minimum of prerequisites. Since most of the interesting examples of Lie groups are matrix Lie groups, there is not too much loss of generality. Furthermore, the proofs of the main results are ultimately similar to standard proofs in the general setting, but with less preparation.
Of course, there is a price to be paid and certain constructions (e.g. covering groups) that are easy in the Lie group setting are problematic in the matrix group setting. (Indeed the universal cover of a matrix Lie group need not be a matrix Lie group.) On the other hand, the matrix approach suffices for a first course. Anyone planning to do research in Lie group theory certainly needs to learn the manifold approach, but even for such a person it might be helpful to start with a more concrete approach. And for those in other fields who simply want to learn the basics of Lie group theory, this approach allows them to do so quickly.
These notes also use an atypical approach to the theory of semisimple Lie algebras, namely one that starts with a detailed calculation of the representations of $`\mathrm{𝗌𝗅}(3;)`$. My own experience was that the theory of Cartan subalgebras, roots, Weyl group, etc., was pretty difficult to absorb all at once. I have tried, then, to motivate these constructions by showing how they are used in the representation theory of the simplest representative Lie algebra. (I also work out the case of $`\mathrm{𝗌𝗅}(2;),`$ but this case does not adequately illustrate the general theory.)
In the interests of making the notes accessible to as wide an audience as possible, I have included a very brief introduction to abstract groups, given in Chapter 1. In fact, not much of abstract group theory is needed, so the quick treatment I give should be sufficient for those who have not seen this material before.
I am grateful to many who have made corrections, large and small, to the notes, including especially Tom Goebeler, Ruth Gornet, and Erdinch Tatar.
## Chapter 1 Groups
### 1.1. Definition of a Group, and Basic Properties
###### Definition 1.1.
A group is a set $`G`$, together with a map of $`G\times G`$ into $`G`$ (denoted $`g_1g_2`$) with the following properties:
First, associativity: for all $`g_1,g_2G`$,
(1.1)
$$g_1(g_2g_3)=(g_1g_2)g_3\text{.}$$
Second, there exists an element $`e`$ in $`G`$ such that for all $`gG`$,
(1.2)
$$ge=eg=g\text{.}$$
and such that for all $`gG`$, there exists $`hG`$ with
(1.3)
$$gh=hg=e\text{.}$$
If $`gh=hg`$ for all $`g,hG`$, then the group is said to be commutative(or abelian).
The element $`e`$ is (as we shall see momentarily) unique, and is called the identity element of the group, or simply the identity. Part of the definition of a group is that multiplying a group element $`g`$ by the identity on either the right or the left must give back $`g`$.
The map of $`G\times G`$ into $`G`$ is called the product operation for the group. Part of the definition of a group $`G`$ is that the product operation map $`G\times G`$ into $`G`$, i.e., that the product of two elements of $`G`$ be again an element of $`G`$. This property is referred to as closure.
Given a group element $`g`$, a group element $`h`$ such that $`gh=hg=e`$ is called an inverse of $`g`$. We shall see momentarily that each group element has a unique inverse.
Given a set and an operation, there are four things that must be checked to show that this is a group: closure, associativity, existence of an identity, and existence of inverses.
###### Proposition 1.2 (Uniqueness of the Identity).
Let $`G`$ be a group, and let $`e,fG`$ be such that for all $`gG`$
$`eg`$ $`=ge=g`$
$`fg`$ $`=gf=g\text{.}`$
Then $`e=f`$.
###### Proof.
Since $`e`$ is an identity, we have
$$ef=f\text{.}$$
On the other hand, since $`f`$ is an identity, we have
$$ef=e\text{.}$$
Thus $`e=ef=f`$. ∎
###### Proposition 1.3 (Uniqueness of Inverses).
Let $`G`$ be a group, $`e`$ the (unique) identity of $`G`$, and $`g,h,k`$ arbitrary elements of $`G`$. Suppose that
$`gh`$ $`=hg=e`$
$`gk`$ $`=kg=e.`$
Then $`h=k`$.
###### Proof.
We know that $`gh=gk`$ $`(=e)`$. Multiplying on the left by $`h`$ gives
$$h(gh)=h(gk)\text{.}$$
By associativity, this gives
$$(hg)h=(hg)k\text{,}$$
and so
$`eh=ek`$
$`h=k\text{.}`$
This is what we wanted to prove.
###### Proposition 1.4.
Let $`G`$ be a group, $`e`$ the identity element of $`G`$, and $`g`$ an arbitrary element of $`G`$. Suppose $`hG`$ satisfies either $`hg=e`$ or $`gh=e`$. Then $`h`$ is the (unique) inverse of $`g`$.
###### Proof.
To show that $`h`$ is the inverse of $`g`$, we must show both that $`hg=e`$ and $`gh=e`$. Suppose we know, say, that $`hg=e`$. Then our goal is to show that this implies that $`gh=e`$.
Since $`hg=e`$,
$$g(hg)=ge=g\text{.}$$
By associativity, we have
$$(gh)g=g\text{.}$$
Now, by the definition of a group, $`g`$ has an inverse. Let $`k`$ be that inverse. (Of course, in the end, we will conclude that $`k=h`$, but we cannot assume that now.) Multiplying on the right by $`k`$ and using associativity again gives
$`((gh)g)k=gk=e`$
$`(gh)(gk)=e`$
$`(gh)e=e`$
$`gh=e\text{.}`$
A similar argument shows that if $`gh=e`$, then $`hg=e`$. ∎
Note that in order to show that $`hg=e`$ implies $`gh=e`$, we used the fact that $`g`$ has an inverse, since it is an element of a group. In more general contexts (that is, in some system which is not a group), one may have $`hg=e`$ but not $`gh=e`$. (See Exercise 11.)
###### Notation 1.5.
For any group element $`g`$, its unique inverse will be denoted $`g^1`$.
###### Proposition 1.6 (Properties of Inverses).
Let $`G`$ be a group, $`e`$ its identity, and $`g,h`$ arbitrary elements of $`G`$. Then
$`\left(g^1\right)^1=g`$
$`\left(gh\right)^1=h^1g^1`$
$`e^1=e\text{.}`$
###### Proof.
Exercise. ∎
### 1.2. Some Examples of Groups
From now on, we will denote the product of two group elements $`g_1`$ and $`g_2`$ simply by $`g_1g_2`$, instead of the more cumbersome $`g_1g_2`$. Moreover, since we have associativity, we will write simply $`g_1g_2g_3`$ in place of $`(g_1g_2)g_3`$ or $`g_1(g_2g_3)`$.
#### 1.2.1. The trivial group
The set with one element, $`e`$, is a group, with the group operation being defined as $`ee=e`$. This group is commutative.
Associativity is automatic, since both sides of (1.1) must be equal to $`e`$. Of course, $`e`$ itself is the identity, and is its own inverse. Commutativity is also automatic.
#### 1.2.2. The integers
The set $``$ of integers forms a group with the product operation being addition. This group is commutative.
First, we check closure, namely, that addition maps $`\times `$ into $``$, i.e., that the sum of two integers is an integer. Since this is obvious, it remains only to check associativity, identity, and inverses. Addition is associative; zero is the additive identity (i.e., $`0+n=n+0=n`$, for all $`n`$); each integer $`n`$ has an additive inverse, namely, $`n`$. Since addition is commutative, $``$ is a commutative group.
#### 1.2.3. The reals and $`^n`$
The set $``$ of real numbers also forms a group under the operation of addition. This group is commutative. Similarly, the $`n`$-dimensional Euclidean space $`^n`$ forms a group under the operation of vector addition. This group is also commutative.
The verification is the same as for the integers.
#### 1.2.4. Non-zero real numbers under multiplication
The set of non-zero real numbers forms a group with respect to the operation of multiplication. This group is commutative.
Again we check closure: the product of two non-zero real numbers is a non-zero real number. Multiplication is associative; one is the multiplicative identity; each non-zero real number $`x`$ has a multiplicative inverse, namely, $`\frac{1}{x}`$. Since multiplication of real numbers is commutative, this is a commutative group.
This group is denoted $`^{}`$.
#### 1.2.5. Non-zero complex numbers under multiplication
The set of non-zero complex numbers forms a group with respect to the operation of complex multiplication. This group is commutative.
This group in denoted $`^{}`$.
#### 1.2.6. Complex numbers of absolute value one under multiplication
The set of complex numbers with absolute value one (i.e., of the form $`e^{i\theta }`$) forms a group under complex multiplication. This group is commutative.
This group is the unit circle, denoted $`S^1`$.
#### 1.2.7. Invertible matrices
For each positive integer $`n`$, the set of all $`n\times n`$ invertible matrices with real entries forms a group with respect to the operation of matrix multiplication. This group in non-commutative, for $`n2`$.
We check closure: the product of two invertible matrices is invertible, since $`\left(AB\right)^1=B^1A^1`$. Matrix multiplication is associative; the identity matrix (with ones down the diagonal, and zeros elsewhere) is the identity element; by definition, an invertible matrix has an inverse. Simple examples show that the group is non-commutative, except in the trivial case $`n=1`$. (See Exercise 8.)
This group is called the general linear group (over the reals), and is denoted $`\mathrm{𝖦𝖫}(n;)`$.
#### 1.2.8. Symmetric group (permutation group)
The set of one-to-one, onto maps of the set $`\{1,2,\mathrm{}n\}`$ to itself forms a group under the operation of composition. This group is non-commutative for $`n3`$.
We check closure: the composition of two one-to-one, onto maps is again one-to-one and onto. Composition of functions is associative; the identity map (which sends 1 to 1, 2 to 2, etc.) is the identity element; a one-to-one, onto map has an inverse. Simple examples show that the group is non-commutative, as long as $`n`$ is at least 3. (See Exercise 10.)
This group is called the symmetric group, and is denoted $`S_n`$. A one-to-one, onto map of $`\{1,2,\mathrm{}n\}`$ is a permutation, and so $`S_n`$ is also called the permutation group. The group $`S_n`$ has $`n!`$ elements.
#### 1.2.9. Integers mod $`n`$
The set $`\{0,1,\mathrm{}n1\}`$ forms a group under the operation of addition mod $`n`$. This group is commutative.
Explicitly, the group operation is the following. Consider $`a,b\{0,1\mathrm{}n1\}`$. If $`a+b<n`$, then $`a+b`$ $`\mathrm{𝐦𝐨𝐝}`$ $`n=a+b`$, if $`a+bn`$, then $`a+b`$ $`\mathrm{𝐦𝐨𝐝}`$ $`n=a+bn`$. (Since $`a`$ and $`b`$ are less than $`n`$, $`a+bn`$ is less than $`n`$; thus we have closure.) To show associativity, note that both $`(a+b\mathrm{𝐦𝐨𝐝}n)+c\mathrm{𝐦𝐨𝐝}n`$ and $`a+(b+c\mathrm{𝐦𝐨𝐝}n)\mathrm{𝐦𝐨𝐝}n`$ are equal to $`a+b+c`$, minus some multiple of $`n`$, and hence differ by a multiple of $`n`$. But since both are in the set $`\{0,1,\mathrm{}n1\}`$, the only possible multiple on $`n`$ is zero. Zero is still the identity for addition mod $`n`$. The inverse of an element $`a\{0,1,\mathrm{}n1\}`$ is $`na`$. (Exercise: check that $`na`$ is in $`\{0,1,\mathrm{}n1\}`$, and that $`a+(na)\mathrm{𝐦𝐨𝐝}n=0`$.) The group is commutative because ordinary addition is commutative.
This group is referred to as “$`\mathrm{𝐦𝐨𝐝}n`$,” and is denoted $`_n`$.
### 1.3. Subgroups, the Center, and Direct Products
###### Definition 1.7.
A subgroup of a group $`G`$ is a subset $`H`$ of $`G`$ with the following properties:
1. The identity is an element of $`H`$.
2. If $`hH`$, then $`h^1H`$.
3. If $`h_1,h_2H`$, then $`h_1h_2H`$ .
The conditions on $`H`$ guarantee that $`H`$ is a group, with the same product operation as $`G`$ (but restricted to $`H`$). Closure is assured by (3), associativity follows from associativity in $`G`$, and the existence of an identity and of inverses is assured by (1) and (2).
#### 1.3.1. Examples
Every group $`G`$ has at least two subgroups: $`G`$ itself, and the one-element subgroup $`\left\{e\right\}`$. (If $`G`$ itself is the trivial group, then these two subgroups coincide.) These are called the trivial subgroups of $`G`$.
The set of even integers is a subgroup of $``$: zero is even, the negative of an even integer is even, and the sum of two even integers is even.
The set $`H`$ of $`n\times n`$ real matrices with determinant one is a subgroup of $`\mathrm{𝖦𝖫}(n;)`$. The set $`H`$ is a subset of $`\mathrm{𝖦𝖫}(n;)`$ because any matrix with determinant one is invertible. The identity matrix has determinant one, so 1 is satisfied. The determinant of the inverse is the reciprocal of the determinant, so 2 is satisfied; and the determinant of a product is the product of the determinants, so 3 is satisfied. This group is called the special linear group (over the reals), and is denoted $`\mathrm{𝖲𝖫}(n;)`$.
Additional examples, as well as some non-examples, are given in Exercise 2.
###### Definition 1.8.
The center of a group $`G`$ is the set of all $`gG`$ such that $`gh=hg`$ for all $`hG`$.
It is not hard to see that the center of any group $`G`$ is a subgroup $`G`$.
###### Definition 1.9.
Let $`G`$ and $`H`$ be groups, and consider the Cartesian product of $`G`$ and $`H`$, i.e., the set of ordered pairs $`(g,h)`$ with $`gG,hH`$. Define a product operation on this set as follows:
$$(g_1,h_1)(g_2,h_2)=(g_1g_2,h_1h_2)\text{.}$$
This operation makes the Cartesian product of $`G`$ and $`H`$ into a group, called the direct product of $`G`$ and $`H`$ and denoted $`G\times H`$.
It is a simple matter to check that this operation truly makes $`G\times H`$ into a group. For example, the identity element of $`G\times H`$ is the pair $`(e_1,e_2)`$, where $`e_1`$ is the identity for $`G`$, and $`e_2`$ is the identity for $`H`$.
### 1.4. Homomorphisms and Isomorphisms
###### Definition 1.10.
Let $`G`$ and $`H`$ be groups. A map $`\varphi :GH`$ is called a homomorphism if $`\varphi (g_1g_2)=\varphi (g_1)\varphi (g_2)`$ for all $`g_1,g_2G`$. If in addition, $`\varphi `$ is one-to-one and onto, then $`\varphi `$ is called an isomorphism. An isomorphism of a group with itself is called an automorphism.
###### Proposition 1.11.
Let $`G`$ and $`H`$ be groups, $`e_1`$ the identity element of $`G`$, and $`e_2`$ the identity element of $`H`$. If $`\varphi :GH`$ is a homomorphism, then $`\varphi (e_1)=e_2`$, and $`\varphi (g^1)=\varphi (g)^1`$ for all $`gG`$.
###### Proof.
Let $`g`$ be any element of $`G`$. Then $`\varphi (g)=\varphi (ge_1)=\varphi (g)\varphi (e_1)`$. Multiplying on the left by $`\varphi (g)^1`$ gives $`e_2=\varphi (e_1)`$. Now consider $`\varphi (g^1)`$. Since $`\varphi (e_1)=e_2`$, we have $`e_2=\varphi (e_1)=\varphi (gg^1)=\varphi (g)\varphi (g^1)`$. In light of Prop. 1.4, we conclude that $`\varphi (g^1)`$ is the inverse of $`\varphi (g)`$. ∎
###### Definition 1.12.
Let $`G`$ and $`H`$ be groups, $`\varphi :GH`$ a homomorphism, and $`e_2`$ the identity element of $`H`$. The kernel of $`\varphi `$ is the set of all $`gG`$ for which $`\varphi (g)=e_2`$.
###### Proposition 1.13.
Let $`G`$ and $`H`$ be groups, and $`\varphi :GH`$ a homomorphism. Then the kernel of $`\varphi `$ is a subgroup of $`G`$.
###### Proof.
Easy. ∎
#### 1.4.1. Examples
Given any two groups $`G`$ and $`H`$, we have the trivial homomorphism from $`G`$ to $`H`$: $`\varphi (g)=e`$ for all $`gG`$. The kernel of this homomorphism is all of $`G`$.
In any group $`G`$, the identity map ($`id(g)=g`$) is an automorphism of $`G`$, whose kernel is just $`\left\{e\right\}`$.
Let $`G=H=`$, and define $`\varphi (n)=2n`$. This is a homomorphism of $``$ to itself, but not an automorphism. The kernel of this homomorphism is just $`\left\{0\right\}`$.
The determinant is a homomorphism of $`\mathrm{𝖦𝖫}(n,)`$ to $`^{}`$. The kernel of this map is $`\mathrm{𝖲𝖫}(n,)`$.
Additional examples are given in Exercises 12 and 7.
If there exists an isomorphism from $`G`$ to $`H`$, then $`G`$ and $`H`$ are said to be isomorphic, and this relationship is denoted $`GH`$. (See Exercise 4.) Two groups which are isomorphic should be thought of as being (for all practical purposes) the same group.
### 1.5. Exercises
Recall the definitions of the groups $`\mathrm{𝖦𝖫}(n;)`$, $`S_n`$, $`^{}`$, and $`_n`$ from Sect. 1.2, and the definition of the group $`\mathrm{𝖲𝖫}(n;)`$ from Sect. 1.3.
1. Show that the center of any group $`G`$ is a subgroup $`G`$.
2. In (a)-(f), you are given a group $`G`$ and a subset $`H`$ of $`G`$. In each case, determine whether $`H`$ is a subgroup of $`G`$.
(a) $`G=,H=\left\{\text{odd integers}\right\}`$
(b) $`G=,H=\left\{\text{multiples of 3}\right\}`$
(c) $`G=\mathrm{𝖦𝖫}(n;),H=\left\{A\mathrm{𝖦𝖫}(n;)\right|detA\text{ is an integer}\}`$
(d) $`G=\mathrm{𝖲𝖫}(n;),H=\left\{A\mathrm{𝖲𝖫}(n;)\right|\text{all the entries of }A\text{ are integers}\}`$
Hint: recall Kramer’s rule for finding the inverse of a matrix.
(e) $`G=\mathrm{𝖦𝖫}(n;),H=\left\{A\mathrm{𝖦𝖫}(n;)\right|\text{all of the entries of }A\text{ are rational}\}`$
(f) $`G=_9,H=\{0,2,4,6,8\}`$
3. Verify the properties of inverses in Prop. 1.6.
4. Let $`G`$ and $`H`$ be groups. Suppose there exists an isomorphism $`\varphi `$ from $`G`$ to $`H`$. Show that there exists an isomorphism from $`H`$ to $`G`$.
5. Show that the set of positive real numbers is a subgroup of $`^{}`$. Show that this group is isomorphic to the group $``$.
6. Show that the set of automorphisms of any group $`G`$ is itself a group, under the operation of composition. This group is the automorphism group of $`G`$, $`Aut(G)`$.
7. Given any group $`G`$, and any element $`g`$ in $`G`$, define $`\varphi _g:GG`$ by $`\varphi _g(h)=ghg^1`$. Show that $`\varphi _g`$ is an automorphism of $`G`$. Show that the map $`g\varphi _g`$ is a homomorphism of $`G`$ into $`Aut(G)`$, and that the kernel of this map is the center of $`G`$.
Note: An automorphism which can be expressed as $`\varphi _g`$ for some $`gG`$ is called an inner automorphism; any automorphism of $`G`$ which is not equal to any $`\varphi _g`$ is called an outer automorphism.
8. Give an example of two $`2\times 2`$ invertible real matrices which do not commute. (This shows that $`\mathrm{𝖦𝖫}(2,𝐑)`$ is not commutative.)
9. Show that in any group $`G`$, the center of $`G`$ is a subgroup.
10. An element $`\sigma `$ of the permutation group $`S_n`$ can be written in two-row form,
$$\sigma =\left(\begin{array}{cccc}1& 2& \mathrm{}& n\\ \sigma _1& \sigma _2& \mathrm{}& \sigma _n\end{array}\right)$$
where $`\sigma _i`$ denotes $`\sigma (i)`$. Thus
$$\sigma =\left(\begin{array}{ccc}1& 2& 3\\ 2& 3& 1\end{array}\right)$$
is the element of $`S_3`$ which sends 1 to 2, 2 to 3, and 3 to 1. When multiplying (i.e., composing) two permutations, one performs the one on the right first, and then the one on the left. (This is the usual convention for composing functions.)
Compute
$$\left(\begin{array}{ccc}1& 2& 3\\ 2& 1& 3\end{array}\right)\left(\begin{array}{ccc}1& 2& 3\\ 1& 3& 2\end{array}\right)$$
and
$$\left(\begin{array}{ccc}1& 2& 3\\ 1& 3& 2\end{array}\right)\left(\begin{array}{ccc}1& 2& 3\\ 2& 1& 3\end{array}\right)$$
Conclude that $`S_3`$ is not commutative.
11. Consider the set $`=\{0,1,2,\mathrm{}\}`$ of natural numbers, and the set $``$ of all functions of $``$ to itself. Composition of functions defines a map of $`\times `$ into $``$, which is associative. The identity ($`id(n)=n`$) has the property that $`idf=fid=f`$, for all $`f`$ in $``$. However, since we do not restrict to functions which are one-to-one and onto, not every element of $``$ has an inverse. Thus $``$ is not a group.
Give an example of two functions $`f,g`$ in $``$ such that $`fg=id`$, but $`gfid`$. (Compare with Prop. 1.4.)
12. Consider the groups $``$ and $`_n`$. For each $`a`$ in $``$, define $`a\mathrm{𝐦𝐨𝐝}n`$ to be the unique element $`b`$ of $`\{0,1,\mathrm{}n1\}`$ such that $`a`$ can be written as $`a=kn+b`$, with $`k`$ an integer. Show that the map $`aa\mathrm{𝐦𝐨𝐝}n`$ is a homomorphism of $``$ into $`_n`$.
13. Let $`G`$ be a group, and $`H`$ a subgroup of $`G`$. $`H`$ is called a normal subgroup of $`G`$ if given any $`gG`$, and $`hH`$, $`ghg^1`$ is in $`H`$.
Show that any subgroup of a commutative group is normal. Show that in any group $`G`$, the trivial subgroups $`G`$ and $`\{e\}`$ are normal. Show that the center of any group is a normal subgroup. Show that if $`\varphi `$ is a homomorphism from $`G`$ to $`H`$, then the kernel of $`\varphi `$ is a normal subgroup of $`G`$.
Show that $`\mathrm{𝖲𝖫}(n;)`$ is a normal subgroup of $`\mathrm{𝖦𝖫}(n;)`$.
Note: a group $`G`$ with no normal subgroups other than $`G`$ and $`\left\{e\right\}`$ is called simple.
## Chapter 2 Matrix Lie Groups
### 2.1. Definition of a Matrix Lie Group
Recall that the general linear group over the reals, denoted $`\mathrm{𝖦𝖫}(n;)`$, is the group of all $`n\times n`$ invertible matrices with real entries. We may similarly define $`\mathrm{𝖦𝖫}(n;)`$ to be the group of all $`n\times n`$ invertible matrices with complex entries. Of course, $`\mathrm{𝖦𝖫}(n;)`$ is contained in $`\mathrm{𝖦𝖫}(n;)`$.
###### Definition 2.1.
Let $`A_n`$ be a sequence of complex matrices. We say that $`A_n`$ converges to a matrix $`A`$ if each entry of $`A_n`$ converges to the corresponding entry of $`A`$, i.e., if $`\left(A_n\right)_{ij}`$ converges to $`A_{ij}`$ for all $`1i,jn`$.
###### Definition 2.2.
A matrix Lie group is any subgroup $`H`$ of $`\mathrm{𝖦𝖫}(n;)`$ with the following property: if $`A_n`$ is any sequence of matrices in $`H`$, and $`A_n`$ converges to some matrix $`A`$, then either $`AH`$, or $`A`$ is not invertible.
The condition on $`H`$ amounts to saying that $`H`$ is a closed subset of $`\mathrm{𝖦𝖫}(n;)`$. (This is not the same as saying that $`H`$ is closed in the space of all matrices.) Thus Definition 2.2 is equivalent to saying that a matrix Lie group is a closed subgroup of $`\mathrm{𝖦𝖫}(n;)`$.
The condition that $`H`$ be a closed subgroup, as opposed to merely a subgroup, should be regarded as a technicality, in that most of the interesting subgroups of $`\mathrm{𝖦𝖫}(n;)`$ have this property. (Almost all of the matrix Lie groups $`H`$ we will consider have the stronger property that if $`A_n`$ is any sequence of matrices in $`H`$, and $`A_n`$ converges to some matrix $`A`$, then $`AH`$.)
There is a topological structure on the set of $`n\times n`$ complex matrices which goes with the above notion of convergence. This topological structure is defined by identifying the space of $`n\times n`$ matrices with $`^{n^2}`$ in the obvious way and using the usual topological structure on $`^{n^2}`$.
#### 2.1.1. Counterexamples
An example of a subgroup of $`\mathrm{𝖦𝖫}(n;)`$ which is not closed (and hence is not a matrix Lie group) is the set of all $`n\times n`$ invertible matrices all of whose entries are real and rational. This is in fact a subgroup of $`\mathrm{𝖦𝖫}(n;)`$, but not a closed subgroup. That is, one can (easily) have a sequence of invertible matrices with rational entries converging to an invertible matrix with some irrational entries. (In fact, every real invertible matrix is the limit of some sequence of invertible matrices with rational entries.)
Another example of a group of matrices which is not a matrix Lie group is the following subgroup of $`\mathrm{𝖦𝖫}(2,)`$. Let $`a`$ be an irrational real number, and let
$$H=\left\{\left(\begin{array}{cc}e^{it}& 0\\ 0& e^{ita}\end{array}\right)\right|t\}$$
Clearly, $`H`$ is a subgroup of $`\mathrm{𝖦𝖫}(2,)`$. Because $`a`$ is irrational, the matrix $`I`$ is not in $`H`$, since to make $`e^{it}`$ equal to $`1`$, we must take $`t`$ to be an odd integer multiple of $`\pi `$, in which case $`ta`$ cannot be an odd integer multiple of $`\pi `$. On the other hand, by taking $`t=(2n+1)\pi `$ for a suitably chosen integer $`n`$, we can make $`ta`$ arbitrarily close to an odd integer multiple of $`\pi `$. (It is left to the reader to verify this.) Hence we can find a sequence of matrices in $`H`$ which converges to $`I`$, and so $`H`$ is not a matrix Lie group. See Exercise 1.
### 2.2. Examples of Matrix Lie Groups
Mastering the subject of Lie groups involves not only learning the general theory, but also familiarizing oneself with examples. In this section, we introduce some of the most important examples of (matrix) Lie groups.
#### 2.2.1. The general linear groups $`\mathrm{𝖦𝖫}(n;)`$ and $`\mathrm{𝖦𝖫}(n;)`$
The general linear groups (over $``$ or $``$) are themselves matrix Lie groups. Of course, $`\mathrm{𝖦𝖫}(n;)`$ is a subgroup of itself. Furthermore, if $`A_n`$ is a sequence of matrices in $`\mathrm{𝖦𝖫}(n;)`$ and $`A_n`$ converges to $`A`$, then by the definition of $`\mathrm{𝖦𝖫}(n;)`$, either $`A`$ is in $`\mathrm{𝖦𝖫}(n;)`$, or $`A`$ is not invertible.
Moreover, $`\mathrm{𝖦𝖫}(n;)`$ is a subgroup of $`\mathrm{𝖦𝖫}(n;)`$, and if $`A_n\mathrm{𝖦𝖫}(n;)`$, and $`A_n`$ converges to $`A`$, then the entries of $`A`$ are real. Thus either $`A`$ is not invertible, or $`A\mathrm{𝖦𝖫}(n;)`$.
#### 2.2.2. The special linear groups SL$`(n;)`$ and SL$`(n;)`$
The special linear group (over $``$ or $``$) is the group of $`n\times n`$ invertible matrices (with real or complex entries) having determinant one. Both of these are subgroups of $`\mathrm{𝖦𝖫}(n;)`$, as noted in Chapter 1. Furthermore, if $`A_n`$ is a sequence of matrices with determinant one, and $`A_n`$ converges to $`A`$, then $`A`$ also has determinant one, because the determinant is a continuous function. Thus SL$`(n;)`$ and SL$`(n;)`$ are matrix Lie groups.
#### 2.2.3. The orthogonal and special orthogonal groups, $`𝖮(n)`$ and SO$`(n)`$
An $`n\times n`$ real matrix $`A`$ is said to be orthogonal if the column vectors that make up $`A`$ are orthonormal, that is, if
$$\underset{i=1}{\overset{n}{}}A_{ij}A_{ik}=\delta _{jk}$$
Equivalently, $`A`$ is orthogonal if it preserves the inner product, namely, if $`x,y=Ax,Ay`$ for all vectors $`x,y`$ in $`^n`$. ( Angled brackets denote the usual inner product on $`^n`$, $`x,y=_ix_iy_i`$.) Still another equivalent definition is that $`A`$ is orthogonal if $`A^{tr}A=I`$, i.e., if $`A^{tr}=A^1`$. ($`A^{tr}`$ is the transpose of $`A`$, $`\left(A^{tr}\right)_{ij}=A_{ji}`$.) See Exercise 2.
Since $`detA^{tr}=detA`$, we see that if $`A`$ is orthogonal, then $`det(A^{tr}A)=\left(detA\right)^2=detI=1`$. Hence $`detA=\pm 1`$, for all orthogonal matrices $`A`$.
This formula tells us, in particular, that every orthogonal matrix must be invertible. But if $`A`$ is an orthogonal matrix, then
$$A^1x,A^1y=A\left(A^1x\right),A\left(A^1x\right)=x,y$$
Thus the inverse of an orthogonal matrix is orthogonal. Furthermore, the product of two orthogonal matrices is orthogonal, since if $`A`$ and $`B`$ both preserve inner products, then so does $`AB`$. Thus the set of orthogonal matrices forms a group.
The set of all $`n\times n`$ real orthogonal matrices is the orthogonal group $`𝖮(n)`$, and is a subgroup of $`\mathrm{𝖦𝖫}(n;)`$. The limit of a sequence of orthogonal matrices is orthogonal, because the relation $`A^{tr}A=I`$ is preserved under limits. Thus $`𝖮(n)`$ is a matrix Lie group.
The set of $`n\times n`$ orthogonal matrices with determinant one is the special orthogonal group $`\mathrm{𝖲𝖮}(n)`$. Clearly this is a subgroup of $`𝖮(n)`$, and hence of $`\mathrm{𝖦𝖫}(n;)`$. Moreover, both orthogonality and the property of having determinant one are preserved under limits, and so $`\mathrm{𝖲𝖮}(n)`$ is a matrix Lie group. Since elements of $`𝖮(n)`$ already have determinant $`\pm 1`$, $`\mathrm{𝖲𝖮}(n)`$ is “half” of $`𝖮(n)`$.
Geometrically, elements of $`𝖮(n)`$ are either rotations, or combinations of rotations and reflections. The elements of $`\mathrm{𝖲𝖮}(n)`$ are just the rotations.
See also Exercise 6.
#### 2.2.4. The unitary and special unitary groups, $`𝖴(n)`$ and $`\mathrm{𝖲𝖴}(n)`$
An $`n\times n`$ complex matrix $`A`$ is said to be unitary if the column vectors of $`A`$ are orthonormal, that is, if
$$\underset{i=1}{\overset{n}{}}\overline{A_{ij}}A_{ik}=\delta _{jk}$$
Equivalently, $`A`$ is unitary if it preserves the inner product, namely, if $`x,y=Ax,Ay`$ for all vectors $`x,y`$ in $`^n`$. (Angled brackets here denote the inner product on $`^n`$, $`x,y=_i\overline{x_i}y_i`$. We will adopt the convention of putting the complex conjugate on the left.) Still another equivalent definition is that $`A`$ is unitary if $`A^{}A=I`$, i.e., if $`A^{}=A^1`$. ($`A^{}`$ is the adjoint of $`A`$, $`\left(A^{}\right)_{ij}=\overline{A_{ji}}`$.) See Exercise 3.
Since $`detA^{}=\overline{detA}`$, we see that if $`A`$ is unitary, then $`det\left(A^{}A\right)=\left|detA\right|^2=detI=1`$. Hence $`\left|detA\right|=1`$, for all unitary matrices $`A`$.
This in particular shows that every unitary matrix is invertible. The same argument as for the orthogonal group shows that the set of unitary matrices forms a group.
The set of all $`n\times n`$ unitary matrices is the unitary group $`𝖴(n)`$, and is a subgroup of $`\mathrm{𝖦𝖫}(n;)`$. The limit of unitary matrices is unitary, so $`𝖴(n)`$ is a matrix Lie group. The set of unitary matrices with determinant one is the special unitary group $`\mathrm{𝖲𝖴}(n)`$. It is easy to check that $`\mathrm{𝖲𝖴}(n)`$ is a matrix Lie group. Note that a unitary matrix can have determinant $`e^{i\theta }`$ for any $`\theta `$, and so $`\mathrm{𝖲𝖴}(n)`$ is a smaller subset of $`𝖴(n)`$ than $`\mathrm{𝖲𝖮}(n)`$ is of $`𝖮(n)`$. (Specifically, $`\mathrm{𝖲𝖮}(n)`$ has the same dimension as $`𝖮(n)`$, whereas $`\mathrm{𝖲𝖴}(n)`$ has dimension one less than that of $`𝖴(n)`$.)
See also Exercise 8.
#### 2.2.5. The complex orthogonal groups, $`𝖮(n;)`$ and $`\mathrm{𝖲𝖮}(n;)`$
Consider the bilinear form $`\left(\right)`$ on $`𝐂^n`$ defined by $`(x,y)=x_iy_i`$. This form is not an inner product, because of the lack of a complex conjugate in the definition. The set of all $`n\times n`$ complex matrices $`A`$ which preserve this form, (i.e., such that $`(Ax,Ay)=(x,y)`$ for all $`x,y𝐂^n`$) is the complex orthogonal group $`𝖮(n;)`$, and is a subgroup of $`\mathrm{𝖦𝖫}(n;)`$. (The proof is the same as for $`𝖮(n)`$.) An $`n\times n`$ complex matrix $`A`$ is in $`𝖮(n;)`$ if and only if $`A^{tr}A=I`$. It is easy to show that $`𝖮(n;)`$ is a matrix Lie group, and that $`detA=\pm 1`$, for all $`A`$ in $`𝖮(n;)`$. Note that $`𝖮(n;)`$ is not the same as the unitary group $`𝖴(n)`$. The group $`\mathrm{𝖲𝖮}(n;)`$ is defined to be the set of all $`A`$ in $`𝖮(n;)`$ with $`detA=1`$. Then $`\mathrm{𝖲𝖮}(n;)`$ is also a matrix Lie group.
#### 2.2.6. The generalized orthogonal and Lorentz groups
Let $`n`$ and $`k`$ be positive integers, and consider $`^{n+k}`$. Define a symmetric bilinear form $`\left[\right]_{n+k}`$ on $`^{n+k}`$ by the formula
(2.1)
$$[x,y]_{n,k}=x_1y_1+\mathrm{}+x_ny_nx_{n+1}y_{n+1}\mathrm{}y_{n+k}x_{n+k}$$
The set of $`(n+k)\times (n+k)`$ real matrices $`A`$ which preserve this form (i.e., such that $`[Ax,Ay]_{n,k}=[x,y]_{n,k}`$ for all $`x,y^{n+k}`$) is the generalized orthogonal group $`𝖮(n;k)`$, and it is a subgroup of $`\mathrm{𝖦𝖫}(n+k;)`$ (Ex. 4). Since $`𝖮(n;k)`$ and $`𝖮(k;n)`$ are essentially the same group, we restrict our attention to the case $`nk`$. It is not hard to check that $`𝖮(n;k)`$ is a matrix Lie group.
If $`A`$ is an $`(n+k)\times (n+k)`$ real matrix, let $`A^{(i)}`$ denote the $`i^{\text{th}}`$ column vector of $`A`$, that is
$$A^{(i)}=\left(\begin{array}{c}A_{1,i}\\ \mathrm{}\\ A_{n+k,i}\end{array}\right)$$
Then $`A`$ is in $`𝖮(n;k)`$ if and only if the following conditions are satisfied:
(2.2)
$$\begin{array}{cccc}[A^{(i)},A^{(j)}]_{n,k}& =& 0& ij\\ [A^{(i)},A^{(i)}]_{n,k}& =& 1& 1in\\ [A^{(i)},A^{(i)}]_{n,k}& =& 1& n+1in+k\end{array}$$
Let $`g`$ denote the $`(n+k)\times (n+k)`$ diagonal matrix with ones in the first $`n`$ diagonal entries, and minus ones in the last $`k`$ diagonal entries. Then $`A`$ is in $`𝖮(n;k)`$ if and only if $`A^{tr}gA=g`$ (Ex. 4). Taking the determinant of this equation gives $`(detA)^2detg=detg`$, or ($`detA)^2=1`$. Thus for any $`A`$ in $`𝖮(n;k)`$, $`detA=\pm 1`$.
The group $`\mathrm{𝖲𝖮}(n;k)`$ is defined to be the set of matrices in $`𝖮(n;k)`$ with $`detA=1`$. This is a subgroup of $`\mathrm{𝖦𝖫}(n+k;)`$, and is a matrix Lie group.
Of particular interest in physics is the Lorentz group $`𝖮(3;1)`$. (Sometimes the phrase Lorentz group is used more generally to refer to the group $`𝖮(n;1)`$ for any $`n1`$.) See also Exercise 7.
#### 2.2.7. The symplectic groups $`\mathrm{𝖲𝗉}(n;)`$, $`\mathrm{𝖲𝗉}(n;)`$, and $`\mathrm{𝖲𝗉}(n)`$
The special and general linear groups, the orthogonal and unitary groups, and the symplectic groups (which will be defined momentarily) make up the classical groups. Of the classical groups, the symplectic groups have the most confusing definition, partly because there are three sets of them ($`\mathrm{𝖲𝗉}(n;)`$, $`\mathrm{𝖲𝗉}(n;)`$, and $`\mathrm{𝖲𝗉}(n)`$), and partly because they involve skew-symmetric bilinear forms rather than the more familiar symmetric bilinear forms. To further confuse matters, the notation for referring to these groups is not consistent from author to author.
Consider the skew-symmetric bilinear form $`B`$ on $`^{2n}`$ defined as follows:
(2.3)
$$B[x,y]=\underset{i=1}{\overset{n}{}}x_iy_{n+i}x_{n+i}y_i$$
The set of all $`2n\times 2n`$ matrices $`A`$ which preserve $`B`$ (i.e., such that $`B[Ax,Ay]=B[x,y]`$ for all $`x,y^{2n}`$) is the real symplectic group $`\mathrm{𝖲𝗉}(n;)`$, and it is a subgroup of $`\mathrm{𝖦𝖫}(2n;)`$. It is not difficult to check that this is a matrix Lie group (Exercise 5). This group arises naturally in the study of classical mechanics. If $`J`$ is the $`2n\times 2n`$ matrix
$$J=\left(\begin{array}{cc}0& I\\ I& 0\end{array}\right)$$
then $`B[x,y]=x,Jy`$, and it is possible to check that a $`2n\times 2n`$ real matrix $`A`$ is in $`\mathrm{𝖲𝗉}(n;)`$ if and only if $`A^{tr}JA=J`$. (See Exercise 5.) Taking the determinant of this identity gives $`\left(detA\right)^2detJ=detJ`$, or $`\left(detA\right)^2=1`$. This shows that $`detA=\pm 1`$, for all $`A\mathrm{𝖲𝗉}(n;)`$. In fact, $`detA=1`$ for all $`A\mathrm{𝖲𝗉}(n;)`$, although this is not obvious.
One can define a bilinear form on $`^n`$ by the same formula (2.3). (This form is bilinear, not Hermitian, and involves no complex conjugates.) The set of $`2n\times 2n`$ complex matrices which preserve this form is the complex symplectic group$`\mathrm{𝖲𝗉}(n;)`$. A $`2n\times 2n`$ complex matrix $`A`$ is in $`\mathrm{𝖲𝗉}(n;)`$ if and only if $`A^{tr}JA=J`$. (Note: this condition involves $`A^{tr}`$, not $`A^{}`$.) This relation shows that $`detA=\pm 1`$, for all $`A\mathrm{𝖲𝗉}(n;)`$. In fact $`detA=1`$, for all $`A\mathrm{𝖲𝗉}(n;)`$.
Finally, we have the compact symplectic group $`\mathrm{𝖲𝗉}(n)`$ defined as
$$\mathrm{𝖲𝗉}(n)=\mathrm{𝖲𝗉}(n;)𝖴(2n).$$
See also Exercise 9. For more information and a proof of the fact that $`detA=1`$, for all $`A\mathrm{𝖲𝗉}(n;)`$, see Miller, Sect. 9.4. What we call $`\mathrm{𝖲𝗉}(n;)`$ Miller calls $`\mathrm{𝖲𝗉}(n)`$, and what we call $`\mathrm{𝖲𝗉}(n)`$, Miller calls $`\mathrm{𝖴𝖲𝗉}(n)`$.
#### 2.2.8. The Heisenberg group $`H`$
The set of all $`3\times 3`$ real matrices $`A`$ of the form
(2.4)
$$A=\left(\begin{array}{ccc}1& a& b\\ 0& 1& c\\ 0& 0& 1\end{array}\right)$$
where $`a`$, $`b`$, and $`c`$ are arbitrary real numbers, is the Heisenberg group. It is easy to check that the product of two matrices of the form (2.4) is again of that form, and clearly the identity matrix is of the form (2.4). Furthermore, direct computation shows that if $`A`$ is as in (2.4), then
$$A^1=\left(\begin{array}{ccc}1& a& acb\\ 0& 1& c\\ 0& 0& 1\end{array}\right)$$
Thus $`H`$ is a subgroup of $`\mathrm{𝖦𝖫}(3;)`$. Clearly the limit of matrices of the form (2.4) is again of that form, and so $`H`$ is a matrix Lie group.
It is not evident at the moment why this group should be called the Heisenberg group. We shall see later that the Lie algebra of $`H`$ gives a realization of the Heisenberg commutation relations of quantum mechanics. (See especially Chapter 5, Exercise 10.)
See also Exercise 10.
#### 2.2.9. The groups $`^{}`$, $`^{}`$, $`S^1`$, $``$, and $`^n`$
Several important groups which are not naturally groups of matrices can (and will in these notes) be thought of as such.
The group $`^{}`$ of non-zero real numbers under multiplication is isomorphic to $`\mathrm{𝖦𝖫}(1,)`$. Thus we will regard $`^{}`$ as a matrix Lie group. Similarly, the group $`^{}`$ of non-zero complex numbers under multiplication is isomorphic to $`\mathrm{𝖦𝖫}(1;)`$, and the group $`S^1`$ of complex numbers with absolute value one is isomorphic to $`𝖴(1)`$.
The group $``$ under addition is isomorphic to $`\mathrm{𝖦𝖫}(1;)^+`$ ($`1\times 1`$ real matrices with positive determinant) via the map $`x\left[e^x\right]`$. The group $`^n`$ (with vector addition) is isomorphic to the group of diagonal real matrices with positive diagonal entries, via the map
$$(x_1,\mathrm{},x_n)\left(\begin{array}{ccc}e^{x_1}& & 0\\ & \mathrm{}& \\ 0& & e^{x_n}\end{array}\right)\text{.}$$
#### 2.2.10. The Euclidean and Poincaré groups
The Euclidean group $`𝖤(n)`$ is by definition the group of all one-to-one, onto, distance-preserving maps of $`^n`$ to itself, that is, maps $`f:^n^n`$ such that $`d(f\left(x\right),f\left(y\right))=d(x,y)`$ for all $`x,y^n.`$ Here $`d`$ is the usual distance on $`^n,`$ $`d(x,y)=\left|xy\right|.`$ Note that we don’t assume anything about the structure of $`f`$ besides the above properties. In particular, $`f`$ need not be linear. The orthogonal group $`𝖮(n)`$ is a subgroup of $`𝖤(n)`$, and is the group of all linear distance-preserving maps of $`^n`$ to itself. The set of translations of $`^n`$ (i.e., the set of maps of the form $`T_x(y)=x+y`$) is also a subgroup of $`𝖤(n)`$.
###### Proposition 2.3.
Every element $`T`$ of $`𝖤(n)`$ can be written uniquely as an orthogonal linear transformation followed by a translation, that is, in the form
$$T=T_xR$$
with $`x^n`$, and $`R𝖮(n)`$.
We will not prove this here. The key step is to prove that every one-to-one, onto, distance-preserving map of $`^n`$ to itself which fixes the origin must be linear.
Following Miller, we will write an element $`T=T_xR`$ of $`𝖤(n)`$ as a pair $`\{x,R\}`$. Note that for $`y^n`$,
$$\{x,R\}y=Ry+x$$
and that
$$\{x_1,R_1\}\{x_2,R_2\}y=R_1(R_2y+x_2)+x_1=R_1R_2y+(x_1+R_1x_2)$$
Thus the product operation for $`𝖤(n)`$ is the following:
(2.5)
$$\{x_1,R_1\}\{x_2,R_2\}=\{x_1+R_1x_2,R_1R_2\}$$
The inverse of an element of $`𝖤(n)`$ is given by
$$\{x,R\}^1=\{R^1x,R^1\}$$
Now, as already noted, $`𝖤(n)`$ is not a subgroup of $`\mathrm{𝖦𝖫}(n;)`$, since translations are not linear maps. However, $`𝖤(n)`$ is isomorphic to a subgroup of $`\mathrm{𝖦𝖫}(n+1;)`$, via the map which associates to $`\{x,R\}𝖤(n)`$ the following matrix
(2.6)
$$\left(\begin{array}{cccc}& & & x_1\\ & R& & \mathrm{}\\ & & & x_n\\ 0& \mathrm{}& 0& 1\end{array}\right)$$
This map is clearly one-to-one, and it is a simple computation to show that it is a homomorphism. Thus $`𝖤(n)`$ is isomorphic to the group of all matrices of the form (2.6) (with $`R𝖮(n)`$). The limit of things of the form (2.6) is again of that form, and so we have expressed the Euclidean group $`𝖤(n)`$ as a matrix Lie group.
We similarly define the Poincaré group $`𝖯(n;1)`$ to be the group of all transformations of $`^{n+1}`$ of the form
$$T=T_xA$$
with $`x^{n+1}`$, $`A𝖮(n;1)`$. This is the group of affine transformations of $`^{n+1}`$ which preserve the Lorentz “distance” $`d_L(x,y)=(x_1y_1)^2+\mathrm{}+(x_ny_n)^2(x_{n+1}y_{n+1})^2`$. (An affine transformation is one of the form $`xAx+b,`$ where $`A`$ is a linear transformation and $`b`$ is constant.) The group product is the obvious analog of the product (2.5) for the Euclidean group.
The Poincaré group $`𝖯(n;1)`$ is isomorphic to the group of $`(n+2)\times (n+2)`$ matrices of the form
(2.7)
$$\left(\begin{array}{cccc}& & & x_1\\ & A& & \mathrm{}\\ & & & x_{n+1}\\ 0& \mathrm{}& 0& 1\end{array}\right)$$
with $`A𝖮(n;1)`$. The set of matrices of the form (2.7) is a matrix Lie group.
### 2.3. Compactness
###### Definition 2.4.
A matrix Lie group $`G`$ is said to be compact if the following two conditions are satisfied:
1. If $`A_n`$ is any sequence of matrices in $`G`$, and $`A_n`$ converges to a matrix $`A`$, then $`A`$ is in $`G`$.
2. There exists a constant $`C`$ such that for all $`AG`$, $`\left|A_{ij}\right|C`$ for all $`1i,jn`$.
This is not the usual topological definition of compactness. However, the set of all $`n\times n`$ complex matrices can be thought of as $`^{n^2}`$. The above definition says that $`G`$ is compact if it is a closed, bounded subset of $`^{n^2}`$. It is a standard theorem from elementary analysis that a subset of $`^m`$ is compact (in the usual sense that every open cover has a finite subcover) if and only if it is closed and bounded.
All of our examples of matrix Lie groups except $`\mathrm{𝖦𝖫}(n;)`$ and $`\mathrm{𝖦𝖫}(n;)`$ have property (1). Thus it is the boundedness condition (2) that is most important.
The property of compactness has very important implications. For example, if $`G`$ is compact, then every irreducible unitary representation of $`G`$ is finite-dimensional.
#### 2.3.1. Examples of compact groups
The groups $`𝖮(n)`$ and $`\mathrm{𝖲𝖮}(n)`$ are compact. Property (1) is satisfied because the limit of orthogonal matrices is orthogonal and the limit of matrices with determinant one has determinant one. Property (2) is satisfied because if $`A`$ is orthogonal, then the column vectors of $`A`$ have norm one, and hence $`\left|A_{ij}\right|1`$, for all $`1i,jn`$. A similar argument shows that $`𝖴(n)`$, $`\mathrm{𝖲𝖴}(n)`$, and $`\mathrm{𝖲𝗉}(n)`$ are compact. (This includes the unit circle, $`S^1𝖴(1)`$.)
#### 2.3.2. Examples of non-compact groups
All of the other examples given of matrix Lie groups are non-compact. $`\mathrm{𝖦𝖫}(n;)`$ and $`\mathrm{𝖦𝖫}(n;)`$ violate property (1), since a limit of invertible matrices may be non-invertible. $`\mathrm{𝖲𝖫}(n;)`$ and $`\mathrm{𝖲𝖫}(n;)`$ violate (2), except in the trivial case $`n=1`$, since
$$A_n=\left(\begin{array}{ccccc}n& & & & \\ & \frac{1}{n}& & & \\ & & 1& & \\ & & & \mathrm{}& \\ & & & & 1\end{array}\right)$$
has determinant one, no matter how big $`n`$ is.
The following groups also violate (2), and hence are non-compact: $`𝖮(n;)`$ and $`\mathrm{𝖲𝖮}(n;)`$; $`𝖮(n;k)`$ and $`\mathrm{𝖲𝖮}(n;k)`$ ($`n1`$, $`k1`$); the Heisenberg group $`H`$; $`\mathrm{𝖲𝗉}(n;)`$ and $`\mathrm{𝖲𝗉}(n;)`$; $`𝖤(n)`$ and $`𝖯(n;1)`$; $``$ and $`^n`$; $`^{}`$ and $`^{}`$. It is left to the reader to provide examples to show that this is the case.
### 2.4. Connectedness
###### Definition 2.5.
A matrix Lie group $`G`$ is said to be connected if given any two matrices $`A`$ and $`B`$ in $`G`$, there exists a continuous path $`A(t)`$, $`atb`$, lying in $`G`$ with $`A(a)=A`$, and $`A(b)=B`$.
This property is what is called path-connected in topology, which is not (in general) the same as connected. However, it is a fact (not particularly obvious at the moment) that a matrix Lie group is connected if and only if it is path-connected. So in a slight abuse of terminology we shall continue to refer to the above property as connectedness. (See Section 2.7.)
A matrix Lie group $`G`$ which is not connected can be decomposed (uniquely) as a union of several pieces, called components, such that two elements of the same component can be joined by a continuous path, but two elements of different components cannot.
###### Proposition 2.6.
If $`G`$ is a matrix Lie group, then the component of $`G`$ containing the identity is a subgroup of $`G`$.
###### Proof.
Saying that $`A`$ and $`B`$ are both in the component containing the identity means that there exist continuous paths $`A(t)`$ and $`B(t)`$ with $`A(0)=B(0)=I`$, $`A(1)=A`$, and $`B(1)=B`$. But then $`A(t)B(t)`$ is a continuous path starting at $`I`$ and ending at $`AB`$. Thus the product of two elements of the identity component is again in the identity component. Furthermore, $`A(t)^1`$ is a continuous path starting at $`I`$ and ending at $`A^1`$, and so the inverse of any element of the identity component is again in the identity component. Thus the identity component is a subgroup. ∎
###### Proposition 2.7.
The group $`\mathrm{𝖦𝖫}(n;)`$ is connected for all $`n1`$.
###### Proof.
Consider first the case $`n=1`$. A $`1\times 1`$ invertible complex matrix $`A`$ is of the form $`A=\left[\lambda \right]`$ with $`\lambda ^{}`$, the set of non-zero complex numbers. But given any two non-zero complex numbers, we can easily find a continuous path which connects them and does not pass through zero.
For the case $`n1`$, we use the Jordan canonical form. Every $`n\times n`$ complex matrix $`A`$ can be written as
$$A=CBC^1$$
where $`B`$ is the Jordan canonical form. The only property of $`B`$ we will need is that $`B`$ is upper-triangular:
$$B=\left(\begin{array}{ccc}\lambda _1& & \\ & \mathrm{}& \\ 0& & \lambda _n\end{array}\right)$$
If $`A`$ is invertible, then all the $`\lambda _i`$’s must be non-zero, since $`detA=detB=\lambda _1\mathrm{}\lambda _n`$.
Let $`B(t)`$ be obtained by multiplying the part of $`B`$ above the diagonal by $`(1t)`$, for $`0t1`$, and let $`A(t)=CB(t)C^1`$. Then $`A(t)`$ is a continuous path which starts at $`A`$ and ends at $`CDC^1`$, where $`D`$ is the diagonal matrix
$$D=\left(\begin{array}{ccc}\lambda _1& & 0\\ & \mathrm{}& \\ 0& & \lambda _n\end{array}\right)$$
This path lies in $`\mathrm{𝖦𝖫}(n;)`$ since $`detA(t)=\lambda _1\mathrm{}\lambda _n`$ for all $`t`$.
But now, as in the case $`n=1`$, we can define $`\lambda _i(t)`$ which connects each $`\lambda _i`$ to 1 in $`𝐂^{}`$, as $`t`$ goes from 1 to 2. Then we can define
$$A(t)=C\left(\begin{array}{ccc}\lambda _1(t)& & 0\\ & \mathrm{}& \\ 0& & \lambda _n(t)\end{array}\right)C^1$$
This is a continuous path which starts at $`CDC^1`$ when $`t=1`$, and ends at $`I`$ ($`=CIC^1`$) when $`t=2`$. Since the $`\lambda _i(t)`$’s are always non-zero, $`A(t)`$ lies in $`\mathrm{𝖦𝖫}(n;)`$.
We see, then, that every matrix $`A`$ in $`\mathrm{𝖦𝖫}(n;)`$ can be connected to the identity by a continuous path lying in $`\mathrm{𝖦𝖫}(n;)`$. Thus if $`A`$ and $`B`$ are two matrices in $`\mathrm{𝖦𝖫}(n;)`$, they can be connected by connecting each of them to the identity. ∎
###### Proposition 2.8.
The group $`\mathrm{𝖲𝖫}(n;)`$ is connected for all $`n1`$.
###### Proof.
The proof is almost the same as for $`\mathrm{𝖦𝖫}(n;)`$, except that we must be careful to preserve the condition $`detA=1`$. Let $`A`$ be an arbitrary element of $`\mathrm{𝖲𝖫}(n;)`$. The case $`n=1`$ is trivial, so we assume $`n2`$. We can define $`A(t)`$ as above for $`0t1`$, with $`A(0)=A`$, and $`A(1)=CDC^1`$, since $`detA(t)=detA=1`$. Now define $`\lambda _i(t)`$ as before for $`1in1`$, and define $`\lambda _n(t)`$ to be $`\left[\lambda _1(t)\mathrm{}\lambda _{n1}(t)\right]^1`$. (Note that since $`\lambda _1\mathrm{}\lambda _n=1`$, $`\lambda _n(0)=\lambda _n`$.) This allows us to connect $`A`$ to the identity while staying within $`\mathrm{𝖲𝖫}(n;)`$. ∎
###### Proposition 2.9.
The groups $`𝖴(n)`$ and $`\mathrm{𝖲𝖴}(n)`$ are connected, for all $`n1`$.
###### Proof.
By a standard result of linear algebra, every unitary matrix has an orthonormal basis of eigenvectors, with eigenvalues of the form $`e^{i\theta }`$. It follows that every unitary matrix $`U`$ can be written as
(2.8)
$$U=U_1\left(\begin{array}{ccc}e^{i\theta _1}& & 0\\ & \mathrm{}& \\ 0& & e^{i\theta _n}\end{array}\right)U_1^1$$
with $`U_1`$ unitary and $`\theta _i`$. Conversely, as is easily checked, every matrix of the form (2.8) is unitary. Now define
$$U(t)=U_1\left(\begin{array}{ccc}e^{i(1t)\theta _1}& & 0\\ & \mathrm{}& \\ 0& & e^{i(1t)\theta _n}\end{array}\right)U_1^1$$
As $`t`$ ranges from 0 to 1, this defines a continuous path in $`𝖴(n)`$ joining $`U`$ to $`I`$. This shows that $`𝖴(n)`$ is connected.
A slight modification of this argument, as in the proof of Proposition 2.8, shows that $`\mathrm{𝖲𝖴}(n)`$ is connected. ∎
###### Proposition 2.10.
The group $`\mathrm{𝖦𝖫}(n;)`$ is not connected, but has two components. These are $`\mathrm{𝖦𝖫}(n;)^+`$, the set of $`n\times n`$ real matrices with positive determinant, and $`\mathrm{𝖦𝖫}(n;)^{}`$, the set of $`n\times n`$ real matrices with negative determinant.
###### Proof.
$`\mathrm{𝖦𝖫}(n;)`$ cannot be connected, for if $`detA>0`$ and $`detB<0`$, then any continuous path connecting $`A`$ to $`B`$ would have to include a matrix with determinant zero, and hence pass outside of $`\mathrm{𝖦𝖫}(n;)`$.
The proof that $`\mathrm{𝖦𝖫}(n;)^+`$ is connected is given in Exercise 14. Once $`\mathrm{𝖦𝖫}(n;)^+`$ is known to be connected, it is not difficult to see that $`\mathrm{𝖦𝖫}(n;)^{}`$ is also connected. For let $`C`$ be any matrix with negative determinant, and take $`A,B`$ in $`\mathrm{𝖦𝖫}(n;)^{}`$. Then $`C^1A`$ and $`C^1B`$ are in $`\mathrm{𝖦𝖫}(n;)^+`$, and can be joined by a continuous path $`D(t)`$ in $`\mathrm{𝖦𝖫}(n;)^+`$. But then $`CD(t)`$ is a continuous path joining $`A`$ and $`B`$ in $`\mathrm{𝖦𝖫}(n;)^{}`$. ∎
The following table lists some matrix Lie groups, indicates whether or not the group is connected, and gives the number of components.
$$\begin{array}{ccc}\text{Group}& \text{Connected?}& \text{Components}\\ \mathrm{𝖦𝖫}(n;)& \text{yes}& 1\\ \mathrm{𝖲𝖫}(n;)& \text{yes}& 1\\ \mathrm{𝖦𝖫}(n;)& \text{no}& 2\\ \mathrm{𝖲𝖫}(n;)& \text{yes}& 1\\ 𝖮(n)& \text{no}& 2\\ \mathrm{𝖲𝖮}(n)& \text{yes}& 1\\ 𝖴(n)& \text{yes}& 1\\ \mathrm{𝖲𝖴}(n)& \text{yes}& 1\\ 𝖮(n;1)& \text{no}& 4\\ \mathrm{𝖲𝖮}(n;1)& \text{no}& 2\\ \text{Heisenberg}& \text{yes}& 1\\ 𝖤\left(n\right)& \text{no}& 2\\ 𝖯(n;1)& \text{no}& 4\end{array}$$
Proofs of some of these results are given in Exercises 7, 11, 13, and 14. (The connectedness of the Heisenberg group is immediate.)
### 2.5. Simple-connectedness
###### Definition 2.11.
A connected matrix Lie group $`G`$ is said to be simply connected if every loop in $`G`$ can be shrunk continuously to a point in $`G`$.
More precisely, $`G`$ is simply connected if given any continuous path $`A(t)`$, $`0t1`$, lying in $`G`$ with $`A(0)=A(1)`$, there exists a continuous function $`A(s,t)`$, $`0s,t1`$, taking values in $`G`$ with the following properties: 1)$`A(s,0)=A(s,1)`$ for all $`s`$, 2) $`A(0,t)=A(t)`$, and 3) $`A(1,t)=A(1,0)`$ for all $`t`$.
You should think of $`A(t)`$ as a loop, and $`A(s,t)`$ as a parameterized family of loops which shrinks $`A(t)`$ to a point. Condition 1) says that for each value of the parameter $`s`$, we have a loop; condition 2) says that when $`s=0`$ the loop is the specified loop $`A(t)`$; and condition 3) says that when $`s=1`$ our loop is a point.
It is customary to speak of simple-connectedness only for connected matrix Lie groups, even though the definition makes sense for disconnected groups.
###### Proposition 2.12.
The group $`\mathrm{𝖲𝖴}(2)`$ is simply connected.
###### Proof.
Exercise 8 shows that $`\mathrm{𝖲𝖴}(2)`$ may be thought of (topologically) as the three-dimensional sphere $`S^3`$ sitting inside $`^4`$. It is well-known that $`S^3`$ is simply connected. ∎
The condition of simple-connectedness is extremely important. One of our most important theorems will be that if $`G`$ is simply connected, then there is a natural one-to-one correspondence between the representations of $`G`$ and the representations of its Lie algebra.
Without proof, we give the following table.
$$\begin{array}{cc}\text{Group}& \text{Simply connected?}\\ \mathrm{𝖦𝖫}(n;)& \text{no}\\ \mathrm{𝖲𝖫}(n;)& \text{yes}\\ \mathrm{𝖦𝖫}(n;)& \text{no}\\ \mathrm{𝖲𝖫}(n;)& \text{no}\\ \mathrm{𝖲𝖮}(n)& \text{no}\\ 𝖴(n)& \text{no}\\ \mathrm{𝖲𝖴}(n)& \text{yes}\\ \mathrm{𝖲𝖮}(1;1)& \text{yes}\\ \mathrm{𝖲𝖮}(n;1)\text{ (}n2\text{)}& \text{no}\\ \text{Heisenberg}& \text{yes}\end{array}$$
### 2.6. Homomorphisms and Isomorphisms
###### Definition 2.13.
Let $`G`$ and $`H`$ be matrix Lie groups. A map $`\varphi `$ from $`G`$ to $`H`$ is called a Lie group homomorphism if 1) $`\varphi `$ is a group homomorphism and 2) $`\varphi `$ is continuous. If in addition, $`\varphi `$ is one-to-one and onto, and the inverse map $`\varphi ^1`$ is continuous, then $`\varphi `$ is called a Lie group isomorphism.
The condition that $`\varphi `$ be continuous should be regarded as a technicality, in that it is very difficult to give an example of a group homomorphism between two matrix Lie groups which is not continuous. In fact, if $`G=`$ and $`H=^{}`$, then any group homomorphism from $`G`$ to $`H`$ which is even measurable (a very weak condition) must be continuous. (See W. Rudin, Real and Complex Analysis, Chap. 9, Ex. 17.)
If $`G`$ and $`H`$ are matrix Lie groups, and there exists a Lie group isomorphism from $`G`$ to $`H`$, then $`G`$ and $`H`$ are said to be isomorphic, and we write $`GH`$. Two matrix Lie groups which are isomorphic should be thought of as being essentially the same group. (Note that by definition, the inverse of Lie group isomorphism is continuous, and so also a Lie group isomorphism.)
#### 2.6.1. Example: $`\mathrm{𝖲𝖴}(2)`$ and $`\mathrm{𝖲𝖮}(3)`$
A very important topic for us will be the relationship between the groups $`\mathrm{𝖲𝖴}(2)`$ and $`\mathrm{𝖲𝖮}(3)`$. This example is designed to show that $`\mathrm{𝖲𝖴}(2)`$ and $`\mathrm{𝖲𝖮}(3)`$ are almost (but not quite!) isomorphic. Specifically, there exists a Lie group homomorphism $`\varphi `$ which maps $`\mathrm{𝖲𝖴}(2)`$ onto $`\mathrm{𝖲𝖮}(3)`$, and which is two-to-one. (See Miller 7.1 and Bröcker, Chap. I, 6.18.)
Consider the space $`V`$ of all $`2\times 2`$ complex matrices which are self-adjoint and have trace zero. This is a three-dimensional real vector space with the following basis
$$\begin{array}{ccc}A_1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right);& A_2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right);& A_3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\end{array}$$
We may define an inner product on $`V`$ by the formula
$$A,B=\frac{1}{2}\mathrm{trace}(AB)$$
(Exercise: check that this is an inner product.)
Direct computation shows that $`\{A_1,A_2,A_3\}`$ is an orthonormal basis for $`V`$. Having chosen an orthonormal basis for $`V`$, we can identify $`V`$ with $`^3`$.
Now, if $`U`$ is an element of $`\mathrm{𝖲𝖴}(2)`$, and $`A`$ is an element of $`V`$, then it is easy to see that $`UAU^1`$ is in $`V`$. Thus for each $`U\mathrm{𝖲𝖴}(2)`$, we can define a linear map $`\varphi _U`$ of $`V`$ to itself by the formula
$$\varphi _U(A)=UAU^1$$
(This definition would work for $`U𝖴(2)`$, but we choose to restrict our attention to $`\mathrm{𝖲𝖴}(2)`$.) Moreover, given $`U\mathrm{𝖲𝖴}(2)`$, and $`A,BV`$, note that
$$\varphi _U(A),\varphi _U(B)=\frac{1}{2}\mathrm{trace}(UAU^1UBU^1)=\frac{1}{2}\mathrm{trace}(AB)=A,B$$
Thus $`\varphi _U`$ is an orthogonal transformation of $`V^3`$, which we can think of as an element of $`𝖮(3)`$.
We see, then, that the map $`U\varphi _U`$ is a map of $`\mathrm{𝖲𝖴}(2)`$ into $`𝖮(3)`$. It is very easy to check that this map is a homomorphism (i.e., $`\varphi _{U_1U_2}=\varphi _{U_1}\varphi _{U_2}`$), and that it is continuous. Thus $`U\varphi _U`$ is a Lie group homomorphism of $`\mathrm{𝖲𝖴}(2)`$ into $`𝖮(3)`$.
Recall that every element of $`𝖮(3)`$ has determinant $`\pm 1`$. Since $`\mathrm{𝖲𝖴}(2)`$ is connected (Exercise 8), and the map $`U\varphi _U`$ is continuous, $`\varphi _U`$ must actually map into $`\mathrm{𝖲𝖮}(3)`$. Thus $`U\varphi _U`$ is a Lie group homomorphism of $`\mathrm{𝖲𝖴}(2)`$ into $`\mathrm{𝖲𝖮}(3)`$.
The map $`U\varphi _U`$ is not one-to-one, since for any $`U\mathrm{𝖲𝖴}(2)`$, $`\varphi _U=\varphi _U`$. (Observe that if $`U`$ is in $`\mathrm{𝖲𝖴}(2)`$, then so is $`U`$.) It is possible to show that $`\varphi _U`$ is a two-to-one map of $`\mathrm{𝖲𝖴}(2)`$ onto $`\mathrm{𝖲𝖮}(3)`$. (See Miller.)
### 2.7. Lie Groups
A Lie group is something which is simultaneously a group and a differentiable manifold (see Definition 2.14). As the terminology suggests, every matrix Lie group is a Lie group, although this requires proof (Theorem 2.15). I have decided to restrict attention to matrix Lie groups, except in emergencies, for three reasons. First, this makes the course accessible to students who are not familiar with the theory of differentiable manifolds. Second, this makes the definition of the Lie algebra and of the exponential mapping far more comprehensible. Third, all of the important examples of Lie groups are (or can easily be represented as) matrix Lie groups.
Alas, there is a price to pay for this simplification. Certain important topics (notably, the universal cover) are considerably complicated by restricting to the matrix case. Nevertheless, I feel that the advantages outweigh the disadvantages in an introductory course such as this.
###### Definition 2.14.
A Lie group is a differentiable manifold $`G`$ which is also a group, and such that the group product
$$G\times GG$$
and the inverse map $`gg^1`$ are differentiable.
For the reader who is not familiar with the notion of a differentiable manifold, here is a brief recap. (I will consider only manifolds embedded in some $`^n`$, which is a harmless assumption.) A subset $`M`$ of $`𝐑^n`$ is called a $`k`$-dimensional differentiable manifold if given any $`m_0M`$, there exists a smooth (non-linear) coordinate system $`(x^1,\mathrm{}x^n)`$ defined in a neighborhood $`U`$ of $`m_0`$ such that
$$MU=\left\{mU\right|x^{k+1}(m)=c_1,\mathrm{},x^n(m)=c_{nk}\}$$
This says that locally, after a suitable change of variables, $`M`$ looks like the $`k`$-dimensional hyperplane in $`^n`$ obtained by setting all but the first $`k`$ coordinates equal to constants.
For example, $`S^1^2`$ is a one-dimensional differentiable manifold because in the usual polar coordinates $`(\theta ,r)`$, $`S^1`$ is the set $`r=1`$. Of course, polar coordinates are not globally defined, because $`\theta `$ is undefined at the origin, and because $`\theta `$ is not “single-valued.” But given any point $`m_0`$ in $`S^1`$, we can define polar coordinates in a neighborhood $`U`$ of $`m_0`$, and then $`S^1U`$ will be the set $`r=1`$.
Note that while we assume that our differentiable manifolds are embedded in some $`^n`$ (a harmless assumption), we are not saying that a Lie group has to be embedded in $`^{n^2}`$, or that the group operation has to have anything to do with matrix multiplication. A Lie group is simply a subset $`G`$ of some $`^n`$ which is a differentiable manifold, together with any map from $`G\times G`$ into $`G`$ which makes $`G`$ into a group (and such that the group operations are smooth). It is remarkable that almost (but not quite!) every Lie group is isomorphic to a matrix Lie group.
Note also that it is far from obvious that a matrix Lie group must be a Lie group, since our definition of a matrix Lie group $`G`$ does not say anything about $`G`$ being a manifold. It is not too difficult to verify that all of our examples of matrix Lie groups are Lie groups, but in fact we have the following result which makes such verifications unnecessary:
###### Theorem 2.15.
Every matrix Lie group is a Lie group.
Although I will not prove this result, I want to discuss what would be involved. Let us consider first the group $`\mathrm{𝖦𝖫}(n;)`$. The space of all $`n\times n`$ real matrices can be thought of as $`^{n^2}`$. Since $`\mathrm{𝖦𝖫}(n;)`$ is the set of all matrices $`A`$ with $`detA0`$, $`\mathrm{𝖦𝖫}(n;)`$ is an open subset of $`^{n^2}`$. (That is, given an invertible matrix $`A`$, there is a neighborhood $`U`$ of $`A`$ such that every matrix $`BU`$ is also invertible.) Thus $`\mathrm{𝖦𝖫}(n;)`$ is an $`n^2`$-dimensional smooth manifold. Furthermore, the matrix product $`AB`$ is clearly a smooth (even polynomial) function of the entries of $`A`$ and $`B`$, and (in light of Kramer’s rule) $`A^1`$ is a smooth function of the entries of $`A`$. Thus $`\mathrm{𝖦𝖫}(n;)`$ is a Lie group.
Similarly, if we think of the space of $`n\times n`$ complex matrices as $`^{n^2}^{2n^2}`$, then the same argument shows that $`\mathrm{𝖦𝖫}(n;)`$ is a Lie group.
Thus, to prove that every matrix Lie group is a Lie group, it suffices to show that a closed subgroup of a Lie group is a Lie group. This is proved in Bröcker and tom Dieck, Chapter I, Theorem 3.11. The proof is not too difficult, but it requires the exponential mapping, which we have not yet introduced. (See Chapter 3.)
It is customary to call a map $`\varphi `$ between two Lie groups a Lie group homomorphism if $`\varphi `$ is a group homomorphism and $`\varphi `$ is smooth, whereas we have (in Definition 2.13) required only that $`\varphi `$ be continuous. However, the following Proposition shows that our definition is equivalent to the more standard one.
###### Proposition 2.16.
Let $`G`$ and $`H`$ be Lie groups, and $`\varphi `$ a group homomorphism from $`G`$ to $`H`$. Then if $`\varphi `$ is continuous it is also smooth.
Thus group homomorphisms from $`G`$ to $`H`$ come in only two varieties: the very bad ones (discontinuous), and the very good ones (smooth). There simply aren’t any intermediate ones. (See, for example, Exercise 16.) For proof, see Bröcker and tom Dieck, Chapter I, Proposition 3.12.
In light of Theorem 2.15, every matrix Lie group is a (smooth) manifold. As such, a matrix Lie group is automatically locally path connected. It follows that a matrix Lie group is path connected if and only if it is connected. (See Remarks following Definition 2.5.)
### 2.8. Exercises
1. Let $`a`$ be an irrational real number. Show that the set of numbers of the form $`e^{2\pi ina}`$, $`n`$, is dense in $`S^1`$. Now let $`G`$ be the following subgroup of $`\mathrm{𝖦𝖫}(2;)`$:
$$G=\left\{\left(\begin{array}{cc}e^{it}& 0\\ 0& e^{iat}\end{array}\right)\right|t\}$$
Show that
$$\overline{G}=\left\{\left(\begin{array}{cc}e^{it}& 0\\ 0& e^{is}\end{array}\right)\right|t,s\},$$
where $`\overline{G}`$ denotes the closure of the set $`G`$ inside the space of $`2\times 2`$ matrices.
Note: The group $`\overline{G}`$ can be thought of as the torus $`S^1\times S^1`$, which in turn can be thought of as $`[0,2\pi ]\times [0,2\pi ]`$, with the ends of the intervals identified. The set $`G[0,2\pi ]\times [0,2\pi ]`$ is called an irrational line. Draw a picture of this set and you should see why $`G`$ is dense in $`[0,2\pi ]\times [0,2\pi ]`$.
2. Orthogonal groups. Let $``$ denote the standard inner product on $`^n`$, $`x,y=_ix_iy_i`$. Show that a matrix $`A`$ preserves inner products if and only if the column vectors of $`A`$ are orthonormal.
Show that for any $`n\times n`$ real matrix $`B`$,
$$Bx,y=x,B^{tr}y$$
where $`\left(B^{tr}\right)_{ij}=B_{ji}`$. Using this fact, show that a matrix $`A`$ preserves inner products if and only if $`A^{tr}A=I`$.
Note: a similar analysis applies to the complex orthogonal groups $`𝖮(n;)`$ and $`\mathrm{𝖲𝖮}(n;)`$.
3. Unitary groups. Let $``$ denote the standard inner product on $`^n`$, $`x,y=_i\overline{x_i}y_i`$. Following Exercise 2, show that $`A^{}A=I`$ if and only if $`Ax,Ay=x,y`$ for all $`x,y^n`$. ($`\left(A^{}\right)_{ij}=\overline{A_{ji}}`$.)
4. Generalized orthogonal groups. Let $`[x,y]_{n,k}`$ be the symmetric bilinear form on $`^{n+k}`$ defined in (2.1). Let $`g`$ be the $`(n+k)\times (n+k)`$ diagonal matrix with first $`n`$ diagonal entries equal to one, and last $`k`$ diagonal entries equal to minus one:
$$g=\left(\begin{array}{cc}I_n& 0\\ 0& I_k\end{array}\right)$$
Show that for all $`x,y^{n+k}`$,
$$[x,y]_{n,k}=x,gy$$
Show that a $`(n+k)\times (n+k)`$ real matrix $`A`$ is in $`𝖮(n;k)`$ if and only if $`A^{tr}gA=g`$. Show that $`𝖮(n;k)`$ and $`\mathrm{𝖲𝖮}(n;k)`$ are subgroups of $`\mathrm{𝖦𝖫}(n+k;)`$, and are matrix Lie groups.
5. Symplectic groups. Let $`B[x,y]`$ be the skew-symmetric bilinear form on $`^{2n}`$ given by $`B[x,y]=_{i=1}^nx_iy_{n+i}x_{n+i}y_i`$. Let $`J`$ be the $`2n\times 2n`$ matrix
$$J=\left(\begin{array}{cc}0& I\\ I& 0\end{array}\right)$$
Show that for all $`x,y^{2n}`$
$$B[x,y]=x,Jy$$
Show that a $`2n\times 2n`$ matrix $`A`$ is in $`\mathrm{𝖲𝗉}(n;)`$ if and only if $`A^{tr}JA=J`$. Show that $`\mathrm{𝖲𝗉}(n;)`$ is a subgroup of $`\mathrm{𝖦𝖫}(2n;)`$, and a matrix Lie group.
Note: a similar analysis applies to $`\mathrm{𝖲𝗉}(n;)`$.
6. The groups$`𝖮(2)`$ and $`\mathrm{𝖲𝖮}(2)`$. Show that the matrix
$$A=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)$$
is in $`\mathrm{𝖲𝖮}(2)`$, and that
$$\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)\left(\begin{array}{cc}\mathrm{cos}\varphi & \mathrm{sin}\varphi \\ \mathrm{sin}\varphi & \mathrm{cos}\varphi \end{array}\right)=\left(\begin{array}{cc}\mathrm{cos}(\theta +\varphi )& \mathrm{sin}(\theta +\varphi )\\ \mathrm{sin}(\theta +\varphi )& \mathrm{cos}(\theta +\varphi )\end{array}\right)$$
Show that every element $`A`$ of $`𝖮(2)`$ is of one of the two forms
$$A=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)$$
$$A=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)$$
(If $`A`$ is of the first form, then $`detA=1`$; if $`A`$ is of the second form, then $`detA=1`$.)
Hint: Recall that for $`A=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)`$ to be in $`𝖮(2)`$, the column vectors $`\left(\begin{array}{c}a\\ c\end{array}\right)`$ and $`\left(\begin{array}{c}b\\ d\end{array}\right)`$ must be unit vectors, and must be orthogonal.
7. The groups $`𝖮(1;1)`$ and $`\mathrm{𝖲𝖮}(1;1)`$. Show that
$$A=\left(\begin{array}{cc}\mathrm{cosh}t& \mathrm{sinh}t\\ \mathrm{sinh}t& \mathrm{cosh}t\end{array}\right)$$
is in $`\mathrm{𝖲𝖮}(1;1)`$, and that
$$\left(\begin{array}{cc}\mathrm{cosh}t& \mathrm{sinh}t\\ \mathrm{sinh}t& \mathrm{cosh}t\end{array}\right)\left(\begin{array}{cc}\mathrm{cosh}s& \mathrm{sinh}s\\ \mathrm{sinh}s& \mathrm{cosh}s\end{array}\right)=\left(\begin{array}{cc}\mathrm{cosh}(t+s)& \mathrm{sinh}(t+s)\\ \mathrm{sinh}(t+s)& \mathrm{cosh}(t+s)\end{array}\right)$$
Show that every element of $`𝖮(1;1)`$ can be written in one of the four forms
$$\left(\begin{array}{cc}\mathrm{cosh}t& \mathrm{sinh}t\\ \mathrm{sinh}t& \mathrm{cosh}t\end{array}\right)$$
$$\left(\begin{array}{cc}\mathrm{cosh}t& \mathrm{sinh}t\\ \mathrm{sinh}t& \mathrm{cosh}t\end{array}\right)$$
$$\left(\begin{array}{cc}\mathrm{cosh}t& \mathrm{sinh}t\\ \mathrm{sinh}t& \mathrm{cosh}t\end{array}\right)$$
$$\left(\begin{array}{cc}\mathrm{cosh}t& \mathrm{sinh}t\\ \mathrm{sinh}t& \mathrm{cosh}t\end{array}\right)$$
(Since $`\mathrm{cosh}t`$ is always positive, there is no overlap among the four cases. Matrices of the first two forms have determinant one; matrices of the last two forms have determinant minus one.)
Hint: For $`\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)`$ to be in $`𝖮(1;1)`$, we must have $`a^2c^2=1`$, $`b^2d^2=1`$, and $`abcd=0`$. The set of points $`(a,c)`$ in the plane with $`a^2c^2=1`$ (i.e., $`a=\pm \sqrt{1+c^2}`$ ) is a hyperbola.
8. The group$`\mathrm{𝖲𝖴}(2)`$. Show that if $`\alpha ,\beta `$ are arbitrary complex numbers satisfying $`\left|\alpha \right|^2+\left|\beta \right|^2=1`$, then the matrix
(2.9)
$$A=\left(\begin{array}{cc}\alpha & \overline{\beta }\\ \beta & \overline{\alpha }\end{array}\right)$$
is in $`\mathrm{𝖲𝖴}(2)`$. Show that every $`A\mathrm{𝖲𝖴}(2)`$ can be expressed in the form (2.9) for a unique pair $`(\alpha ,\beta )`$ satisfying $`\left|\alpha \right|^2+\left|\beta \right|^2=1`$. (Thus $`\mathrm{𝖲𝖴}(2)`$ can be thought of as the three-dimensional sphere $`S^3`$ sitting inside $`𝐂^2=^4`$. In particular, this shows that $`\mathrm{𝖲𝖴}(2)`$ is connected and simply connected.)
9. The groups$`\mathrm{𝖲𝗉}(1;)`$, $`\mathrm{𝖲𝗉}(1;𝐂)`$, and$`\mathrm{𝖲𝗉}\left(1\right)`$. Show that $`\mathrm{𝖲𝗉}(1;)=\mathrm{𝖲𝖫}(2;)`$, $`\mathrm{𝖲𝗉}(1;𝐂)=\mathrm{𝖲𝖫}(2;𝐂)`$, and $`\mathrm{𝖲𝗉}(1)=\mathrm{𝖲𝖴}(2)`$.
10. The Heisenberg group. Determine the center $`Z(H)`$ of the Heisenberg group $`H`$. Show that the quotient group $`H/Z(H)`$ is abelian.
11. Connectedness of$`\mathrm{𝖲𝖮}(n)`$. Show that $`\mathrm{𝖲𝖮}(n)`$ is connected, following the outline below.
For the case $`n=1`$, there is not much to show, since a $`1\times 1`$ matrix with determinant one must be $`\left[1\right]`$. Assume, then, that $`n2`$. Let $`e_1`$ denote the vector
$$e_1=\left(\begin{array}{c}1\\ 0\\ \mathrm{}\\ 0\end{array}\right)$$
in $`^n`$. Given any unit vector $`v^n`$, show that there exists a continuous path $`R(t)`$ in $`\mathrm{𝖲𝖮}(n)`$ with $`R(0)=I`$ and $`R(1)v=e_1`$. (Thus any unit vector can be “continuously rotated” to $`e_1`$.)
Now show that any element $`R`$ of $`\mathrm{𝖲𝖮}(n)`$ can be connected to an element of $`\mathrm{𝖲𝖮}(n1)`$, and proceed by induction.
12. The polar decomposition of$`\mathrm{𝖲𝖫}(n;)`$. Show that every element $`A`$ of $`\mathrm{𝖲𝖫}(n;)`$ can be written uniquely in the form $`A=RH`$, where $`R`$ is in $`\mathrm{𝖲𝖮}(n)`$, and $`H`$ is a symmetric, positive-definite matrix with determinant one. (That is, $`H^{tr}=H`$, and $`x,Hx0`$ for all $`x^n`$).
Hint: If $`A`$ could be written in this form, then we would have
$$A^{tr}A=H^{tr}R^{tr}RH=HR^1RH=H^2$$
Thus $`H`$ would have to be the unique positive-definite symmetric square root of $`A^{tr}A`$.
Note: A similar argument gives polar decompositions for $`\mathrm{𝖦𝖫}(n;)`$, $`\mathrm{𝖲𝖫}(n;)`$, and $`\mathrm{𝖦𝖫}(n;)`$. For example, every element $`A`$ of $`\mathrm{𝖲𝖫}(n;)`$ can be written uniquely as $`A=UH`$, with $`U`$ in $`\mathrm{𝖲𝖴}(n)`$, and $`H`$ a self-adjoint positive-definite matrix with determinant one.
13. The connectedness of$`\mathrm{𝖲𝖫}(n;)`$. Using the polar decomposition of $`\mathrm{𝖲𝖫}(n;)`$ (Ex. 12) and the connectedness of $`\mathrm{𝖲𝖮}(n)`$ (Ex. 11), show that $`\mathrm{𝖲𝖫}(n;)`$ is connected.
Hint: Recall that if $`H`$ is a real, symmetric matrix, then there exists a real orthogonal matrix $`R_1`$ such that $`H=R_1DR_1^1`$, where $`D`$ is diagonal.
14. The connectedness of$`\mathrm{𝖦𝖫}(n;)^+`$. Show that $`\mathrm{𝖦𝖫}(n;)^+`$ is connected.
15. Show that the set of translations is a normal subgroup of the Euclidean group, and also of the Poincaré group. Show that $`\left(𝖤(n)/\mathrm{translations}\right)𝖮(n)`$.
16. Harder. Show that every Lie group homomorphism $`\varphi `$ from $``$ to $`S^1`$ is of the form $`\varphi (x)=e^{iax}`$ for some $`a`$. In particular, every such homomorphism is smooth.
## Chapter 3 Lie Algebras and the Exponential Mapping
### 3.1. The Matrix Exponential
The exponential of a matrix plays a crucial role in the theory of Lie groups. The exponential enters into the definition of the Lie algebra of a matrix Lie group (Section 3.5 below), and is the mechanism for passing information from the Lie algebra to the Lie group. Since many computations are done much more easily at the level of the Lie algebra, the exponential is indispensable.
Let $`X`$ be an $`n\times n`$ real or complex matrix. We wish to define the exponential of $`X`$, $`e^X`$ or $`\mathrm{exp}X`$, by the usual power series
(3.1)
$$e^X=\underset{m=0}{\overset{\mathrm{}}{}}\frac{X^m}{m!}\text{.}$$
We will follow the convention of using letters such as $`X`$ and $`Y`$ for the variable in the matrix exponential.
###### Proposition 3.1.
For any $`n\times n`$ real or complex matrix $`X`$, the series (3.1) converges. The matrix exponential $`e^X`$ is a continuous function of $`X`$.
Before proving this, let us review some elementary analysis. Recall that the norm of a vector $`x`$ in $`^n`$ is defined to be
$$x=\sqrt{x,x}=\sqrt{\left|x_i\right|^2}\text{.}$$
This norm satisfies the triangle inequality
$$x+yx+y\text{.}$$
The norm of a matrix $`A`$ is defined to be
$$A=\underset{x0}{sup}\frac{Ax}{x}\text{.}$$
Equivalently, $`A`$ is the smallest number $`\lambda `$ such that $`Ax\lambda x`$ for all $`x𝐂^n`$.
It is not hard to see that for any $`n\times n`$ matrix $`A`$, $`A`$ is finite. Furthermore, it is easy to see that for any matrices $`A,B`$
(3.2) $`AB`$ $`AB`$
(3.3) $`A+B`$ $`A+B\text{.}`$
It is also easy to see that a sequence of matrices $`A_m`$ converges to a matrix $`A`$ if and only if $`A_mA0`$. (Compare this with Definition 2.1 of Chapter 2.)
A sequence of matrices $`A_m`$ is said to be a Cauchy sequence if $`A_mA_l0`$ as $`m,l\mathrm{}`$. Thinking of the space of matrices as $`^{n^2}`$ or $`^{n^2}`$, and using a standard result from analysis, we have the following:
###### Proposition 3.2.
If $`A_m`$ is a sequence of $`n\times n`$ real or complex matrices, and $`A_m`$ is a Cauchy sequence, then there exists a unique matrix $`A`$ such that $`A_m`$ converges to $`A`$.
That is, every Cauchy sequence converges.
Now, consider an infinite series whose terms are matrices:
(3.4)
$$A_0+A_1+A_2+\mathrm{}\text{.}$$
If
$$\underset{m=0}{\overset{\mathrm{}}{}}A_m<\mathrm{}$$
then the series (3.4) is said to converge absolutely. If a series converges absolutely, then it is not hard to show that the partial sums of the series form a Cauchy sequence, and hence by Proposition 3.2, the series converges. That is, any series which converges absolutely also converges. (The converse is not true; a series of matrices can converge without converging absolutely.)
###### Proof.
In light of (3.2), we see that
$$X^mX^m\text{,}$$
and hence
$$\underset{m=0}{\overset{\mathrm{}}{}}\frac{X^m}{m!}\underset{m=0}{\overset{\mathrm{}}{}}\frac{X^m}{m!}=e^X<\mathrm{}\text{.}$$
Thus the series (3.1) converges absolutely, and so it converges.
To show continuity, note that since $`X^m`$ is a continuous function of $`X`$, the partial sums of (3.1) are continuous. But it is easy to see that (3.1) converges uniformly on each set of the form $`\left\{XR\right\}`$, and so the sum is again continuous. ∎
###### Proposition 3.3.
Let $`X,Y`$ be arbitrary $`n\times n`$ matrices. Then
1. $`e^0=I`$.
2. $`e^X`$ is invertible, and $`\left(e^X\right)^1=e^X`$.
3. $`e^{(\alpha +\beta )X}=e^{\alpha X}e^{\beta X}`$ for all real or complex numbers $`\alpha ,\beta `$.
4. If $`XY=YX`$, then $`e^{X+Y}=e^Xe^Y=e^Ye^X`$.
5. If $`C`$ is invertible, then $`e^{CXC^1}=Ce^XC^1`$.
6. $`e^Xe^X`$.
It is not true in general that $`e^{X+Y}=e^Xe^Y`$, although by 4) it is true if $`X`$ and $`Y`$ commute. This is a crucial point, which we will consider in detail later. (See the Lie product formula in Section 3.4 and the Baker-Campbell-Hausdorff formula in Chapter 4.)
###### Proof.
Point 1) is obvious. Points 2) and 3) are special cases of point 4). To verify point 4), we simply multiply power series term by term. (It is left to the reader to verify that this is legal.) Thus
$$e^Xe^Y=\left(I+X+\frac{X^2}{2!}+\mathrm{}\right)\left(I+Y+\frac{Y^2}{2!}+\mathrm{}\right)\text{.}$$
Multiplying this out and collecting terms where the power of $`X`$ plus the power of $`Y`$ equals $`m`$, we get
(3.5)
$$e^Xe^Y=\underset{m=0}{\overset{\mathrm{}}{}}\underset{k=0}{\overset{m}{}}\frac{X^k}{k!}\frac{Y^{mk}}{(mk)!}=\underset{m=0}{\overset{\mathrm{}}{}}\frac{1}{m!}\underset{k=0}{\overset{m}{}}\frac{m!}{k!(mk)!}X^kY^{mk}\text{.}$$
Now because (and only because) $`X`$ and $`Y`$ commute,
$$(X+Y)^n=\underset{k=0}{\overset{m}{}}\frac{m!}{k!(mk)!}X^kY^{mk}\text{,}$$
and so (3.5) becomes
$$e^Xe^Y=\underset{m=0}{\overset{\mathrm{}}{}}\frac{1}{m!}(X+Y)^m=e^{X+Y}\text{.}$$
To prove 5), simply note that
$$\left(CXC^1\right)^m=CX^mC^1$$
and so the two sides of 5) are the same term by term.
Point 6) is evident from the proof of Proposition 3.1. ∎
###### Proposition 3.4.
Let $`X`$ be a $`n\times n`$ complex matrix, and view the space of all $`n\times n`$ complex matrices as $`^{n^2}`$. Then $`e^{tX}`$ is a smooth curve in $`^{n^2}`$, and
$$\frac{d}{dt}e^{tX}=Xe^{tX}=e^{tX}X\text{.}$$
In particular,
$$\frac{d}{dt}|_{t=0}e^{tX}=X\text{.}$$
###### Proof.
Differentiate the power series for $`e^{tX}`$ term-by-term. (You might worry whether this is valid, but you shouldn’t. For each $`i,j`$, $`\left(e^{tX}\right)_{ij}`$ is given by a convergent power series in $`t`$, and it is a standard theorem that you can differentiate power series term-by-term.) ∎
### 3.2. Computing the Exponential of a Matrix
#### 3.2.1. Case 1: $`X`$ is diagonalizable
Suppose that $`X`$ is a $`n\times n`$ real or complex matrix, and that $`X`$ is diagonalizable over $``$, that is, that there exists an invertible complex matrix $`C`$ such that $`X=CDC^1`$, with
$$D=\left(\begin{array}{ccc}\lambda _1& & 0\\ & \mathrm{}& \\ 0& & \lambda _n\end{array}\right)\text{.}$$
Observe that $`e^D`$ is the diagonal matrix with eigenvalues $`e^{\lambda _1},\mathrm{},e^{\lambda _n}`$, and so in light of Proposition 3.3, we have
$$e^X=C\left(\begin{array}{ccc}e^{\lambda _1}& & 0\\ & \mathrm{}& \\ 0& & e^{\lambda _n}\end{array}\right)C^1\text{.}$$
Thus if you can explicitly diagonalize $`X`$, you can explicitly compute $`e^X`$. Note that if $`X`$ is real, then although $`C`$ may be complex and the $`\lambda _i`$’s may be complex, $`e^X`$ must come out to be real, since each term in the series (3.1) is real.
For example, take
$$X=\left(\begin{array}{cc}0& a\\ a& 0\end{array}\right)\text{.}$$
Then the eigenvectors of $`X`$ are $`\left(\begin{array}{c}1\\ i\end{array}\right)`$ and $`\left(\begin{array}{c}i\\ 1\end{array}\right)`$, with eigenvalues $`ia`$ and $`ia`$, respectively. Thus the invertible matrix
$$C=\left(\begin{array}{cc}1& i\\ i& 1\end{array}\right)$$
maps the basis vectors $`\left(\begin{array}{c}1\\ 0\end{array}\right)`$ and $`\left(\begin{array}{c}0\\ 1\end{array}\right)`$ to the eigenvectors of $`X`$, and so (check) $`C^1XC`$ is a diagonal matrix $`D`$. Thus $`X=CDC^1`$:
$`e^X`$ $`=\left(\begin{array}{cc}1& i\\ i& 1\end{array}\right)\left(\begin{array}{cc}e^{ia}& 0\\ 0& e^{ia}\end{array}\right)\left(\begin{array}{cc}1/2& i/2\\ i/2& 1/2\end{array}\right)`$
$`=\left(\begin{array}{cc}\mathrm{cos}a& \mathrm{sin}a\\ \mathrm{sin}a& \mathrm{cos}a\end{array}\right)\text{.}`$
Note that explicitly if $`X`$ (and hence $`a`$) is real, then $`e^X`$ is real.
#### 3.2.2. Case 2: $`X`$ is nilpotent
An $`n\times n`$ matrix $`X`$ is said to be nilpotent if $`X^m=0`$ for some positive integer $`m`$. Of course, if $`X^m=0`$, then $`X^l=0`$ for all $`l>m`$. In this case the series (3.1) which defines $`e^X`$ terminates after the first $`m`$ terms, and so can be computed explicitly.
For example, compute $`e^{tX}`$, where
$$X=\left(\begin{array}{ccc}0& a& b\\ 0& 0& c\\ 0& 0& 0\end{array}\right)\text{.}$$
Note that
$$X^2=\left(\begin{array}{ccc}0& 0& ac\\ 0& 0& 0\\ 0& 0& 0\end{array}\right)$$
and that $`X^3=0`$. Thus
$$e^{tX}=\left(\begin{array}{ccc}1& ta& tb+\frac{1}{2}t^2ac\\ 0& 1& tc\\ 0& 0& 1\end{array}\right)\text{.}$$
#### 3.2.3. Case 3: $`X`$ arbitrary
A general matrix $`X`$ may be neither nilpotent nor diagonalizable. However, it follows from the Jordan canonical form that $`X`$ can be written (Exercise 2) in the form $`X=S+N`$ with $`S`$ diagonalizable, $`N`$ nilpotent, and $`SN=NS`$. (See Exercise 2.) Then, since $`N`$ and $`S`$ commute,
$$e^X=e^{S+N}=e^Se^N$$
and $`e^S`$ and $`e^N`$ can be computed as above.
For example, take
$$X=\left(\begin{array}{cc}a& b\\ 0& a\end{array}\right)\text{.}$$
Then
$$X=\left(\begin{array}{cc}a& 0\\ 0& a\end{array}\right)+\left(\begin{array}{cc}0& b\\ 0& 0\end{array}\right)\text{.}$$
The two terms clearly commute (since the first one is a multiple of the identity), and so
$$e^X=\left(\begin{array}{cc}e^a& 0\\ 0& e^a\end{array}\right)\left(\begin{array}{cc}1& b\\ 0& 1\end{array}\right)=\left(\begin{array}{cc}e^a& e^ab\\ 0& e^a\end{array}\right)\text{.}$$
### 3.3. The Matrix Logarithm
We wish to define a matrix logarithm, which should be an inverse function to the matrix exponential. Defining a logarithm for matrices should be at least as difficult as defining a logarithm for complex numbers, and so we cannot hope to define the matrix logarithm for all matrices, or even for all invertible matrices. We will content ourselves with defining the logarithm in a neighborhood of the identity matrix.
The simplest way to define the matrix logarithm is by a power series. We recall the situation for complex numbers:
###### Lemma 3.5.
The function
$$\mathrm{log}z=\underset{m=1}{\overset{\mathrm{}}{}}(1)^{m+1}\frac{(z1)^m}{m}$$
is defined and analytic in a circle of radius one about $`z=1`$.
For all $`z`$ with $`\left|z1\right|<1`$,
$$e^{\mathrm{log}z}=z\text{.}$$
For all $`u`$ with $`\left|u\right|<\mathrm{log}2`$, $`\left|e^u1\right|<1`$ and
$$\mathrm{log}e^u=u\text{.}$$
###### Proof.
The usual logarithm for real, positive numbers satisfies
$$\frac{d}{dx}\mathrm{log}(1x)=\frac{1}{1x}=\left(1+x+x^2+\mathrm{}\right)$$
for $`\left|x\right|<1`$. Integrating term-by-term and noting that $`\mathrm{log}1=0`$ gives
$$\mathrm{log}(1x)=\left(x+\frac{x^2}{2}+\frac{x^3}{3}+\mathrm{}\right)\text{.}$$
Taking $`z=1x`$ (so that $`x=1z`$), we have
$`\mathrm{log}z=\left((1z)+\frac{(1z)^2}{2}+\frac{(1z)^3}{3}+\mathrm{}\right)`$
$`={\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}(1)^{m+1}{\displaystyle \frac{(z1)^m}{m}}\text{.}`$
This series has radius of convergence one, and defines a complex analytic function on the set $`\left\{\left|z1\right|<1\right\}`$, which coincides with the usual logarithm for real $`z`$ in the interval $`(0,2)`$. Now, $`\mathrm{exp}(\mathrm{log}z)=z`$ for $`z(0,2)`$, and by analyticity this identity continues to hold on the whole set $`\left\{\left|z1\right|<1\right\}`$.
On the other hand, if $`\left|u\right|<\mathrm{log}2`$, then
$$\left|e^u1\right|=\left|u+\frac{u^2}{2!}+\mathrm{}\right|\left|u\right|+\frac{\left|u\right|^2}{2!}+\mathrm{}$$
so that
$$\left|e^u1\right|e^{\left|u\right|}1<1\text{.}$$
Thus $`\mathrm{log}(\mathrm{exp}u)`$ makes sense for all such $`u`$. Since $`\mathrm{log}(\mathrm{exp}u)=u`$ for real $`u`$ with $`\left|u\right|<\mathrm{log}2`$, it follows by analyticity that $`\mathrm{log}(\mathrm{exp}u)=u`$ for all complex numbers with $`\left|u\right|<\mathrm{log}2`$. ∎
###### Theorem 3.6.
The function
(3.6)
$$\mathrm{log}A=\underset{m=1}{\overset{\mathrm{}}{}}(1)^{m+1}\frac{(AI)^m}{m}$$
is defined and continuous on the set of all $`n\times n`$ complex matrices $`A`$ with $`AI<1`$, and $`\mathrm{log}A`$ is real if $`A`$ is real.
For all $`A`$ with $`AI<1`$,
$$e^{\mathrm{log}A}=A\text{.}$$
For all $`X`$ with $`X<\mathrm{log}2`$, $`e^X1<1`$ and
$$\mathrm{log}e^X=X\text{.}$$
###### Proof.
It is easy to see that the series (3.6) converges absolutely whenever $`AI<1`$. The proof of continuity is essentially the same as for the exponential. If $`A`$ is real, then every term in the series (3.6) is real, and so $`\mathrm{log}A`$ is real.
We will now show that $`\mathrm{exp}(\mathrm{log}A)=A`$ for all $`A`$ with $`AI<1`$. We do this by considering two cases.
Case 1: $`A`$ is diagonalizable.
Suppose that $`A=CDC^1`$, with $`D`$ diagonal. Then $`AI=CDC^1I=C(DI)C^1`$. It follows that $`(AI)^m`$ is of the form
$$(AI)^m=C\left(\begin{array}{ccc}(z_11)^m& & 0\\ & \mathrm{}& \\ 0& & (z_n1)^m\end{array}\right)C^1\text{,}$$
where $`z_1,\mathrm{},z_n`$ are the eigenvalues of $`A`$.
Now, if $`AI<1`$, then certainly $`\left|z_i1\right|<1`$ for $`i=1,\mathrm{},n`$. (Think about it.) Thus
$$\underset{m=1}{\overset{\mathrm{}}{}}(1)^{m+1}\frac{(AI)^m}{m}=C\left(\begin{array}{ccc}\mathrm{log}z_1& & 0\\ & \mathrm{}& \\ 0& & \mathrm{log}z_n\end{array}\right)C^1$$
and so by the Lemma
$$e^{\mathrm{log}A}=C\left(\begin{array}{ccc}e^{\mathrm{log}z_1}& & 0\\ & \mathrm{}& \\ 0& & e^{\mathrm{log}z_n}\end{array}\right)C^1=A\text{.}$$
Case 2: $`A`$ is not diagonalizable.
If $`A`$ is not diagonalizable, then, using the Jordan canonical form, it is not difficult to construct a sequence $`A_m`$ of diagonalizable matrices with $`A_mA`$. (See Exercise 4.) If $`AI<1`$, then $`A_mI<1`$ for all sufficiently large $`m`$. By Case 1, $`\mathrm{exp}(\mathrm{log}A_m)=A_m`$, and so by the continuity of $`\mathrm{exp}`$ and $`\mathrm{log}`$, $`\mathrm{exp}(\mathrm{log}A)=A`$.
Thus we have shown that $`\mathrm{exp}(\mathrm{log}A)=A`$ for all $`A`$ with $`AI<1`$. Now, the same argument as in the complex case shows that if $`X<\mathrm{log}2`$, then $`e^XI<1`$. But then the same two-case argument as above shows that $`\mathrm{log}(\mathrm{exp}X)=X`$ for all such $`X`$. ∎
###### Proposition 3.7.
There exists a constant $`c`$ such that for all $`n\times n`$ matrices $`B`$ with $`B<\frac{1}{2}`$
$$\mathrm{log}(I+B)BcB^2\text{.}$$
###### Proof.
Note that
$$\mathrm{log}(I+B)B=\underset{m=2}{\overset{\mathrm{}}{}}(1)^m\frac{B^m}{m}=B^2\underset{m=2}{\overset{\mathrm{}}{}}(1)^m\frac{B^{m2}}{m}$$
so that
$$\mathrm{log}(I+B)BB^2\underset{m=2}{\overset{\mathrm{}}{}}\frac{\left(\frac{1}{2}\right)^m}{m}\text{.}$$
This is what we want. ∎
###### Proposition 3.8.
Let $`X`$ be any $`n\times n`$ complex matrix, and let $`C_m`$ be a sequence of matrices such that $`C_m\frac{\mathrm{const}.}{m^2}`$. Then
$$\underset{m\mathrm{}}{lim}\left[I+\frac{X}{m}+C_m\right]^m=e^X\text{.}$$
###### Proof.
The expression inside the brackets is clearly tending to $`I`$ as $`m\mathrm{}`$, and so is in the domain of the logarithm for all sufficiently large $`m`$. Now
$$\mathrm{log}\left(I+\frac{X}{m}+C_m\right)=\frac{X}{m}+C_m+E_m$$
where $`E_m`$ is an error term which, by Proposition 3.7 satisfies $`E_mc\frac{X}{m}+C_m^2\frac{\mathrm{const}.}{m^2}`$. But then
$$I+\frac{X}{m}+C_m=\mathrm{exp}\left(\frac{X}{m}+C_m+E_m\right)\text{,}$$
and so
$$\left[I+\frac{X}{m}+C_m\right]^m=\mathrm{exp}\left(X+mC_m+mE_m\right)\text{.}$$
Since both $`C_m`$ and $`E_m`$ are of order $`\frac{1}{m^2}`$, we obtain the desired result by letting $`m\mathrm{}`$ and using the continuity of the exponential. ∎
### 3.4. Further Properties of the Matrix Exponential
In this section we give three additional results involving the exponential of a matrix, which will be important in our study of Lie algebras.
###### Theorem 3.9 (Lie Product Formula).
Let $`X`$ and $`Y`$ be $`n\times n`$ complex matrices. Then
$$e^{X+Y}=\underset{m\mathrm{}}{lim}\left(e^{\frac{X}{m}}e^{\frac{Y}{m}}\right)^m\text{.}$$
This theorem has a big brother, called the Trotter product formula, which gives the same result in the case where $`X`$ and $`Y`$ are suitable unbounded operators on an infinite-dimensional Hilbert space. The Trotter formula is described, for example, in M. Reed and B. Simon, Methods of Modern Mathematical Physics, Vol. I, VIII.8.
###### Proof.
Using the power series for the exponential and multiplying, we get
$$e^{\frac{X}{m}}e^{\frac{Y}{m}}=I+\frac{X}{m}+\frac{Y}{m}+C_m\text{,}$$
where (check!) $`C_m\frac{const.}{m^2}`$. Since $`e^{\frac{X}{m}}e^{\frac{Y}{m}}I`$ as $`m\mathrm{}`$, $`e^{\frac{X}{m}}e^{\frac{Y}{m}}`$ is in the domain of the logarithm for all sufficiently large $`m`$. But
$`\mathrm{log}\left(e^{\frac{X}{m}}e^{\frac{Y}{m}}\right)`$ $`=\mathrm{log}\left(I+{\displaystyle \frac{X}{m}}+{\displaystyle \frac{Y}{m}}+C_m\right)`$
$`={\displaystyle \frac{X}{m}}+{\displaystyle \frac{Y}{m}}+C_m+E_m`$
where by Proposition 3.7 $`C_mconst.\frac{X}{m}+\frac{Y}{m}+C_m^2\frac{const.}{m^2}`$. Exponentiating the logarithm gives
$$e^{\frac{X}{m}}e^{\frac{Y}{m}}=\mathrm{exp}\left(\frac{X}{m}+\frac{Y}{m}+C_m+E_m\right)$$
and
$$\left(e^{\frac{X}{m}}e^{\frac{Y}{m}}\right)^m=\mathrm{exp}\left(X+Y+mC_m+mE_m\right)\text{.}$$
Since both $`C_m`$ and $`E_m`$ are of order $`\frac{1}{m^2}`$, we have (using the continuity of the exponential)
$$\underset{m\mathrm{}}{lim}\left(e^{\frac{X}{m}}e^{\frac{Y}{m}}\right)^m=\mathrm{exp}\left(X+Y\right)$$
which is the Lie product formula. ∎
###### Theorem 3.10.
Let $`X`$ be an $`n\times n`$ real or complex matrix. Then
$$det\left(e^X\right)=e^{\mathrm{trace}(X)}\text{.}$$
###### Proof.
There are three cases, as in Section 3.2.
Case 1: $`A`$ is diagonalizable. Suppose there is a complex invertible matrix $`C`$ such that
$$X=C\left(\begin{array}{ccc}\lambda _1& & 0\\ & \mathrm{}& \\ 0& & \lambda _n\end{array}\right)C^1\text{.}$$
Then
$$e^X=C\left(\begin{array}{ccc}e^{\lambda _1}& & 0\\ & \mathrm{}& \\ 0& & e^{\lambda _n}\end{array}\right)C^1\text{.}$$
Thus $`\mathrm{trace}(X)=\lambda _i`$, and $`det(e^X)=e^{\lambda _i}=e^{{\scriptscriptstyle \lambda _i}}`$. (Recall that $`\mathrm{trace}(CDC^1)=\mathrm{trace}(D)`$.)
Case 2: $`X`$ is nilpotent. If $`X`$ is nilpotent, then it cannot have any non-zero eigenvalues (check!), and so all the roots of the characteristic polynomial must be zero. Thus the Jordan canonical form of $`X`$ will be strictly upper triangular. That is, $`X`$ can be written as
$$X=C\left(\begin{array}{ccc}0& & \\ & \mathrm{}& \\ 0& & 0\end{array}\right)C^1\text{.}$$
In that case (it is easy to see) $`e^X`$ will be upper triangular, with ones on the diagonal:
$$e^X=C\left(\begin{array}{ccc}1& & \\ & \mathrm{}& \\ 0& & 1\end{array}\right)C^1\text{.}$$
Thus if $`X`$ is nilpotent, $`\mathrm{trace}(X)=0`$, and $`det(e^X)=1`$.
Case 3: $`X`$ arbitrary. As pointed out in Section 3.2, every matrix $`X`$ can be written as the sum of two commuting matrices $`S`$ and $`N`$, with $`S`$ diagonalizable (over $``$) and $`N`$ nilpotent. Since $`S`$ and $`N`$ commute, $`e^X=e^Se^N`$. So by the two previous cases
$$det\left(e^X\right)=det\left(e^S\right)det\left(e^N\right)=e^{\mathrm{trace}(S)}e^{\mathrm{trace}(N)}=e^{\mathrm{trace}(X)}\text{,}$$
which is what we want. ∎
###### Definition 3.11.
A function $`A:\mathrm{𝖦𝖫}(n;)`$ is called a one-parameter group if
1. $`A`$ is continuous,
2. $`A(0)=I`$,
3. $`A(t+s)=A(t)A(s)`$ for all $`t,s`$.
###### Theorem 3.12 (One-parameter Subgroups).
If $`A`$ is a one-parameter group in $`\mathrm{𝖦𝖫}(n;)`$, then there exists a unique $`n\times n`$ complex matrix $`X`$ such that
$$A(t)=e^{tX}\text{.}$$
By taking $`n=1`$, and noting that $`\mathrm{𝖦𝖫}(1;)^{}`$, this Theorem provides an alternative method of solving Exercise 16 in Chapter 2.
###### Proof.
The uniqueness is immediate, since if there is such an $`X`$, then $`X=\frac{d}{dt}|_{t=0}A(t)`$. So we need only worry about existence.
The first step is to show that $`A(t)`$ must be smooth. This follows from Proposition 2.16 in Chapter 2 (which we did not prove), but we give a self-contained proof. Let $`f(s)`$ be a smooth real-valued function supported in a small neighborhood of zero, with $`f(s)0`$ and $`f(s)𝑑s=1`$. Now look at
(3.7)
$$B(t)=A(t+s)f(s)𝑑s\text{.}$$
Making the change-of-variable $`u=t+s`$ gives
$$B(t)=A(u)f(ut)𝑑u\text{.}$$
It follows that $`B(t)`$ is differentiable, since derivatives in the $`t`$ variable go onto $`f`$, which is smooth.
On the other hand, if we use the identity $`A(t+s)=A(t)A(s)`$ in (3.7), we have
$$B(t)=A(t)A(s)f(s)𝑑s\text{.}$$
Now, the conditions on the function $`f`$, together with the continuity of $`A`$, guarantee that $`A(s)f(s)𝑑s`$ is close to $`A(0)=I`$, and hence is invertible. Thus we may write
(3.8)
$$A(t)=B(t)\left(A(s)f(s)𝑑s\right)^1\text{.}$$
Since $`B\left(t\right)`$ is smooth and $`A(s)f(s)𝑑s`$ is just a constant matrix, this shows that $`A\left(t\right)`$ is smooth.
Now that $`A(t)`$ is known to be differentiable, we may define
$$X=\frac{d}{dt}|_{t=0}A(t)\text{.}$$
Our goal is to show that $`A(t)=e^{tX}`$. Since $`A(t)`$ is smooth, a standard calculus result (extended trivially to handle matrix-valued functions) says
$$A(t)(I+tX)\mathrm{const}.t^2\text{.}$$
It follows that for each fixed $`t`$,
$$A\left(\frac{t}{m}\right)=I+\frac{t}{m}X+O\left(\frac{1}{m^2}\right)\text{.}$$
Then, since $`A`$ is a one-parameter group
$$A(t)=\left[A\left(\frac{t}{m}\right)\right]^m=\left[I+\frac{t}{m}X+O\left(\frac{1}{m^2}\right)\right]^m\text{.}$$
Letting $`m\mathrm{}`$ and using Proposition 3.8 from Section 3.3 shows that $`A(t)=e^{tX}`$. ∎
### 3.5. The Lie Algebra of a Matrix Lie Group
The Lie algebra is an indispensable tool in studying matrix Lie groups. On the one hand, Lie algebras are simpler than matrix Lie groups, because (as we will see) the Lie algebra is a linear space. Thus we can understand much about Lie algebras just by doing linear algebra. On the other hand, the Lie algebra of a matrix Lie group contains much information about that group. (See for example, Proposition 3.23 in Section 3.7, and the Baker-Campbell-Hausdorff Formula (Chapter 4).) Thus many questions about matrix Lie groups can be answered by considering a similar but easier problem for the Lie algebra.
###### Definition 3.13.
Let $`G`$ be a matrix Lie group. Then the Lie algebra of $`G`$, denoted $`𝔤`$, is the set of all matrices $`X`$ such that $`e^{tX}`$ is in $`G`$ for all real numbers $`t`$.
Note that even if $`G`$ is a subgroup of $`\mathrm{𝖦𝖫}(n;)`$ we do not require that $`e^{tX}`$ be in $`G`$ for all complex $`t`$, but only for all real $`t`$. Also, it is definitely not enough to have just $`e^X`$ in $`G`$. That is, it is easy to give an example of an $`X`$ and a $`G`$ such that $`e^XG`$ but $`e^{tX}G`$ for some values of $`t`$. Such an $`X`$ is not in the Lie algebra of $`G`$.
It is customary to use lower case Gothic (Fraktur) characters such as $`𝔤`$ and $`𝔥`$ to refer to Lie algebras.
#### 3.5.1. Physicists’ Convention
Physicists are accustomed to considering the map $`Xe^{iX}`$ instead of $`Xe^X`$. Thus a physicist would think of the Lie algebra of $`G`$ as the set of all matrices $`X`$ such that $`e^{itX}G`$ for all real $`t`$. In the physics literature, the Lie algebra is frequently referred to as the space of “infinitesimal group elements.” See Bröcker and tom Dieck, Chapter I, 2.21. The physics literature does not always distinguish clearly between a matrix Lie group and its Lie algebra.
Before examining general properties of the Lie algebra, let us compute the Lie algebras of the matrix Lie groups introduced in the previous chapter.
#### 3.5.2. The general linear groups
If $`X`$ is any $`n\times n`$ complex matrix, then by Proposition 3.3, $`e^{tX}`$ is invertible. Thus the Lie algebra of $`\mathrm{𝖦𝖫}(n;)`$ is the space of all $`n\times n`$ complex matrices. This Lie algebra is denoted $`\mathrm{𝗀𝗅}(n;)`$.
If $`X`$ is any $`n\times n`$ real matrix, then $`e^{tX}`$ will be invertible and real. On the other hand, if $`e^{tX}`$ is real for all real $`t`$, then $`X=\frac{d}{dt}|_{t=0}e^{tX}`$ will also be real. Thus the Lie algebra of $`\mathrm{𝖦𝖫}(n;)`$ is the space of all $`n\times n`$ real matrices, denoted $`\mathrm{𝗀𝗅}(n;)`$.
Note that the preceding argument shows that if $`G`$ is a subgroup of $`\mathrm{𝖦𝖫}(n;)`$, then the Lie algebra of $`G`$ must consist entirely of real matrices. We will use this fact when appropriate in what follows.
#### 3.5.3. The special linear groups
Recall Theorem 3.10: $`det\left(e^X\right)=e^{\mathrm{trace}X}`$. Thus if $`\mathrm{trace}X=0`$, then $`det\left(e^{tX}\right)=1`$ for all real $`t`$. On the other hand, if $`X`$ is any $`n\times n`$ matrix such that $`det\left(e^{tX}\right)=1`$ for all $`t`$, then $`e^{(t)(\mathrm{trace}X)}=1`$ for all $`t`$. This means that $`(t)(\mathrm{trace}X)`$ is an integer multiple of $`2\pi i`$ for all $`t`$, which is only possible if $`\mathrm{trace}X=0`$. Thus the Lie algebra of $`\mathrm{𝖲𝖫}(n;)`$ is the space of all $`n\times n`$ complex matrices with trace zero, denoted $`\mathrm{𝗌𝗅}(n;)`$.
Similarly, the Lie algebra of $`\mathrm{𝖲𝖫}(n;)`$ is the space of all $`n\times n`$ real matrices with trace zero, denoted $`\mathrm{𝗌𝗅}(n;)`$.
#### 3.5.4. The unitary groups
Recall that a matrix $`U`$ is unitary if and only if $`U^{}=U^1`$. Thus $`e^{tX}`$ is unitary if and only if
(3.9)
$$\left(e^{tX}\right)^{}=\left(e^{tX}\right)^1=e^{tX}\text{.}$$
But by taking adjoints term-by-term, we see that $`\left(e^{tX}\right)^{}=e^{tX^{}}`$, and so (3.9) becomes
(3.10)
$$e^{tX^{}}=e^{tX}\text{.}$$
Clearly, a sufficient condition for (3.10) to hold is that $`X^{}=X`$. On the other hand, if (3.10) holds for all $`t`$, then by differentiating at $`t=0`$, we see that $`X^{}=X`$ is necessary.
Thus the Lie algebra of $`𝖴(n)`$ is the space of all $`n\times n`$ complex matrices $`X`$ such that $`X^{}=X`$, denoted $`𝗎(n)`$.
By combining the two previous computations, we see that the Lie algebra of $`\mathrm{𝖲𝖴}(n)`$ is the space of all $`n\times n`$ complex matrices $`X`$ such that $`X^{}=X`$ and $`\mathrm{trace}X=0`$, denoted $`\mathrm{𝗌𝗎}(n)`$.
#### 3.5.5. The orthogonal groups
The identity component of $`𝖮(n)`$ is just $`\mathrm{𝖲𝖮}(n)`$. Since (Proposition 3.14) the exponential of a matrix in the Lie algebra is automatically in the identity component, the Lie algebra of $`𝖮(n)`$ is the same as the Lie algebra of $`\mathrm{𝖲𝖮}(n)`$.
Now, an $`n\times n`$ real matrix $`R`$ is orthogonal if and only if $`R^{tr}=R^1`$. So, given an $`n\times n`$ real matrix $`X`$, $`e^{tX}`$ is orthogonal if and only if $`(e^{tX})^{tr}=(e^{tX})^1`$, or
(3.11)
$$e^{tX^{tr}}=e^{tX}\text{.}$$
Clearly, a sufficient condition for this to hold is that $`X^{tr}=X`$. If (3.11) holds for all $`t`$, then by differentiating at $`t=0`$, we must have $`X^{tr}=X`$.
Thus the Lie algebra of $`𝖮(n)`$, as well as the Lie algebra of $`\mathrm{𝖲𝖮}(n)`$, is the space of all $`n\times n`$ real matrices $`X`$ with $`X^{tr}=X`$, denoted $`\mathrm{𝗌𝗈}(n)`$. Note that the condition $`X^{tr}=X`$ forces the diagonal entries of $`X`$ to be zero, and so explicitly the trace of $`X`$ is zero.
The same argument shows that the Lie algebra of $`\mathrm{𝖲𝖮}(n;)`$ is the space of $`n\times n`$ complex matrices satisfying $`X^{tr}=X`$, denoted $`\mathrm{𝗌𝗈}(n;)`$. This is not the same as $`\mathrm{𝗌𝗎}(n)`$.
#### 3.5.6. The generalized orthogonal groups
A matrix $`A`$ is in $`𝖮(n;k)`$ if and only if $`A^{tr}gA=g`$, where $`g`$ is the $`(n+k)\times (n+k)`$ diagonal matrix with the first $`n`$ diagonal entries equal to one, and the last $`k`$ diagonal entries equal to minus one. This condition is equivalent to the condition $`g^1A^{tr}g=A^1`$, or, since explicitly $`g^1=g`$, $`gA^{tr}g=A^1`$. Now, if $`X`$ is an $`(n+k)\times (n+k)`$ real matrix, then $`e^{tX}`$ is in $`𝖮(n;k)`$ if and only if
$$ge^{tX^{tr}}g=e^{tgX^{tr}g}=e^{tX}\text{.}$$
This condition holds for all real $`t`$ if and only if $`gX^{tr}g=X`$. Thus the Lie algebra of $`𝖮(n;k)`$, which is the same as the Lie algebra of $`\mathrm{𝖲𝖮}(n;k)`$, consists of all $`(n+k)\times (n+k)`$ real matrices $`X`$ with $`gX^{tr}g=X`$. This Lie algebra is denoted $`\mathrm{𝗌𝗈}(n;k)`$.
(In general, the group $`\mathrm{𝖲𝖮}(n;k)`$ will not be connected, in contrast to the group $`\mathrm{𝖲𝖮}(n)`$. The identity component of $`\mathrm{𝖲𝖮}(n;k)`$, which is also the identity component of $`𝖮(n;k)`$, is denoted $`\mathrm{𝖲𝖮}(n;k)_I`$. The Lie algebra of $`\mathrm{𝖲𝖮}(n;k)_I`$ is the same as the Lie algebra of $`\mathrm{𝖲𝖮}(n;k)`$.)
#### 3.5.7. The symplectic groups
These are denoted $`\mathrm{𝗌𝗉}(n;),`$ sp$`(n;),`$ and $`\mathrm{𝗌𝗉}\left(n\right).`$ The calculation of these Lie algebras is similar to that of the generalized orthogonal groups, and I will just record the result here. Let $`J`$ be the matrix in the definition of the symplectic groups. Then $`\mathrm{𝗌𝗉}(n;)`$ is the space of $`2n\times 2n`$ real matrices $`X`$ such that $`JX^{tr}J=X,`$ sp$`(n;)`$ is the space of $`2n\times 2n`$ complex matrices satisfying the same condition, and $`\mathrm{𝗌𝗉}\left(n\right)=`$sp$`(n;)𝗎\left(2n\right).`$
#### 3.5.8. The Heisenberg group
Recall the Heisenberg group $`H`$ is the group of all $`3\times 3`$ real matrices $`A`$ of the form
(3.12)
$$A=\left(\begin{array}{ccc}1& a& b\\ 0& 1& c\\ 0& 0& 1\end{array}\right)$$
Recall also that in Section 3.2, Case 2, we computed the exponential of a matrix of the form
(3.13)
$$X=\left(\begin{array}{ccc}0& \alpha & \beta \\ 0& 0& \gamma \\ 0& 0& 0\end{array}\right)$$
and saw that $`e^X`$ was in $`H`$. On the other hand, if $`X`$ is any matrix such that $`e^{tX}`$ is of the form (3.12), then all of the entries of $`X=\frac{d}{dt}|_{t=0}e^{tX}`$ which are on or below the diagonal must be zero, so that $`X`$ is of form (3.13).
Thus the Lie algebra of the Heisenberg group is the space of all $`3\times 3`$ real matrices which are strictly upper triangular.
#### 3.5.9. The Euclidean and Poincaré groups
Recall that the Euclidean group $`𝖤(n)`$ is (or can be thought of as) the group of $`(n+1)\times (n+1)`$ real matrices of the form
$$\left(\begin{array}{cccc}& & & x_1\\ & R& & \mathrm{}\\ & & & x_n\\ 0& \mathrm{}& 0& 1\end{array}\right)$$
with $`R𝖮(n)`$. Now if $`X`$ is an $`(n+1)\times (n+1)`$ real matrix such that $`e^{tX}`$ is in $`𝖤(n)`$ for all $`t`$, then $`X=\frac{d}{dt}|_{t=0}e^{tX}`$ must be zero along the bottom row:
(3.14)
$$X=\left(\begin{array}{cccc}& & & y_1\\ & Y& & \mathrm{}\\ & & & y_n\\ 0& \mathrm{}& & 0\end{array}\right)$$
Our goal, then, is to determine which matrices of the form (3.14) are actually in the Lie algebra of the Euclidean group. A simple computation shows that for $`n1`$
$$\left(\begin{array}{cccc}& & & y_1\\ & Y& & \mathrm{}\\ & & & y_n\\ 0& \mathrm{}& & 0\end{array}\right)^n=\left(\begin{array}{cccc}& & & \\ & Y^n& & Y^{n1}y\\ & & & \\ 0& \mathrm{}& & 0\end{array}\right),$$
where $`y`$ is the column vector with entries $`y_1,\mathrm{},y_n.`$ It follows that if $`X`$ is as in (3.14), then $`e^{tX}`$ is of the form
$$e^{tX}=\left(\begin{array}{cccc}& & & \\ & e^{tY}& & \mathrm{}\\ & & & \\ 0& \mathrm{}& 0& 1\end{array}\right)\text{.}$$
Now, we have already established that $`e^{tY}`$ is in $`𝖮(n)`$ for all $`t`$ if and only if $`Y^{tr}=Y`$. Thus we see that the Lie algebra of $`𝖤(n)`$ is the space of all $`(n+1)\times (n+1)`$ real matrices of the form (3.14) with $`Y`$ satisfying $`Y^{tr}=Y`$.
A similar argument shows that the Lie algebra of $`𝖯(n;1)`$ is the space of all $`(n+2)\times (n+2)`$ real matrices of the form
$$\left(\begin{array}{cccc}& & & y_1\\ & Y& & \mathrm{}\\ & & & y_{n+1}\\ 0& \mathrm{}& & 0\end{array}\right)$$
with $`Y\mathrm{𝗌𝗈}(n;1)`$.
### 3.6. Properties of the Lie Algebra
We will now establish various basic properties of the Lie algebra of a matrix Lie group. The reader is invited to verify by direct calculation that these general properties hold for the examples computed in the previous section.
###### Proposition 3.14.
Let $`G`$ be a matrix Lie group, and $`X`$ an element of its Lie algebra. Then $`e^X`$ is an element of the identity component of $`G.`$
###### Proof.
By definition of the Lie algebra, $`e^{tX}`$ lies in $`G`$ for all real $`t`$. But as $`t`$ varies from $`0`$ to $`1`$, $`e^{tX}`$ is a continuous path connecting the identity to $`e^X`$. ∎
###### Proposition 3.15.
Let $`G`$ be a matrix Lie group, with Lie algebra $`𝔤`$. Let $`X`$ be an element of $`𝔤`$, and $`A`$ an element of $`G`$. Then $`AXA^1`$ is in $`𝔤`$.
###### Proof.
This is immediate, since by Proposition 3.3,
$$e^{t(AXA^1)}=Ae^{tX}A^1\text{,}$$
and $`Ae^{tX}A^1G`$. ∎
###### Theorem 3.16.
Let $`G`$ be a matrix Lie group, $`𝔤`$ its Lie algebra, and $`X,Y`$ elements of $`𝔤`$. Then
1. $`sX𝔤`$ for all real numbers $`s`$,
2. $`X+Y𝔤`$,
3. $`XYYX𝔤`$.
If you are following the physics convention for the definition of the Lie algebra, then condition 3 should be replaced with the condition $`i\left(XYYX\right)𝔤`$.
###### Proof.
Point 1 is immediate, since $`e^{t(sX)}=e^{(ts)X}`$, which must be in $`G`$ if $`X`$ is in $`𝔤`$. Point 2 is easy to verify if $`X`$ and $`Y`$ commute, since then $`e^{t(X+Y)}=e^{tX}e^{tY}`$. If $`X`$ and $`Y`$ do not commute, this argument does not work. However, the Lie product formula says that
$$e^{t(X+Y)}=\underset{m\mathrm{}}{lim}\left(e^{tX/m}e^{tY/m}\right)^m\text{.}$$
Because $`X`$ and $`Y`$ are in the Lie algebra, $`e^{tX/m}`$ and $`e^{tY/m}`$ are in $`G`$, as is $`\left(e^{tX/m}e^{tY/m}\right)^m`$, since $`G`$ is a group. But now because $`G`$ is a matrix Lie group, the limit of things in $`G`$ must be again in $`G`$, provided that the limit is invertible. Since $`e^{t(X+Y)}`$ is automatically invertible, we conclude that it must be in $`G`$. This shows that $`X+Y`$ is in $`𝔤`$.
Now for point 3. Recall (Proposition 3.4) that $`\frac{d}{dt}|_{t=0}e^{tX}=X`$. It follows that $`\frac{d}{dt}|_{t=0}e^{tX}Y=XY`$, and hence by the product rule (Exercise 1)
$`{\displaystyle \frac{d}{dt}}|_{t=0}\left(e^{tX}Ye^{tX}\right)=(XY)e^0+(e^0Y)(X)`$
$`=XYYX\text{.}`$
But now, by Proposition 3.15, $`e^{tX}Ye^{tX}`$ is in $`𝔤`$ for all $`t`$. Since we have (by points 1 and 2) established that $`𝔤`$ is a real vector space, it follows that the derivative of any smooth curve lying in $`𝔤`$ must be again in $`𝔤`$. Thus $`XYYX`$ is in $`𝔤`$. ∎
###### Definition 3.17.
Given two $`n\times n`$ matrices $`A`$ and $`B`$, the bracket (or commutator) of $`A`$ and $`B`$ is defined to be simply
$$[A,B]=ABBA\text{.}$$
According to Theorem 3.16, the Lie algebra of any matrix Lie group is closed under brackets.
The following very important theorem tells us that a Lie group homomorphism between two Lie groups gives rise in a natural way to a map between the corresponding Lie algebras. In particular, this will tell us that two isomorphic Lie groups have “the same” Lie algebras. (That is, the Lie algebras are isomorphic in the sense of Section 3.8.) See Exercise 6.
###### Theorem 3.18.
Let $`G`$ and $`H`$ be matrix Lie groups, with Lie algebras $`𝔤`$ and $`𝔥`$, respectively. Suppose that $`\varphi :GH`$ be a Lie group homomorphism. Then there exists a unique real linear map $`\stackrel{~}{\varphi }:𝔤𝔥`$ such that
$$\varphi (e^X)=e^{\stackrel{~}{\varphi }(X)}$$
for all $`X𝔤`$. The map $`\stackrel{~}{\varphi }`$ has following additional properties
1. $`\stackrel{~}{\varphi }\left(AXA^1\right)=\varphi (A)\stackrel{~}{\varphi }(X)\varphi (A)^1`$, for all $`X𝔤`$, $`AG`$.
2. $`\stackrel{~}{\varphi }([X,Y])=[\stackrel{~}{\varphi }(X),\stackrel{~}{\varphi }(Y)]`$, for all $`X,Y𝔤`$.
3. $`\stackrel{~}{\varphi }(X)=\frac{d}{dt}|_{t=0}\varphi (e^{tX})`$, for all $`X𝔤`$.
If $`G`$, $`H`$, and $`K`$ are matrix Lie groups and $`\varphi :HK`$ and $`\psi :GH`$ are Lie group homomorphisms, then
$$\stackrel{~}{\varphi \psi }=\stackrel{~}{\varphi }\stackrel{~}{\psi }\text{.}$$
In practice, given a Lie group homomorphism $`\varphi `$, the way one goes about computing $`\stackrel{~}{\varphi }`$ is by using Property 3. Of course, since $`\stackrel{~}{\varphi }`$ is (real) linear, it suffices to compute $`\stackrel{~}{\varphi }`$ on a basis for $`𝔤`$. In the language of differentiable manifolds, Property 3 says that $`\stackrel{~}{\varphi }`$ is the derivative (or differential) of $`\varphi `$ at the identity, which is the standard definition of $`\stackrel{~}{\varphi }`$. (See also Exercise 19.)
A linear map with property (2) is called a Lie algebra homomorphism. (See Section 3.8.) This theorem says that every Lie group homomorphism gives rise to a Lie algebra homomorphism. We will see eventually that the converse is true under certain circumstances. Specifically, suppose that $`G`$ and $`H`$ are Lie groups, and $`\stackrel{~}{\varphi }:𝔤𝔥`$ is a Lie algebra homomorphism. If $`G`$ is connected and simply connected, then there exists a unique Lie group homomorphism $`\varphi :GH`$ such that $`\varphi `$ and $`\stackrel{~}{\varphi }`$ are related as in Theorem 3.18.
###### Proof.
The proof is similar to the proof of Theorem 3.16. Since $`\varphi `$ is a continuous group homomorphism, $`\varphi (e^{tX})`$ will be a one-parameter subgroup of $`H`$, for each $`X𝔤`$. Thus by Theorem 3.12, there is a unique $`Z`$ such that
(3.15)
$$\varphi \left(e^{tX}\right)=e^{tZ}$$
for all $`t`$. This $`Z`$ must lie in $`𝔥`$ since $`e^{tZ}=\varphi \left(e^{tX}\right)H`$.
We now define $`\stackrel{~}{\varphi }(X)=Z`$, and check in several steps that $`\stackrel{~}{\varphi }`$ has the required properties.
Step 1: $`\varphi (e^X)=e^{\stackrel{~}{\varphi }(X)}`$.
This follows from (3.15) and our definition of $`\stackrel{~}{\varphi }`$, by putting $`t=1`$.
Step 2: $`\stackrel{~}{\varphi }(sX)=s\stackrel{~}{\varphi }(X)`$ for all $`s`$.
This is immediate, since if $`\varphi (e^{tX})=e^{tZ}`$, then $`\varphi (e^{tsX})=e^{tsZ}`$.
Step 3: $`\stackrel{~}{\varphi }(X+Y)=\stackrel{~}{\varphi }(X)+\stackrel{~}{\varphi }(Y)`$.
By Steps 1 and 2,
$$e^{t\stackrel{~}{\varphi }(X+Y)}=e^{\stackrel{~}{\varphi }[t(X+Y)]}=\varphi \left(e^{t(X+Y)}\right)\text{.}$$
By the Lie product formula, and the fact that $`\varphi `$ is a continuous homomorphism:
$`=\varphi \left(\underset{m\mathrm{}}{lim}\left(e^{tX/m}e^{tY/m}\right)^m\right)`$
$`=\underset{m\mathrm{}}{lim}\left(\varphi \left(e^{tX/m}\right)\varphi (e^{tY/m})\right)^m\text{.}`$
But then we have
$$e^{t\stackrel{~}{\varphi }(X+Y)}=\underset{m\mathrm{}}{lim}\left(e^{t\stackrel{~}{\varphi }(X)/m}e^{t\stackrel{~}{\varphi }(Y)/m}\right)^m=e^{t\left(\stackrel{~}{\varphi }(X)+\stackrel{~}{\varphi }(Y)\right)}\text{.}$$
Differentiating this result at $`t=0`$ gives the desired result.
Step 4: $`\stackrel{~}{\varphi }\left(AXA^1\right)=\varphi (A)\stackrel{~}{\varphi }(X)\varphi (A)^1`$.
By Steps 1 and 2,
$$\mathrm{exp}t\stackrel{~}{\varphi }(AXA^1)=\mathrm{exp}\stackrel{~}{\varphi }(tAXA^1)=\varphi \left(\mathrm{exp}tAXA^1\right)\text{.}$$
Using a property of the exponential and Step 1, this becomes
$`\mathrm{exp}t\stackrel{~}{\varphi }(AXA^1)=\varphi \left(Ae^{tX}A^1\right)=\varphi (A)\varphi (e^{tX})\varphi (A)^1`$
$`=\varphi (A)e^{t\stackrel{~}{\varphi }(X)}\varphi (A)^1\text{.}`$
Differentiating this at $`t=0`$ gives the desired result.
Step 5: $`\stackrel{~}{\varphi }([X,Y])=[\stackrel{~}{\varphi }(X),\stackrel{~}{\varphi }(Y)]`$.
Recall from the proof of Theorem 3.16 that
$$[X,Y]=\frac{d}{dt}|_{t=0}e^{tX}Ye^{tX}\text{.}$$
Hence
$$\stackrel{~}{\varphi }\left([X,Y]\right)=\stackrel{~}{\varphi }\left(\frac{d}{dt}|_{t=0}e^{tX}Ye^{tX}\right)=\frac{d}{dt}|_{t=0}\stackrel{~}{\varphi }\left(e^{tX}Ye^{tX}\right)$$
where we have used the fact that a derivative commutes with a linear transformation.
But then by Step 4,
$`\stackrel{~}{\varphi }\left([X,Y]\right)`$ $`=\frac{d}{dt}|_{t=0}\varphi (e^{tX})\stackrel{~}{\varphi }(Y)\varphi (e^{tX})`$
$`=\frac{d}{dt}|_{t=0}e^{t\stackrel{~}{\varphi }(X)}\stackrel{~}{\varphi }(Y)e^{t\stackrel{~}{\varphi }(X)}`$
$`=[\stackrel{~}{\varphi }(X),\stackrel{~}{\varphi }(Y)]\text{.}`$
Step 6: $`\stackrel{~}{\varphi }(X)=\frac{d}{dt}|_{t=0}\varphi (e^{tX})`$.
This follows from (3.15) and our definition of $`\stackrel{~}{\varphi }`$.
Step 7: $`\stackrel{~}{\varphi }`$ is the unique real-linear map such that $`\varphi (e^X)=e^{\stackrel{~}{\varphi }(X)}`$.
Suppose that $`\psi `$ is another such map. Then
$$e^{t\psi (X)}=e^{\psi (tX)}=\varphi (e^{tX})$$
so that
$$\psi (X)=\frac{d}{dt}|_{t=0}\varphi (e^{tX})\text{.}$$
Thus by Step 6, $`\psi `$ coincides with $`\stackrel{~}{\varphi }`$.
Step 8: $`\stackrel{~}{\varphi \psi }=\stackrel{~}{\varphi }\stackrel{~}{\psi }`$.
For any $`X𝔤`$,
$$\varphi \psi \left(e^{tX}\right)=\varphi \left(\psi \left(e^{tX}\right)\right)=\varphi \left(e^{t\stackrel{~}{\psi }(X)}\right)=e^{t\stackrel{~}{\varphi }(\stackrel{~}{\psi }(X))}\text{.}$$
Thus $`\stackrel{~}{\varphi \psi }(X)=\stackrel{~}{\varphi }\stackrel{~}{\psi }(X)`$. ∎
###### Definition 3.19 (The Adjoint Mapping).
Let $`G`$ be a matrix Lie group, with Lie algebra $`𝔤`$. Then for each $`AG`$, define a linear map $`\mathrm{Ad}A:𝔤𝔤`$ by the formula
$$\mathrm{Ad}A(X)=AXA^1\text{.}$$
We will let $`\mathrm{Ad}`$ denote the map $`A\mathrm{Ad}A`$.
###### Proposition 3.20.
Let $`G`$ be a matrix Lie group, with Lie algebra $`𝔤`$. Then for each $`AG`$, $`\mathrm{Ad}A`$ is an invertible linear transformation of $`𝔤`$ with inverse $`\mathrm{Ad}A^1`$, and $`\mathrm{Ad}:G\mathrm{𝖦𝖫}(𝔤)`$ is a group homomorphism.
###### Proof.
Easy. Note that Proposition 3.15 guarantees that $`\mathrm{Ad}A(X)`$ is actually in $`𝔤`$ for all $`X𝔤`$. ∎
Since $`𝔤`$ is a real vector space with some dimension $`k`$, $`\mathrm{𝖦𝖫}(𝔤)`$ is essentially the same as $`\mathrm{𝖦𝖫}(k;)`$. Thus we will regard $`\mathrm{𝖦𝖫}(𝔤)`$ as a matrix Lie group. It is easy to show that $`\mathrm{Ad}:G\mathrm{𝖦𝖫}(𝔤)`$ is continuous, and so is a Lie group homomorphism. By Theorem 3.18, there is an associated real linear map $`\stackrel{~}{\mathrm{Ad}}`$ from the Lie algebra of $`G`$ to the Lie algebra of $`\mathrm{𝖦𝖫}(𝔤)`$, i.e., from $`𝔤`$ to gl$`(𝔤)`$, with the property that
$$e^{\stackrel{~}{\mathrm{Ad}}X}=\mathrm{Ad}\left(e^X\right)\text{.}$$
###### Proposition 3.21.
Let $`G`$ be a matrix Lie group, let $`𝔤`$ its Lie algebra, and let $`\mathrm{Ad}:G\mathrm{𝖦𝖫}(𝔤)`$ be the Lie group homomorphism defined above. Let $`\stackrel{~}{\mathrm{Ad}}:𝔤`$gl$`(𝔤)`$ be the associated Lie algebra map. Then for all $`X,Y𝔤`$
$$\stackrel{~}{\mathrm{Ad}}X(Y)=[X,Y]\text{.}$$
###### Proof.
Recall that by Theorem 3.18, $`\stackrel{~}{\mathrm{Ad}}`$ can be computed as follows:
$$\stackrel{~}{\mathrm{Ad}}X=\frac{d}{dt}|_{t=0}\mathrm{Ad}(e^{tX})\text{.}$$
Thus
$`\stackrel{~}{\mathrm{Ad}}X(Y)`$ $`=\frac{d}{dt}|_{t=0}\mathrm{Ad}(e^{tX})(Y)=\frac{d}{dt}|_{t=0}e^{tX}Ye^{tX}`$
$`=[X,Y]`$
which is what we wanted to prove. See also Exercise 13. ∎
### 3.7. The Exponential Mapping
###### Definition 3.22.
If $`G`$ is a matrix Lie group with Lie algebra $`𝔤`$, then the exponential mapping for $`G`$ is the map
$$\mathrm{exp}:𝔤G\text{.}$$
In general the exponential mapping is neither one-to-one nor onto. Nevertheless, it provides an crucial mechanism for passing information between the group and the Lie algebra. The following result says that the exponential mapping is locally one-to-one and onto, a result that will be essential later.
###### Theorem 3.23.
Let $`G`$ be a matrix Lie group with Lie algebra $`𝔤`$. Then there exist a neighborhood $`U`$ of zero in $`𝔤`$ and a neighborhood $`V`$ of $`I`$ in $`G`$ such that the exponential mapping takes $`U`$ homeomorphically onto $`V`$.
###### Proof.
We follow the proof of Theorem I.3.11 in Bröcker and tom Dieck. In view of what we have proved about the matrix logarithm, we know this result for the case of $`\mathrm{𝖦𝖫}(n;)`$. To prove the general case, we consider a matrix Lie group $`G<\mathrm{𝖦𝖫}(n;)`$, with Lie algebra $`𝔤`$.
###### Lemma 3.24.
Suppose $`g_n`$ are elements of $`G`$, and that $`g_nI`$. Let $`Y_n=\mathrm{log}g_n`$, which is defined for all sufficiently large $`n`$. Suppose $`Y_n/Y_nY\mathrm{𝗀𝗅}(n;)`$. Then $`Y𝔤`$.
###### Proof.
To show that $`Y𝔤`$, we must show that $`\mathrm{exp}tYG`$ for all $`t`$. As $`n\mathrm{}`$, $`\left(t/Y_n\right)Y_ntY`$. Note that since $`g_nI`$, $`Y_n0`$, and so $`Y_n0`$. Thus we can find integers $`m_n`$ such that $`\left(m_nY_n\right)t`$. Then $`\mathrm{exp}\left(m_nY_n\right)=\mathrm{exp}\left[\left(m_nY_n\right)\left(Y_n/Y_n\right)\right]\mathrm{exp}\left(tY\right)`$. But $`\mathrm{exp}\left(m_nY_n\right)=\mathrm{exp}\left(Y_n\right)^{m_n}=\left(g_n\right)^{m_n}G`$, and $`G`$ is closed, so $`\mathrm{exp}\left(tY\right)G`$. ∎
We think of $`\mathrm{𝗀𝗅}(n;)`$ as $`^{n^2}^{2n^2}`$. Then $`𝔤`$ is a subspace of $`^{2n^2}`$. Let $`D`$ denote the orthogonal complement of $`𝔤`$ with respect to the usual inner product on $`^{2n^2}`$. Consider the map $`\mathrm{\Phi }:𝔤D\mathrm{𝖦𝖫}(n;)`$ given by
$$\mathrm{\Phi }(X,Y)=e^Xe^Y\text{.}$$
Of course, we can identify $`𝔤D`$ with $`^{2n^2}`$. Moreover, $`\mathrm{𝖦𝖫}(n;)`$ is an open subset of $`\mathrm{𝗀𝗅}(n;)^{2n^2}`$. Thus we can regard $`\mathrm{\Phi }`$ as a map from $`^{2n^2}`$ to itself.
Now, using the properties of the matrix exponential, we see that
$`{\displaystyle \frac{d}{dt}}|_{t=0}\mathrm{\Phi }(tX,0)`$ $`=X`$
$`{\displaystyle \frac{d}{dt}}|_{t=0}\mathrm{\Phi }(0,tY)`$ $`=Y\text{.}`$
This shows that the derivative of $`\mathrm{\Phi }`$ at the point $`0^{2n^2}`$ is the identity. (Recall that the derivative at a point of a function from $`^{2n^2}`$ to itself is a linear map of $`^{2n^2}`$ to itself, in this case the identity map.) In particular, the derivative of $`\mathrm{\Phi }`$ at 0 is invertible. Thus the inverse function theorem says that $`\mathrm{\Phi }`$ has a continuous local inverse, defined in a neighborhood of $`I`$.
Now let $`U`$ be any neighborhood of zero in $`𝔤`$. I want to show that $`\mathrm{exp}\left(U\right)`$ contains a neighborhood of $`I`$ in $`G`$. Suppose not. Then we can find a sequence $`g_nG`$ with $`g_nI`$ such that no $`g_n`$ is in $`\mathrm{exp}\left(U\right)`$. Since $`\mathrm{\Phi }`$ is locally invertible, we can write $`g_n`$ (for large $`n`$) uniquely as $`g_n=\mathrm{exp}\left(X_n\right)\mathrm{exp}\left(Y_n\right)`$, with $`X_n𝔤`$ and $`Y_nD`$. Since $`g_nI`$ and $`\mathrm{\Phi }^1`$ is continuous, $`X_n`$ and $`Y_n`$ tend to zero. Thus (for large $`n`$), $`X_nU`$. So we must have (for large $`n`$) $`Y_n0`$, otherwise $`g_n`$ would be in $`\mathrm{exp}\left(U\right)`$.
Let $`\stackrel{~}{g}_n=\mathrm{exp}\left(Y_n\right)=\mathrm{exp}\left(X_n\right)g_n`$. Note that $`\stackrel{~}{g}_nG`$ and $`\stackrel{~}{g}_nI`$. Since the unit ball in $`D`$ is compact, we can choose a subsequence of $`\left\{Y_n\right\}`$ (still called $`\left\{Y_n\right\}`$) so that $`Y_n/Y_n`$ converges to some $`YD`$, with $`Y=1`$. But then by the Lemma, $`Y𝔤`$! This is a contradiction, because $`D`$ is the orthogonal complement of $`𝔤`$.
So for every neighborhood $`U`$ of zero in $`𝔤`$, $`\mathrm{exp}\left(U\right)`$ contains a neighborhood of the identity in $`G`$. If we make $`U`$ small enough, then the exponential will be one-to-one on $`\overline{U}`$. (The existence of the matrix logarithm implies that the exponential is one-to-one near zero.) Let $`\mathrm{log}`$ denote the inverse map, defined on $`\mathrm{exp}\left(\overline{U}\right)`$. Since $`\overline{U}`$ is compact, and $`\mathrm{exp}`$ is one-to-one and continuous on $`\overline{U}`$, log will be continuous. (This is a standard topological result.) So take $`V`$ to be a neighborhood of $`I`$ contained in $`\mathrm{exp}\left(\overline{U}\right)`$, and let $`U^{}=\mathrm{exp}^1\left(V\right)U`$. Then $`U^{}`$ is open and the exponential takes $`U^{}`$ homeomorphically onto $`V`$. ∎
###### Definition 3.25.
If $`U`$ and $`V`$ are as in Proposition 3.23, then the inverse map $`\mathrm{exp}^1:V𝔤`$ is called the logarithm for $`G`$.
###### Corollary 3.26.
If $`G`$ is a connected matrix Lie group, then every element $`A`$ of $`G`$ can be written in the form
(3.16)
$$A=e^{X_1}e^{X_2}\mathrm{}e^{X_n}$$
for some $`X_1,X_2,\mathrm{}X_n`$ in $`𝔤`$.
###### Proof.
Recall that for us, saying $`G`$ is connected means that $`G`$ is path-connected. This certainly means that $`G`$ is connected in the usual topological sense, namely, the only non-empty subset of $`G`$ that is both open and closed is $`G`$ itself. So let $`E`$ denote the set of all $`AG`$ that can be written in the form (3.16). In light of the Proposition, $`E`$ contains a neighborhood $`V`$ of the identity. In particular, $`E`$ is non-empty.
We first claim that $`E`$ is open. To see this, consider $`AE`$. Then look at the set of matrices of the form $`AB`$, with $`BV`$. This will be a neighborhood of $`A`$. But every such $`B`$ can be written as $`B=e^X`$ and $`A`$ can be written as $`A=e^{X_1}e^{X_2}\mathrm{}e^{X_n}`$, so $`AB=e^{X_1}e^{X_2}\mathrm{}e^{X_n}e^X`$.
Now we claim that $`E`$ is closed (in $`G`$). Suppose $`AG`$, and there is a sequence $`A_nE`$ with $`A_nA`$. Then $`AA_n^1I`$. Thus we can choose some $`n_0`$ such that $`AA_{n_0}^1V`$. Then $`AA_{n_0}^1=e^X`$ and $`A=A_{n_0}e^X`$. But by assumption, $`A_{n_0}=e^{X_1}e^{X_2}\mathrm{}e^{X_n}`$, so $`A=e^{X_1}e^{X_2}\mathrm{}e^{X_n}e^X`$. Thus $`AE`$, and $`E`$ is closed.
Thus $`E`$ is both open and closed, so $`E=G`$. ∎
### 3.8. Lie Algebras
###### Definition 3.27.
A finite-dimensional real or complex Lie algebra is a finite-dimensional real or complex vector space $`𝔤`$, together with a map $`\left[\right]`$ from $`𝔤\times 𝔤`$ into $`𝔤`$, with the following properties:
1. $`\left[\right]`$ is bilinear.
2. $`[X,Y]=[Y,X]`$ for all $`X,Y𝔤`$.
3. $`[X,[Y,Z]]+[Y,[Z,X]]+[Z,[X,Y]]=0`$ for all $`X,Y,Z𝔤`$.
Condition 3 is called the Jacobi identity. Note also that Condition 2 implies that $`[X,X]=0`$ for all $`X𝔤`$. The same three conditions define a Lie algebra over an arbitrary field $`𝐅`$, except that if $`𝐅`$ has characteristic two, then one should add the condition $`[X,X]=0`$, which doesn’t follow from skew-symmetry in characteristic two. We will deal only with finite-dimensional Lie algebras, and will from now on interpret “Lie algebra” as “finite-dimensional Lie algebra.”
A Lie algebra is in fact an algebra in the usual sense, but the product operation $`\left[\right]`$ for this algebra is neither commutative nor associative. The Jacobi identity should be thought of as a substitute for associativity.
###### Proposition 3.28.
The space $`\mathrm{𝗀𝗅}(n;)`$ of all $`n\times n`$ real matrices is a real Lie algebra with respect to the bracket operation $`[A,B]=ABBA`$. The space $`\mathrm{𝗀𝗅}(n;)`$ of all $`n\times n`$ complex matrices is a complex Lie algebra with respect to the analogous bracket operation.
Let $`V`$ is a finite-dimensional real or complex vector space, and let $`\mathrm{𝗀𝗅}(V)`$ denote the space of linear maps of $`V`$ into itself. Then $`\mathrm{𝗀𝗅}(V)`$ becomes a real or complex Lie algebra with the bracket operation $`[A,B]=ABBA`$.
###### Proof.
The only non-trivial point is the Jacobi identity. The only way to prove this is to write everything out and see, and this is best left to the reader. Note that each triple bracket generates four terms, for a total of twelve. Each of the six orderings of $`\{X,Y,Z\}`$ occurs twice, once with a plus sign and once with a minus sign. ∎
###### Definition 3.29.
A subalgebra of a real or complex Lie algebra $`𝔤`$ is a subspace $`𝔥`$ of $`𝔤`$ such that $`[H_1,H_2]𝔥`$ for all $`H_1,H_2𝔥`$. If $`𝔤`$ is a complex Lie algebra, and $`𝔥`$ is a real subspace of $`𝔤`$ which is closed under brackets, then $`𝔥`$ is said to be a real subalgebra of $`𝔤`$.
If $`𝔤`$ and $`𝔥`$ are Lie algebras, then a linear map $`\varphi :𝔤𝔥`$ is called a Lie algebra homomorphism if $`\varphi \left([X,Y]\right)=[\varphi (X),\varphi (Y)]`$ for all $`X,Y𝔤`$. If in addition $`\varphi `$ is one-to-one and onto, then $`\varphi `$ is called a Lie algebra isomorphism. A Lie algebra isomorphism of a Lie algebra with itself is called a Lie algebra automorphism.
A subalgebra of a Lie algebra is again a Lie algebra. A real subalgebra of a complex Lie algebra is a real Lie algebra. The inverse of a Lie algebra isomorphism is again a Lie algebra isomorphism.
###### Proposition 3.30.
The Lie algebra $`𝔤`$ of a matrix Lie group $`G`$ is a real Lie algebra.
###### Proof.
By Theorem 3.16, $`𝔤`$ is a real subalgebra of $`\mathrm{𝗀𝗅}(n;)`$ complex matrices, and is thus a real Lie algebra. ∎
###### Theorem 3.31 (Ado).
Every finite-dimensional real Lie algebra is isomorphic to a subalgebra of $`\mathrm{𝗀𝗅}(n;)`$. Every finite-dimensional complex Lie algebra is isomorphic to a (complex) subalgebra of $`\mathrm{𝗀𝗅}(n;)`$.
This remarkable theorem is proved in Varadarajan. The proof is well beyond the scope of this course (which is after all a course on Lie groups), and requires a deep understanding of the structure of complex Lie algebras. The theorem tells us that every Lie algebra is (isomorphic to) a Lie algebra of matrices. (This is in contrast to the situation for Lie groups, where most but not all Lie groups are matrix Lie groups.)
###### Definition 3.32.
Let $`𝔤`$ be a Lie algebra. For $`X𝔤`$, define a linear map $`\mathrm{ad}X:𝔤𝔤`$ by
$$\mathrm{ad}X(Y)=[X,Y]\text{.}$$
Thus “$`\mathrm{ad}`$” (i.e., the map $`X\mathrm{ad}X`$) can be viewed as a linear map from $`𝔤`$ into $`\mathrm{𝗀𝗅}(𝔤)`$, where $`\mathrm{𝗀𝗅}(𝔤)`$ denotes the space of linear operators from $`𝔤`$ to $`𝔤`$.
Since $`\mathrm{ad}X(Y)`$ is just $`[X,Y]`$, it might seem foolish to introduce the additional “$`\mathrm{ad}`$” notation. However, thinking of $`[X,Y]`$ as a linear map in $`Y`$ for each fixed $`X`$, gives a somewhat different perspective. In any case, the “$`\mathrm{ad}`$” notation is extremely useful in some situations. For example, instead of writing
$$[X,[X,[X,[X,Y]]]]$$
we can now write
$$\left(\mathrm{ad}X\right)^4(Y)\text{.}$$
This kind of notation will be essential in Section 4.1.
###### Proposition 3.33.
If $`𝔤`$ is a Lie algebra, then
$$\mathrm{ad}[X,Y]=\mathrm{ad}X\mathrm{ad}Y\mathrm{ad}Y\mathrm{ad}X=[\mathrm{ad}X,\mathrm{ad}Y]\text{.}$$
That is, ad$`:𝔤\mathrm{𝗀𝗅}(𝔤)`$ is a Lie algebra homomorphism.
###### Proof.
Observe that
$$\mathrm{ad}[X,Y](Z)=[[X,Y],Z]$$
whereas
$$[\mathrm{ad}X,\mathrm{ad}Y](Z)=[X,[Y,Z]][Y,[X,Z]]\text{.}$$
So we require that
$$[[X,Y],Z]=[X,[Y,Z]][Y,[X,Z]]$$
or equivalently
$$0=[X,[Y,Z]]+[Y,[Z,X]]+[Z,[X,Y]]$$
which is exactly the Jacobi identity. ∎
Recall that for any $`X𝔤`$, and any $`AG`$, we define
$$\mathrm{Ad}A(X)=AXA^1$$
and that Ad$`:G\mathrm{𝖦𝖫}(𝔤)`$ is a Lie group homomorphism. We showed (Proposition 3.21) that the associated Lie algebra homomorphism $`\stackrel{~}{\mathrm{Ad}}:𝔤\mathrm{𝗀𝗅}(𝔤)`$ is given by
$$\stackrel{~}{\mathrm{Ad}}X(Y)=[X,Y]\text{.}$$
In our new notation, we may say
$$\stackrel{~}{\mathrm{Ad}}=\mathrm{ad}$$
By the defining property of $`\stackrel{~}{\mathrm{Ad}}`$, we have the following identity: For all $`X𝔤`$,
(3.17)
$$\mathrm{Ad}(e^X)=e^{\mathrm{ad}X}\text{.}$$
Note that both sides of (3.17) are linear operators on the Lie algebra $`𝔤`$. This is an important relation, which can also be verified directly, by expanding out both sides. (See Exercise 13.)
#### 3.8.1. Structure Constants
Let $`𝔤`$ be a finite-dimensional real or complex Lie algebra, and let $`X_1,\mathrm{},X_n`$ be a basis for $`𝔤`$ (as a vector space). Then for each $`i,j`$, $`[X_i,X_j]`$ can be written uniquely in the form
$$[X_i,X_j]=\underset{k=1}{\overset{n}{}}c_{ijk}X_k\text{.}$$
The constants $`c_{ijk}`$ are called the structure constants of $`𝔤`$ (with respect to the chosen basis). Clearly, the structure constants determine the bracket operation on $`𝔤`$. In some of the literature, the structure constants play an important role, although we will not have occasion to use them in this course. (In the physics literature, the structure constants are defined as $`[X_i,X_j]=\sqrt{1}_kc_{ijk}X_k`$, reflecting the factor of $`\sqrt{1}`$ difference between the physics definition of the Lie algebra and our own.)
The structure constants satisfy the following two conditions,
$`c_{ijk}+c_{jik}`$ $`=0`$
$`{\displaystyle \underset{m}{}}(c_{ijm}c_{mkl}+c_{jkm}c_{mil}+c_{kim}c_{mjl})`$ $`=0`$
for all $`i,j,k,l`$. The first of these conditions comes from the skew-symmetry of the bracket, and the second comes from the Jacobi identity. (The reader is invited to verify these conditions for himself.)
### 3.9. The Complexification of a Real Lie Algebra
###### Definition 3.34.
If $`V`$ is a finite-dimensional real vector space, then the complexification of $`V`$, denoted $`V_{}`$, is the space of formal linear combinations
$$v_1+iv_2$$
with $`v_1,v_2V`$. This becomes a real vector space in the obvious way, and becomes a complex vector space if we define
$$i(v_1+iv_2)=v_2+iv_1\text{.}$$
We could more pedantically define $`V_{}`$ to be the space of ordered pairs $`(v_1,v_2)`$, but this is notationally cumbersome. It is straightforward to verify that the above definition really makes $`V_{}`$ into a complex vector space. We will regard $`V`$ as a real subspace of $`V_{}`$ in the obvious way.
###### Proposition 3.35.
Let $`𝔤`$ be a finite-dimensional real Lie algebra, and $`𝔤_{}`$ its complexification (as a real vector space). Then the bracket operation on $`𝔤`$ has a unique extension to $`𝔤_{}`$ which makes $`𝔤_{}`$ into a complex Lie algebra. The complex Lie algebra $`𝔤_{}`$ is called the complexification of the real Lie algebra $`𝔤`$.
###### Proof.
The uniqueness of the extension is obvious, since if the bracket operation on $`𝔤_{}`$ is to be bilinear, then it must be given by
(3.18)
$$[X_1+iX_2,Y_1+iY_2]=\left([X_1,Y_1][X_2,Y_2]\right)+i\left([X_1,Y_2]+[X_2,Y_1]\right)\text{.}$$
To show existence, we must now check that (3.18) is really bilinear and skew-symmetric, and that it satisfies the Jacobi identity. It is clear that (3.18) is real bilinear, and skew-symmetric. The skew-symmetry means that if (3.18) is complex linear in the first factor, it is also complex linear in the second factor. Thus we need only show that
(3.19)
$$[i(X_1+iX_2),Y_1+iY_2]=i[X_1+iX_2,Y_1+iY_2]\text{.}$$
Well, the left side of (3.19) is
$$[X_2+iX_1,Y_1+iY_2]=\left([X_2,Y_1][X_1,Y_2]\right)+i\left([X_1,Y_1][X_2,Y_2]\right)$$
whereas the right side of (3.19) is
$`i\left\{\left([X_1,Y_1][X_2,Y_2]\right)+i\left([X_2,Y_1]+[X_1,Y_2]\right)\right\}`$
$`=\left([X_2,Y_1][X_1,Y_2]\right)+i\left([X_1,Y_1][X_2,Y_2]\right)\text{,}`$
and indeed these are equal.
It remains to check the Jacobi identity. Of course, the Jacobi identity holds if $`X,Y,`$ and $`Z`$ are in $`𝔤`$. But now observe that the expression on the left side of the Jacobi identity is (complex!) linear in $`X`$ for fixed $`Y`$ and $`Z`$. It follows that the Jacobi identity holds if $`X`$ is in $`𝔤_{}`$, and $`Y,Z`$ in $`𝔤`$. The same argument then shows that we can extend to $`Y`$ in $`𝔤_{}`$, and then to $`Z`$ in $`𝔤_{}`$. Thus the Jacobi identity holds in $`𝔤_{}`$. ∎
###### Proposition 3.36.
The Lie algebras $`\mathrm{𝗀𝗅}(n;)`$, $`\mathrm{𝗌𝗅}(n;)`$, $`\mathrm{𝗌𝗈}(n;)`$, and $`\mathrm{𝗌𝗉}(n;)`$ are complex Lie algebras, as is the Lie algebra of the complex Heisenberg group. In addition, we have the following isomorphisms of complex Lie algebras
$$\begin{array}{ccc}\mathrm{𝗀𝗅}(n;)_{}& & \mathrm{𝗀𝗅}(n;)\\ 𝗎(n)_{}& & \mathrm{𝗀𝗅}(n;)\\ \mathrm{𝗌𝗅}(n;)_{}& & \mathrm{𝗌𝗅}(n;)\\ \mathrm{𝗌𝗈}(n)_{}& & \mathrm{𝗌𝗈}(n;)\\ \mathrm{𝗌𝗉}(n;)_{}& & \mathrm{𝗌𝗉}(n;)\\ \mathrm{𝗌𝗉}(n)_{}& & \mathrm{𝗌𝗉}(n;)\text{.}\end{array}$$
###### Proof.
From the computations in the previous section we see easily that the specified Lie algebras are in fact complex subalgebras of $`\mathrm{𝗀𝗅}(n;)`$, and hence are complex Lie algebras.
Now, $`\mathrm{𝗀𝗅}(n;)`$ is the space of all $`n\times n`$ complex matrices, whereas $`\mathrm{𝗀𝗅}(n;)`$ is the space of all $`n\times n`$ real matrices. Clearly, then, every $`X\mathrm{𝗀𝗅}(n;)`$ can be written uniquely in the form $`X_1+iX_2`$, with $`X_1,X_2\mathrm{𝗀𝗅}(n;)`$. This gives us a complex vector space isomorphism of $`\mathrm{𝗀𝗅}(n;)_{}`$ with $`\mathrm{𝗀𝗅}(n;)`$, and it is a triviality to check that this is a Lie algebra isomorphism.
On the other hand, $`𝗎(n)`$ is the space of all $`n\times n`$ complex skew-self-adjoint matrices. But if $`X`$ is any $`n\times n`$ complex matrix, then
$`X`$ $`={\displaystyle \frac{XX^{}}{2}}+{\displaystyle \frac{X+X^{}}{2}}`$
$`={\displaystyle \frac{XX^{}}{2}}+i{\displaystyle \frac{(iX)(iX)^{}}{2}}\text{.}`$
Thus $`X`$ can be written as a skew matrix plus $`i`$ times a skew matrix, and it is easy to see that this decomposition is unique. Thus every $`X`$ in $`\mathrm{𝗀𝗅}(n;)`$ can be written uniquely as $`X_1+iX_2`$, with $`X_1`$ and $`X_2`$ in $`𝗎(n)`$. It follows that $`𝗎(n)_{}\mathrm{𝗀𝗅}(n;)`$.
The verification of the remaining isomorphisms is similar, and is left as an exercise to the reader. ∎
Note that $`𝗎(n)_{}\mathrm{𝗀𝗅}(n;)_{}`$ $`\mathrm{𝗀𝗅}(n;)`$. However, $`𝗎(n)`$ is not isomorphic to $`\mathrm{𝗀𝗅}(n;)`$, except when $`n=1`$. The real Lie algebras $`𝗎(n)`$ and $`\mathrm{𝗀𝗅}(n;)`$ are called real forms of the complex Lie algebra $`\mathrm{𝗀𝗅}(n;)`$. A given complex Lie algebra may have several non-isomorphic real forms. See Exercise 11.
Physicists do not always clearly distinguish between a matrix Lie group and its (real) Lie algebra, or between a real Lie algebra and its complexification. Thus, for example, some references in the physics literature to SU$`(2)`$ actually refer to the complexified Lie algebra, $`\mathrm{𝗌𝗅}(2;)`$.
### 3.10. Exercises
1. The product rule*.* Recall that a matrix-valued function $`A(t)`$ is smooth if each $`A_{ij}(t)`$ is smooth. The derivative of such a function is defined as
$$\left(\frac{dA}{dt}\right)_{ij}=\frac{dA_{ij}}{dt}$$
or equivalently,
$$\frac{d}{dt}A(t)=\underset{h0}{lim}\frac{A(t+h)A(t)}{h}\text{.}$$
Let $`A(t)`$ and $`B(t)`$ be two such functions. Prove that $`A(t)B(t)`$ is again smooth, and that
$$\frac{d}{dt}\left[A(t)B(t)\right]=\frac{dA}{dt}B(t)+A(t)\frac{dB}{dt}\text{.}$$
2. Using the Jordan canonical form, show that every $`n\times n`$ matrix $`A`$ can be written as $`A=S+N`$, with $`S`$ diagonalizable (over $``$), $`N`$ nilpotent, and $`SN=NS`$. Recall that the Jordan canonical form is block diagonal, with each block of the form
$$\left(\begin{array}{ccc}\lambda & & \\ & \mathrm{}& \\ 0& & \lambda \end{array}\right)\text{.}$$
3. Let $`X`$ and $`Y`$ be $`n\times n`$ matrices. Show that there exists a constant $`C`$ such that
$$e^{(X+Y)/m}e^{X/m}e^{Y/m}\frac{C}{m^2}$$
for all integers $`m1`$.
4. Using the Jordan canonical form, show that every $`n\times n`$ complex matrix $`A`$ is the limit of a sequence of diagonalizable matrices.
Hint: If the characteristic polynomial of $`A`$ has $`n`$ distinct roots, then $`A`$ is diagonalizable.
5. Give an example of a matrix Lie group $`G`$ and a matrix $`X`$ such that $`e^XG`$, but $`X𝔤`$.
6. Show that two isomorphic matrix Lie groups have isomorphic Lie algebras.
7. The Lie algebraso$`(3;1)`$. Write out explicitly the general form of a $`4\times 4`$ real matrix in so$`(3;1)`$.
8. Verify directly that Proposition 3.15 and Theorem 3.16 hold for the Lie algebra of $`\mathrm{𝖲𝖴}(n)`$.
9. The Lie algebra$`\mathrm{𝗌𝗎}(2)`$. Show that the following matrices form a basis for the real Lie algebra $`\mathrm{𝗌𝗎}(2)`$:
$$\begin{array}{ccc}E_1=\frac{1}{2}\left(\begin{array}{cc}i& 0\\ 0& i\end{array}\right)& E_2=\frac{1}{2}\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)& E_3=\frac{1}{2}\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right)\end{array}\text{.}$$
Compute $`[E_1,E_2]`$, $`[E_2,E_3]`$, and $`[E_3,E_1]`$. Show that there is an invertible linear map $`\varphi :\mathrm{𝗌𝗎}(2)^3`$ such that $`\varphi ([X,Y])=\varphi (X)\times \varphi (Y)`$ for all $`X,Y`$ $`\mathrm{𝗌𝗎}(2),`$ where $`\times `$ denotes the cross-product on $`^3`$.
10. The Lie algebras $`\mathrm{𝗌𝗎}(2)`$ and$`\mathrm{𝗌𝗈}(3)`$. Show that the real Lie algebras $`\mathrm{𝗌𝗎}(2)`$ and $`\mathrm{𝗌𝗈}(3)`$ are isomorphic.
Note: Nevertheless, the corresponding groups $`\mathrm{𝖲𝖴}(2)`$ and $`\mathrm{𝖲𝖮}(3)`$ are not isomorphic. (Although $`\mathrm{𝖲𝖮}(3)`$ is isomorphic to $`\mathrm{𝖲𝖴}(2)/\{I,I\}`$.)
11. The Lie algebras$`\mathrm{𝗌𝗎}(2)`$ *and* $`\mathrm{𝗌𝗅}(2;)`$. Show that $`\mathrm{𝗌𝗎}(2)`$and $`\mathrm{𝗌𝗅}(2;)`$ are not isomorphic Lie algebras, even though $`\mathrm{𝗌𝗎}(2)_{}\mathrm{𝗌𝗅}(2;)_{}`$.
Hint: Using Exercise 9, show that $`\mathrm{𝗌𝗎}(2)`$ has no two-dimensional subalgebras.
12. Let $`G`$ be a matrix Lie group, and $`𝔤`$ its Lie algebra. For each $`AG`$, show that $`\mathrm{Ad}A`$ is a Lie algebra automorphism of $`𝔤`$.
13. Ad and ad. Let $`X`$ and $`Y`$ be matrices. Show by induction that
$$\left(\mathrm{ad}X\right)^n(Y)=\underset{k=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{k}\right)X^kY(X)^{nk}\text{.}$$
Now show by direct computation that
$$e^{\mathrm{ad}X}(Y)=\mathrm{Ad}(e^X)Y=e^XYe^X\text{.}$$
You may assume that it is legal to multiply power series term-by-term. (This result was obtained indirectly in Equation 3.17.)
Hint: Recall that Pascal’s Triangle gives a relationship between things of the form $`\left(\genfrac{}{}{0pt}{}{n+1}{k}\right)`$ and things of the form $`\left(\genfrac{}{}{0pt}{}{n}{k}\right)`$.
14. The complexification of a real Lie algebra. Let $`𝔤`$ be a real Lie algebra, $`𝔤_{}`$ its complexification, and $`𝔥`$ an arbitrary complex Lie algebra. Show that every real Lie algebra homomorphism of $`𝔤`$ into $`𝔥`$ extends uniquely to a complex Lie algebra homomorphism of $`𝔤_{}`$ into $`𝔥`$. (This is the universal property of the complexification of a real Lie algebra. This property can be used as an alternative definition of the complexification.)
15. The exponential mapping for$`\mathrm{𝖲𝖫}(2;)`$. Show that the image of the exponential mapping for $`\mathrm{𝖲𝖫}(2;)`$ consists of precisely those matrices $`A\mathrm{𝖲𝖫}(2;)`$ such that $`\mathrm{trace}\left(A\right)>2,`$ together with the matrix $`I`$ (which has trace $`2`$). You will need to consider the possibilities for the eigenvalues of a matrix in the Lie algebra $`\mathrm{𝗌𝗅}(2;)`$ and in the group $`\mathrm{𝖲𝖫}(2;)`$. In the Lie algebra, show that the eigenvalues are of the form $`(\lambda ,\lambda )`$ or $`(i\lambda ,i\lambda )`$ with $`\lambda `$ real. In the group, show that the eigenvalues are of the form $`(\alpha ,1/a)`$ or $`(a,1/a)`$ with $`a`$ real and positive, or else of the form $`(e^{i\theta },e^{i\theta }),`$ with $`\theta `$ real. The case of a repeated eigenvalue ($`(0,0)`$ in the Lie algebra and $`(1,1)`$ or $`(1,1)`$ in the group) will have to be treated separately.
Show that the image of the exponential mapping is not dense in $`\mathrm{𝖲𝖫}(2;)`$.
16. Using Exercise 4, show that the exponential mapping for $`\mathrm{𝖦𝖫}(n;)`$ maps onto a dense subset of $`\mathrm{𝖦𝖫}(n;)`$.
17. The exponential mapping for the Heisenberg group. Show that the exponential mapping from the Lie algebra of the Heisenberg group to the Heisenberg group is one-to-one and onto.
18. The exponential mapping for$`𝖴(n)`$. Show that the exponential mapping from $`𝗎(n)`$ to $`𝖴(n)`$ is onto, but not one-to-one. (Note that this shows that $`𝖴(n)`$ is connected.)
Hint: Every unitary matrix has an orthonormal basis of eigenvectors.
19. Let $`G`$ be a matrix Lie group, and $`𝔤`$ its Lie algebra. Let $`A(t)`$ be a smooth curve lying in $`G`$, with $`A(0)=I`$. Let $`X=\frac{d}{dt}|_{t=0}A(t)`$. Show that $`X𝔤`$.
Hint: Use Proposition 3.8.
Note: This shows that the Lie algebra $`𝔤`$ coincides with what would be called the tangent space at the identity in the language of differentiable manifolds.
20. Consider the space $`\mathrm{𝗀𝗅}(n;)`$ of all $`n\times n`$ complex matrices. As usual, for $`X\mathrm{𝗀𝗅}(n;)`$, define $`\mathrm{ad}X:\mathrm{𝗀𝗅}(n;)\mathrm{𝗀𝗅}(n;)`$ by $`\mathrm{ad}X(Y)=[X,Y]`$. Suppose that $`X`$ is a diagonalizable matrix. Show, then, that $`\mathrm{ad}X`$ is diagonalizable as an operator on $`\mathrm{𝗀𝗅}(n;)`$.
Hint: Consider first the case where $`X`$ is actually diagonal.
Note: The problem of diagonalizing $`\mathrm{ad}X`$ is an important one that we will encounter again in Chapter 6, when we consider semisimple Lie algebras.
## Chapter 4 The Baker-Campbell-Hausdorff Formula
### 4.1. The Baker-Campbell-Hausdorff Formula for the Heisenberg Group
A crucial result of Chapter 5 will be the following: Let $`G`$ and $`H`$ be matrix Lie groups, with Lie algebras $`𝔤`$ and $`𝔥`$, and suppose that $`G`$ is connected and simply connected. Then if $`\stackrel{~}{\varphi }:𝔤𝔥`$ is a Lie algebra homomorphism, there exists a unique Lie group homomorphism $`\varphi :GH`$ such that $`\varphi `$ and $`\stackrel{~}{\varphi }`$ are related as in Theorem 3.18. This result is extremely important because it implies that if $`G`$ is connected and simply connected, then there is a natural one-to-one correspondence between the representations of $`G`$ and the representations of its Lie algebra $`𝔤`$ (as explained in Chapter 5). In practice, it is much easier to determine the representations of the Lie algebra than to determine directly the representations of the corresponding group.
This result (relating Lie algebra homomorphisms and Lie group homomorphisms) is deep. The “modern” proof (e.g., Varadarajan, Theorem 2.7.5) makes use of the Frobenius theorem, which is both hard to understand and hard to prove (Varadarajan, Section 1.3). Our proof will instead use the Baker-Campbell-Hausdorff formula, which is more easily stated and more easily motivated than the Frobenius theorem, but still deep.
The idea is the following. The desired group homomorphism $`\varphi :GH`$ must satisfy
(4.1)
$$\varphi \left(e^X\right)=e^{\stackrel{~}{\varphi }(X)}\text{.}$$
We would like, then, to define $`\varphi `$ by this relation. This approach has two serious difficulties. First, a given element of $`G`$ may not be expressible as $`e^X`$, and even if it is, the $`X`$ may not be unique. Second, it is very far from clear why the $`\varphi `$ in (4.1) (even to the extent it is well-defined) should be a group homomorphism.
It is the second issue which the Baker-Campbell-Hausdorff formula addresses. (The first issue will be addressed in the next chapter; it is there that the simple connectedness of $`G`$ comes into play.) Specifically, (one form of) the Baker-Campbell-Hausdorff formula says that if $`X`$ and $`Y`$ are sufficiently small, then
(4.2)
$$\mathrm{log}(e^Xe^Y)=X+Y+\frac{1}{2}[X,Y]+\frac{1}{12}[X,[X,Y]]\frac{1}{12}[Y,[X,Y]]+\mathrm{}\text{.}$$
It is not supposed to be evident at the moment what “$`\mathrm{}`$” refers to. The only important point is that all of the terms in (4.2) are given in terms of $`X`$ and $`Y`$, brackets of $`X`$ and $`Y`$, brackets of brackets involving $`X`$ and $`Y`$, etc. Then because $`\stackrel{~}{\varphi }`$ is a Lie algebra homomorphism,
$`\stackrel{~}{\varphi }\left(\mathrm{log}\left(e^Xe^Y\right)\right)`$ $`=\stackrel{~}{\varphi }(X)+\stackrel{~}{\varphi }(Y)+\frac{1}{2}[\stackrel{~}{\varphi }(X),\stackrel{~}{\varphi }(Y)]`$
$`+\frac{1}{12}[\stackrel{~}{\varphi }(X),[\stackrel{~}{\varphi }(X),\stackrel{~}{\varphi }(Y)]]\frac{1}{12}[\stackrel{~}{\varphi }(Y),[\stackrel{~}{\varphi }(X),\stackrel{~}{\varphi }(Y)]]+\mathrm{}`$
(4.3) $`=\mathrm{log}\left(e^{\stackrel{~}{\varphi }(X)}e^{\stackrel{~}{\varphi }(Y)}\right)`$
The relation (4.3) is extremely significant. For of course
$$e^Xe^Y=e^{\mathrm{log}(e^Xe^Y)}$$
and so by (4.1),
$$\varphi \left(e^Xe^Y\right)=e^{\stackrel{~}{\varphi }(\mathrm{log}(e^Xe^Y))}\text{.}$$
Thus (4.3) tells us that
$$\varphi \left(e^Xe^Y\right)=e^{\mathrm{log}\left(e^{\stackrel{~}{\varphi }(X)}e^{\stackrel{~}{\varphi }(Y)}\right)}=e^{\stackrel{~}{\varphi }(X)}e^{\stackrel{~}{\varphi }(Y)}=\varphi (e^X)\varphi (e^Y)\text{.}$$
Thus, the Baker-Campbell-Hausdorff formula shows that on elements of the form $`e^X`$, with $`X`$ small, $`\varphi `$ is a group homomorphism. (See Corollary 4.4 below.)
The Baker-Campbell-Hausdorff formula shows that all the information about the group product, at least near the identity, is “encoded” in the Lie algebra. Thus if $`\stackrel{~}{\varphi }`$ is a Lie algebra homomorphism (which by definition preserves the Lie algebra structure), and if we define $`\varphi `$ near the identity by (4.1), then we can expect $`\varphi `$ to preserve the group structure, i.e., to be a group homomorphism.
In this section we will look at how all of this works out in the very special case of the Heisenberg group. In the next section we will consider the general situation.
###### Theorem 4.1.
Suppose $`X`$ and $`Y`$ are $`n\times n`$ complex matrices, and that $`X`$ and $`Y`$ commute with their commutator. That is, suppose that
$$[X,[X,Y]]=[Y,[X,Y]]=0\text{.}$$
Then
$$e^Xe^Y=e^{X+Y+\frac{1}{2}[X,Y]}\text{.}$$
This is the special case of (4.2) in which the series terminates after the $`[X,Y]`$ term.
###### Proof.
Let $`X`$ and $`Y`$ be as in the statement of the theorem. We will prove that in fact
$$e^{tX}e^{tY}=\mathrm{exp}\left(tX+tY+\frac{t^2}{2}[X,Y]\right)\text{,}$$
which reduces to the desired result in the case $`t=1`$. Since by assumption $`[X,Y]`$ commutes with everything in sight, the above relation is equivalent to
(4.4)
$$e^{tX}e^{tY}e^{\frac{t^2}{2}[X,Y]}=e^{t\left(X+Y\right)}\text{.}$$
Let us call the left side of (4.4) $`A(t)`$ and the right side $`B\left(t\right)`$. Our strategy will be to show that $`A\left(t\right)`$ and $`B\left(t\right)`$ satisfy the same differential equation, with the same initial conditions. We can see right away that
$$\frac{dB}{dt}=B\left(t\right)\left(X+Y\right)\text{.}$$
On the other hand, differentiating $`A\left(t\right)`$ by means of the product rule gives
(4.5)
$$\frac{dA}{dt}=e^{tX}Xe^{tY}e^{\frac{t^2}{2}[X,Y]}+e^{tX}e^{tY}Ye^{\frac{t^2}{2}[X,Y]}+e^{tX}e^{tY}e^{\frac{t^2}{2}[X,Y]}\left(t[X,Y]\right)\text{.}$$
(You can verify that the last term on the right is correct by differentiating term-by-term.)
Now, since $`X`$ and $`Y`$ commute with $`[X,Y]`$, they also commute with $`e^{\frac{t^2}{2}[X,Y]}`$. Thus the second term on the right in (4.5) can be rewritten as
$$e^{tX}e^{tY}e^{\frac{t^2}{2}[X,Y]}Y\text{.}$$
The first term on the right in (4.5) is more complicated, since $`X`$ does not necessarily commute with $`e^{tY}`$. However,
$`Xe^{tY}=e^{tY}e^{tY}Xe^{tY}`$
$`=e^{tY}\mathrm{Ad}\left(e^{tY}\right)\left(X\right)`$
$`=e^{tY}e^{t\mathrm{ad}Y}\left(X\right)\text{.}`$
But since $`[Y,[Y,X]]=[Y,[X,Y]]=0`$,
$$e^{t\mathrm{ad}Y}\left(X\right)=Xt[Y,X]=X+t[X,Y]$$
with all higher terms being zero. Using the fact that everything commutes with $`e^{\frac{t^2}{2}[X,Y]}`$ gives
$$e^{tX}Xe^{tY}e^{\frac{t^2}{2}[X,Y]}=e^{tX}e^{tY}e^{\frac{t^2}{2}[X,Y]}\left(X+t[X,Y]\right)$$
Making these substitutions into (4.5) gives
$`{\displaystyle \frac{dA}{dt}}=e^{tX}e^{tY}e^{\frac{t^2}{2}[X,Y]}\left(X+t[X,Y]\right)+e^{tX}e^{tY}e^{\frac{t^2}{2}[X,Y]}Y+e^{tX}e^{tY}e^{\frac{t^2}{2}[X,Y]}\left(t[X,Y]\right)`$
$`=e^{tX}e^{tY}e^{\frac{t^2}{2}[X,Y]}\left(X+Y\right)`$
$`=A\left(t\right)\left(X+Y\right)\text{.}`$
Thus $`A\left(t\right)`$ and $`B\left(t\right)`$ satisfy the same differential equation. Moreover, $`A\left(0\right)=B\left(0\right)=I`$. Thus by standard uniqueness results for ordinary differential equations, $`A\left(t\right)=B\left(t\right)`$ for all $`t`$. ∎
###### Theorem 4.2.
Let $`H`$ denote the Heisenberg group, and $`𝔥`$ its Lie algebra. Let $`G`$ be a matrix Lie group with Lie algebra $`𝔤`$, and let $`\stackrel{~}{\varphi }:𝔥𝔤`$ be a Lie algebra homomorphism. Then there exists a unique Lie group homomorphism $`\varphi :HG`$ such that
$$\varphi \left(e^X\right)=e^{\stackrel{~}{\varphi }\left(X\right)}$$
for all $`X𝔥`$.
###### Proof.
Recall that the Heisenberg group has the very special property that its exponential mapping is one-to-one and onto. Let “log” denote the inverse of this map. Define $`\varphi :HG`$ by the formula
$$\varphi \left(A\right)=e^{\stackrel{~}{\varphi }\left(\mathrm{log}A\right)}\text{.}$$
We will show that $`\varphi `$ is a Lie group homomorphism.
If $`X`$ and $`Y`$ are in the Lie algebra of the Heisenberg group ($`3\times 3`$ strictly upper-triangular matrices), then $`[X,Y]`$ is of the form
$$\left(\begin{array}{ccc}0\hfill & 0\hfill & a\hfill \\ 0\hfill & 0\hfill & 0\hfill \\ 0\hfill & 0\hfill & 0\hfill \end{array}\right);$$
such a matrix commutes with both $`X`$ and $`Y`$. That is, $`X`$ and $`Y`$ commute with their commutator. Since $`\stackrel{~}{\varphi }`$ is a Lie algebra homomorphism, $`\stackrel{~}{\varphi }\left(X\right)`$ and $`\stackrel{~}{\varphi }\left(Y\right)`$ will also commute with their commutator:
$`[\stackrel{~}{\varphi }\left(X\right),[\stackrel{~}{\varphi }\left(X\right),\stackrel{~}{\varphi }\left(Y\right)]]=\stackrel{~}{\varphi }\left([X,[X,Y]]\right)=0`$
$`[\stackrel{~}{\varphi }\left(Y\right),[\stackrel{~}{\varphi }\left(X\right),\stackrel{~}{\varphi }\left(Y\right)]]=\stackrel{~}{\varphi }\left([Y,[X,Y]]\right)=0\text{.}`$
We want to show that $`\varphi `$ is a homomorphism, i.e., that $`\varphi \left(AB\right)=\varphi \left(A\right)\varphi \left(B\right)`$. Well, $`A`$ can be written as $`e^X`$ for a unique $`X𝔥`$ and $`B`$ can be written as $`e^Y`$ for a unique $`Y𝔥`$. Thus by Theorem 4.1
$$\varphi \left(AB\right)=\varphi \left(e^Xe^Y\right)=\varphi \left(e^{X+Y+\frac{1}{2}[X,Y]}\right)\text{.}$$
Using the definition of $`\varphi `$ and the fact that $`\stackrel{~}{\varphi }`$ is a Lie algebra homomorphism:
$$\varphi \left(AB\right)=\mathrm{exp}\left(\stackrel{~}{\varphi }\left(X\right)+\stackrel{~}{\varphi }\left(Y\right)+\frac{1}{2}[\stackrel{~}{\varphi }\left(X\right),\stackrel{~}{\varphi }\left(Y\right)]\right)\text{.}$$
Finally, using Theorem 4.1 again we have
$$\varphi \left(AB\right)=e^{\stackrel{~}{\varphi }\left(X\right)}e^{\stackrel{~}{\varphi }\left(Y\right)}=\varphi \left(A\right)\varphi \left(B\right)\text{.}$$
Thus $`\varphi `$ is a group homomorphism. It is easy to check that $`\varphi `$ is continuous (by checking that $`\mathrm{log}`$, exp, and $`\stackrel{~}{\varphi }`$ are all continuous), and so $`\varphi `$ is a Lie group homomorphism. Moreover, $`\varphi `$ by definition has the right relationship to $`\stackrel{~}{\varphi }`$. Furthermore, since the exponential mapping is one-to-one and onto, there can be at most one $`\varphi `$ with $`\varphi \left(e^X\right)=e^{\stackrel{~}{\varphi }\left(X\right)}`$. So we have uniqueness. ∎
### 4.2. The General Baker-Campbell-Hausdorff Formula
The importance of the Baker-Campbell-Hausdorff formula lies not in the details of the formula, but in the fact that there is one, and the fact that it gives $`\mathrm{log}(e^Xe^Y)`$ in terms of brackets of $`X`$ and $`Y`$, brackets of brackets, etc. This tells us something very important, namely that (at least for elements of the form $`e^X`$, $`X`$ small) the group product for a matrix Lie group $`G`$ is completely expressible in terms of the Lie algebra. (This is because $`\mathrm{log}\left(e^Xe^Y\right)`$, and hence also $`e^Xe^Y`$ itself, can be computed in Lie-algebraic terms by (4.2).)
We will actually state and prove an integral form of the Baker-Campbell-Hausdorff formula, rather than the series form (4.2). However, the integral form is sufficient to obtain the desired result (4.3). (See Corollary 4.4.) The series form of the Baker-Campbell-Hausdorff formula is stated precisely and proved in Varadarajan, Sec. 2.15.
Consider the function
$$g(z)=\frac{\mathrm{log}z}{1\frac{1}{z}}\text{.}$$
This function is defined and analytic in the disk $`\left\{\left|z1\right|<1\right\}`$, and thus for $`z`$ in this set, $`g(z)`$ can be expressed as
$$g(z)=\underset{m=0}{\overset{\mathrm{}}{}}a_m(z1)^m\text{.}$$
This series has radius of convergence one.
Now suppose $`V`$ is a finite-dimensional complex vector space. Choose an arbitrary basis for $`V`$, so that $`V`$ can be identified with $`^n`$ and thus the norm of a linear operator on $`V`$ can be defined. Then for any operator $`A`$ on $`V`$ with $`AI<1`$, we can define
$$g(A)=\underset{m=0}{\overset{\mathrm{}}{}}a_m(A1)^m\text{.}$$
We are now ready to state the integral form of the Baker-Campbell-Hausdorff formula.
###### Theorem 4.3 (Baker-Campbell-Hausdorff).
For all $`n\times n`$ complex matrices $`X`$ and $`Y`$ with $`X`$ and $`Y`$ sufficiently small,
(4.6)
$$\mathrm{log}\left(e^Xe^Y\right)=X+_0^1g(e^{\mathrm{ad}X}e^{t\mathrm{ad}Y})(Y)𝑑t\text{.}$$
###### Corollary 4.4.
Let $`G`$ be a matrix Lie group and $`𝔤`$ its Lie algebra. Suppose that $`\stackrel{~}{\varphi }:𝔤\mathrm{𝗀𝗅}(n;𝐂)`$ is a Lie algebra homomorphism. Then for all sufficiently small $`X,Y`$ in $`𝔤`$, $`\mathrm{log}\left(e^Xe^Y\right)`$ is in $`𝔤`$, and
(4.7)
$$\stackrel{~}{\varphi }\left[\mathrm{log}\left(e^Xe^Y\right)\right]=\mathrm{log}\left(e^{\stackrel{~}{\varphi }(X)}e^{\stackrel{~}{\varphi }(Y)}\right)\text{.}$$
Note that $`e^{\mathrm{ad}X}e^{t\mathrm{ad}Y}`$, and hence also $`g(e^{\mathrm{ad}X}e^{t\mathrm{ad}Y})`$, is a linear operator on the space $`\mathrm{𝗀𝗅}(n;)`$ of all $`n\times n`$ complex matrices. In (4.6), this operator is being applied to the matrix $`Y`$. The fact that $`X`$ and $`Y`$ are assumed small guarantees that $`e^{\mathrm{ad}X}e^{t\mathrm{ad}Y}`$ is close to the identity operator on $`\mathrm{𝗀𝗅}(n;)`$ for all $`0t1`$. This ensures that $`g(e^{\mathrm{ad}X}e^{t\mathrm{ad}Y})`$ is well defined.
If $`X`$ and $`Y`$ commute, then we expect to have $`\mathrm{log}\left(e^Xe^Y\right)=\mathrm{log}(e^{X+Y})=X+Y`$. Exercise 3 asks you to verify that the Baker-Campbell-Hausdorff formula indeed gives $`X+Y`$ in that case.
Formula (4.6) is admittedly horrible-looking. However, we are interested not in the details of the formula, but in the fact that it expresses $`\mathrm{log}\left(e^Xe^Y\right)`$ (and hence $`e^Xe^Y`$) in terms of the Lie-algebraic quantities $`\mathrm{ad}X`$ and $`\mathrm{ad}Y`$.
The goal of the Baker-Campbell-Hausdorff theorem is to compute $`\mathrm{log}\left(e^Xe^Y\right)`$. You may well ask, “Why don’t we simply expand both exponentials and the logarithm in power series and multiply everything out?” Well, you can do this, and if you do it for the first several terms you will get the same answer as B-C-H. However, there is a serious problem with this approach, namely: How do you know that the terms in such an expansion are expressible in terms of commutators? Consider for example the quadratic term. It is clear that this will be a linear combination of $`X^2`$, $`Y^2`$, $`XY`$, and $`YX`$. But to be expressible in terms of commutators it must actually be a constant times $`\left(XYYX\right)`$. Of course, for the quadratic term you can just multiply it out and see, and indeed you get $`\frac{1}{2}\left(XYYX\right)=\frac{1}{2}[X,Y]`$. But it is far from clear how to prove that a similar result occurs for all the higher terms. See Exercise 4.
###### Proof.
We begin by proving that the corollary follows from the integral form of the Baker-Campbell-Hausdorff formula. The proof is conceptually similar to the reasoning in Equation (4.3). Note that if $`X`$ and $`Y`$ lie in some Lie algebra $`𝔤`$ then $`\mathrm{ad}X`$ and $`\mathrm{ad}Y`$ will preserve $`𝔤`$, and so also will $`g(e^{\mathrm{ad}X}e^{t\mathrm{ad}Y})(Y)`$. Thus whenever formula (4.6) holds, $`\mathrm{log}\left(e^Xe^Y\right)`$ will lie in $`𝔤`$. It remains only to verify (4.7). The idea is that if $`\stackrel{~}{\varphi }`$ is Lie algebra homomorphism, then it will take a big horrible looking expression involving ‘ad’ and $`X`$ and $`Y`$, and turn it into the same expression with $`X`$ and $`Y`$ replaced by $`\stackrel{~}{\varphi }\left(X\right)`$ and $`\stackrel{~}{\varphi }\left(Y\right)`$.
More precisely, since $`\stackrel{~}{\varphi }`$ is a Lie algebra homomorphism,
$$\stackrel{~}{\varphi }[Y,X]=[\stackrel{~}{\varphi }(Y),\stackrel{~}{\varphi }(X)]$$
or
$$\stackrel{~}{\varphi }\left(\mathrm{ad}Y\left(X\right)\right)=\mathrm{ad}\stackrel{~}{\varphi }\left(Y\right)\left(\stackrel{~}{\varphi }\left(X\right)\right)\text{.}$$
More generally,
$$\stackrel{~}{\varphi }\left(\left(\mathrm{ad}Y\right)^n\left(X\right)\right)=\left(\mathrm{ad}\stackrel{~}{\varphi }\left(Y\right)\right)^n\left(\stackrel{~}{\varphi }\left(X\right)\right)\text{.}$$
This being the case,
$`\stackrel{~}{\varphi }\left(e^{\mathrm{ad}Y}\left(X\right)\right)`$ $`={\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{t^m}{m!}}\stackrel{~}{\varphi }\left(\left(\mathrm{ad}Y\right)^n\left(X\right)\right)`$
$`={\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{t^m}{m!}}\left(\mathrm{ad}\stackrel{~}{\varphi }\left(Y\right)\right)^n\left(\stackrel{~}{\varphi }\left(X\right)\right)`$
$`=e^{t\mathrm{ad}\stackrel{~}{\varphi }(Y)}\left(\stackrel{~}{\varphi }(X)\right)\text{.}`$
Similarly,
$$\stackrel{~}{\varphi }\left(e^{\mathrm{ad}X}e^{t\mathrm{ad}Y}(X)\right)=e^{\mathrm{ad}\stackrel{~}{\varphi }(X)}e^{t\mathrm{ad}\stackrel{~}{\varphi }(Y)}\left(\stackrel{~}{\varphi }(X)\right)\text{.}$$
Assume now that $`X`$ and $`Y`$ are small enough that B-C-H applies to $`X`$ and $`Y`$, and to $`\stackrel{~}{\varphi }(X)`$ and $`\stackrel{~}{\varphi }(Y)`$. Then, using the linearity of the integral and reasoning similar to the above, we have:
$`\stackrel{~}{\varphi }\left(\mathrm{log}\left(e^Xe^Y\right)\right)`$ $`=\stackrel{~}{\varphi }(X)+{\displaystyle _0^1}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}a_m\stackrel{~}{\varphi }\left[\left(e^{\mathrm{ad}X}e^{t\mathrm{ad}Y}I\right)^n(X)\right]dt`$
$`=\stackrel{~}{\varphi }(X)+{\displaystyle _0^1}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}a_m\left(e^{\mathrm{ad}\stackrel{~}{\varphi }(X)}e^{t\mathrm{ad}\stackrel{~}{\varphi }(Y)}I\right)^n(\stackrel{~}{\varphi }(X))dt`$
$`=\mathrm{log}\left(e^{\stackrel{~}{\varphi }(X)}e^{\stackrel{~}{\varphi }(Y)}\right)\text{.}`$
This is what we wanted to show. ∎
Before coming to the proof Baker-Campbell-Hausdorff formula itself, we will obtain a result concerning derivatives of the exponential mapping. This result is valuable in its own right, and will play a central role in our proof of the Baker-Campbell-Hausdorff formula.
Observe that if $`X`$ and $`Y`$ commute, then
$$e^{X+tY}=e^Xe^{tY}$$
and so
$$\frac{d}{dt}|_{t=0}e^{X+tY}=e^X\frac{d}{dt}|_{t=0}e^{tY}=e^XY\text{.}$$
In general, $`X`$ and $`Y`$ do not commute, and
$$\frac{d}{dt}|_{t=0}e^{X+tY}e^XY\text{.}$$
This, as it turns out, is an important point. In particular, note that in the language of multivariate calculus
(4.8)
$$\frac{d}{dt}|_{t=0}e^{X+tY}=\text{ }\{\begin{array}{c}\text{directional derivative of }\text{exp}\text{ at }X\text{,}\hfill \\ \text{in the direction of }Y\hfill \end{array}\text{.}$$
Thus computing the left side of (4.8) is the same as computing all of the directional derivatives of the (matrix-valued) function exp. We expect the directional derivative to be a linear function of $`Y`$, for each fixed $`X`$.
Now, the function
$$\frac{1e^z}{z}=\frac{1(1z+\frac{z^2}{2!}\mathrm{})}{z}$$
is an entire analytic function of $`z`$, even at $`z=0`$, and is given by the power series
$$\frac{1e^z}{z}=\underset{n=1}{\overset{\mathrm{}}{}}(1)^{n1}\frac{z^{n1}}{n!}=1\frac{z}{2!}+\frac{z^2}{3!}\mathrm{}\text{.}$$
This series (which has infinite radius of convergence), make sense when $`z`$ is replaced by a linear operator $`A`$ on some finite-dimensional vector space.
###### Theorem 4.5 (Derivative of Exponential).
Let $`X`$ and $`Y`$ be $`n\times n`$ complex matrices. Then
$`{\displaystyle \frac{d}{dt}}|_{t=0}e^{X+tY}`$ $`=e^X\left\{{\displaystyle \frac{Ie^{\mathrm{ad}X}}{\mathrm{ad}X}}(Y)\right\}`$
(4.9) $`=e^X\left\{Y{\displaystyle \frac{[X,Y]}{2!}}+{\displaystyle \frac{[X,[X,Y]]}{3!}}\mathrm{}\right\}\text{.}`$
More generally, if $`X\left(t\right)`$ is a smooth matrix-valued function, then
(4.10)
$$\frac{d}{dt}|_{t=0}e^{X(t)}=e^{X(0)}\left\{\frac{Ie^{\mathrm{ad}X(0)}}{\mathrm{ad}X(0)}\left(\frac{dX}{dt}|_{t=0}\right)\right\}\text{.}$$
Note that the directional derivative in (4.9) is indeed linear in $`Y`$ for each fixed $`X`$. Note also that (4.9) is just a special case of (4.10), by taking $`X(t)=X+tY`$, and evaluating at $`t=0`$.
Furthermore, observe that if $`X`$ and $`Y`$ commute, then only the first term in the series (4.9) survives. In that case, we obtain $`\frac{d}{dt}|_{t=0}e^{X+tY}=e^XY`$ as expected.
###### Proof.
It is possible to prove this Theorem by expanding everything in a power series and differentiating term-by-term; we will not take that approach. We will prove only form (4.9) of the derivative formula, but the form (4.10) follows by the chain rule.
Let us use the Lie product formula, and let us assume for the moment that it is legal to interchange limit and derivative. (We will consider this issue at the end.) Then we have
$$e^X\frac{d}{dt}|_{t=0}e^{X+tY}=e^X\underset{n\mathrm{}}{lim}\frac{d}{dt}|_{t=0}\left(e^{X/n}e^{tY/n}\right)^n\text{.}$$
We now apply the product rule (generalized to $`n`$ factors) to obtain
$`e^X{\displaystyle \frac{d}{dt}}|_{t=0}e^{X+tY}`$ $`=e^X\underset{n\mathrm{}}{lim}{\displaystyle \underset{k=0}{\overset{n1}{}}}\left[\left(e^{X/n}e^{tY/n}\right)^{nk1}\left(e^{X/n}e^{tY/n}Y/n\right)\left(e^{X/n}e^{tY/n}\right)^k\right]_{t=0}`$
$`=e^X\underset{n\mathrm{}}{lim}{\displaystyle \underset{k=0}{\overset{n1}{}}}\left(e^{X/n}\right)^{nk1}\left(e^{X/n}Y/n\right)\left(e^{X/n}\right)^k`$
$`=\underset{n\mathrm{}}{lim}{\displaystyle \frac{1}{n}}{\displaystyle \underset{k=0}{\overset{n1}{}}}\left(e^{X/n}\right)^kY\left(e^{X/n}\right)^k\text{.}`$
But
$`\left(e^{X/n}\right)^kY\left(e^{X/n}\right)^k`$ $`=\left[\mathrm{Ad}\left(e^{X/n}\right)\right]^k\left(Y\right)`$
$`=\left(e^{\mathrm{ad}X/n}\right)^k(Y)`$
(where we have used the relationship between Ad and ad). So we have
(4.11)
$$e^X\frac{d}{dt}|_{t=0}e^{X+tY}=\underset{n\mathrm{}}{lim}\frac{1}{n}\underset{k=0}{\overset{n1}{}}\left(e^{\mathrm{ad}X/n}\right)^k(Y)\text{.}$$
Observe now that $`_{k=0}^{n1}\left(e^{\mathrm{ad}X/n}\right)^k`$ is a geometric series. Let us now reason for a moment at the purely formal level. Using the usual formula for geometric series, we get
$`e^X{\displaystyle \frac{d}{dt}}|_{t=0}e^{X+tY}=\underset{n\mathrm{}}{lim}{\displaystyle \frac{1}{n}}{\displaystyle \frac{I\left(e^{\mathrm{ad}X/n}\right)^n}{Ie^{\mathrm{ad}X/n}}}(Y)`$
$`=\underset{n\mathrm{}}{lim}{\displaystyle \frac{Ie^{\mathrm{ad}X}}{n\left[I\left(I\frac{\mathrm{ad}X}{n}+\frac{(\mathrm{ad}X)^2}{n^22!}\mathrm{}\right)\right]}}(Y)`$
$`=\underset{n\mathrm{}}{lim}{\displaystyle \frac{Ie^{\mathrm{ad}X}}{\mathrm{ad}X\frac{(\mathrm{ad}X)^2}{n2!}+\mathrm{}}}(Y)`$
$`={\displaystyle \frac{Ie^{\mathrm{ad}X}}{\mathrm{ad}X}}(Y)\text{.}`$
This is what we wanted to show!
Does this argument make sense at any rigorous level? In fact it does. As usual, let us consider first the diagonalizable case. That is, assume that $`\mathrm{ad}X`$ is diagonalizable as an operator on $`\mathrm{𝗀𝗅}(n;)`$, and assume that $`Y`$ is an eigenvector for $`\mathrm{ad}X`$. This means that $`\mathrm{ad}X(Y)=[X,Y]=\lambda Y`$, for some $`\lambda `$. Now, there are two cases, $`\lambda =0`$ and $`\lambda 0`$. The $`\lambda =0`$ case corresponds to the case in which $`X`$ and $`Y`$ commute, and we have already observed that the Theorem holds trivially in that case.
The interesting case, then, is the case $`\lambda 0`$. Note that $`\left(\mathrm{ad}X\right)^n(Y)=\lambda ^nY`$, and so
$$\left(e^{\mathrm{ad}X/n}\right)^k(Y)=\left(e^{\lambda /n}\right)^k(Y)\text{.}$$
Thus the geometric series in (4.11) becomes an ordinary complex-valued series, with ratio $`e^{\lambda /n}`$. Since $`\lambda 0`$, this ratio will be different from one for all sufficiently large $`n`$. Thus we get
$$e^X\frac{d}{dt}|_{t=0}e^{X+tY}=\left(\underset{n\mathrm{}}{lim}\frac{1}{n}\frac{I\left(e^{\lambda /n}\right)^n}{Ie^{\lambda /n}}\right)Y\text{.}$$
There is now no trouble in taking the limit as we did formally above to get
$`e^X{\displaystyle \frac{d}{dt}}|_{t=0}e^{X+tY}={\displaystyle \frac{1e^\lambda }{\lambda }}Y`$
$`={\displaystyle \frac{Ie^{\mathrm{ad}X}}{\mathrm{ad}X}}(Y)\text{.}`$
We see then that the Theorem holds in the case that $`\mathrm{ad}X`$ is diagonalizable and $`Y`$ is an eigenvector of $`\mathrm{ad}X`$. If $`\mathrm{ad}X`$ is diagonalizable but $`Y`$ is not an eigenvector, then $`Y`$ is a linear combination of eigenvectors and applying the above computation to each of those eigenvectors gives the desired result.
We need, then, to consider the case where $`\mathrm{ad}X`$ is not diagonalizable. But (Exercise 20), if $`X`$ is a diagonalizable matrix, then $`\mathrm{ad}X`$ will be diagonalizable as an operator on $`\mathrm{𝗀𝗅}(n;)`$. Since, as we have already observed, every matrix is the limit of diagonalizable matrices, we are essentially done. For it is easy to see by differentiating the power series term-by-term that $`e^X\frac{d}{dt}|_{t=0}e^{X+tY}`$ exists and varies continuously with $`X`$. Thus once we have the Theorem for all diagonalizable $`X`$ we have it for all $`X`$ by passing to the limit.
The only unresolved issue, then, is the interchange of limit and derivative which we performed at the very beginning of the argument. I do not want to spell this out in detail, but let us see what would be involved in justifying this. A standard theorem in elementary analysis says that if $`f_n(t)f(t)`$ pointwise, and in addition $`df_n/dt`$ converges uniformly to some function $`g(t)`$, then $`f(t)`$ is differentiable and $`df/dt=g(t)`$. (E.g., Theorem 7.17 in W. Rudin’s Principles of Mathematical Analysis.) The key requirement is that the derivatives converge uniformly. Uniform convergence of the $`f_n`$’s themselves is definitely not sufficient.
In our case, $`f_n(t)=e^X\left(e^{X/n}e^{tY/n}\right)^n`$. The Lie product formula says that this converges pointwise to $`e^Xe^{X+tY}`$. We need, then, to show that
$$\frac{d}{dt}e^X\left(e^{X/n}e^{tY/n}\right)^n$$
converges uniformly to some $`g(t)`$, say on the interval $`1t1`$. This computation is similar to what we did above, with relatively minor modifications to account for the fact that we do not take $`t=0`$ and to make sure the convergence is uniform. This part of the proof is left as an exercise to the reader. ∎
#### 4.2.1. Proof of the Baker-Campbell-Hausdorff Formula
We now turn to the proof of the Baker-Campbell-Hausdorff formula itself. Our argument follows Miller, Sec. 5.1, with minor differences of convention. (Warning: Miller’s “Ad” is what we call “ad.”) Define
$$Z(t)=\mathrm{log}\left(e^Xe^{tY}\right)$$
If $`X`$ and $`Y`$ are sufficiently small, then $`Z\left(t\right)`$ is defined for $`0t1`$. It is left as an exercise to verify that $`Z(t)`$ is smooth. Our goal is to compute $`Z(1).`$
By definition
$$e^{Z(t)}=e^Xe^{tY}$$
so that
$$e^{Z(t)}\frac{d}{dt}e^{Z(t)}=\left(e^Xe^{tY}\right)^1e^Xe^{tY}Y=Y\text{.}$$
On the other hand, by Theorem 4.5,
$$e^{Z(t)}\frac{d}{dt}e^{Z(t)}=\left\{\frac{Ie^{\mathrm{ad}Z(t)}}{\mathrm{ad}Z(t)}\right\}\left(\frac{dZ}{dt}\right)\text{.}$$
Hence
$$\left\{\frac{Ie^{\mathrm{ad}Z(t)}}{\mathrm{ad}Z(t)}\right\}\left(\frac{dZ}{dt}\right)=Y\text{.}$$
If $`X`$ and $`Y`$ are small enough, then $`Z(t)`$ will also be small, so that $`\left(Ie^{\mathrm{ad}Z(t)}\right)/\mathrm{ad}Z(t)`$ will be close to the identity and thus invertible. So
(4.12)
$$\frac{dZ}{dt}=\left\{\frac{Ie^{\mathrm{ad}Z(t)}}{\mathrm{ad}Z(t)}\right\}^1(Y)\text{.}$$
Recall that $`e^{Z(t)}=e^Xe^{tY}`$. Applying the homomorphism ‘Ad’ gives
$$\mathrm{Ad}\left(e^{Z(t)}\right)=\mathrm{Ad}\left(e^X\right)\mathrm{Ad}\left(e^{tY}\right)\text{.}$$
By the relationship (3.17) between ‘Ad’ and ‘ad,’ this becomes
$$e^{\mathrm{ad}Z(t)}=e^{\mathrm{ad}X}e^{t\mathrm{ad}Y}$$
or
$$\mathrm{ad}Z(t)=\mathrm{log}\left(e^{\mathrm{ad}X}e^{t\mathrm{ad}Y}\right)\text{.}$$
Plugging this into (4.12) gives
(4.13)
$$\frac{dZ}{dt}=\left\{\frac{I\left(e^{\mathrm{ad}X}e^{t\mathrm{ad}Y}\right)^1}{\mathrm{log}\left(e^{\mathrm{ad}X}e^{t\mathrm{ad}Y}\right)}\right\}^1(Y)\text{.}$$
But now observe that
$$g(z)=\left\{\frac{1z^1}{\mathrm{log}z}\right\}^1$$
so, formally, (4.13) is the same as
(4.14)
$$\frac{dZ}{dt}=g\left(e^{\mathrm{ad}X}e^{t\mathrm{ad}Y}\right)(Y)\text{.}$$
Reasoning as in the proof of Theorem 4.5 shows easily that this formal argument is actually correct.
Now we are essentially done, for if we note that $`Z(0)=X`$ and integrate (4.14), we get
$$Z(1)=X+_0^1g(e^{\mathrm{ad}X}e^{t\mathrm{ad}Y})(Y)𝑑t$$
which is the Baker-Campbell-Hausdorff formula.
### 4.3. The Series Form of the Baker-Campbell-Hausdorff Formula
Let us see how to get the first few terms of the series form of B-C-H from the integral form. Recall the function
$`g\left(z\right)`$ $`={\displaystyle \frac{z\mathrm{log}z}{z1}}`$
$`={\displaystyle \frac{\left[1+\left(z1\right)\right]\left[\left(z1\right)\frac{\left(z1\right)^2}{2}+\frac{\left(z1\right)^3}{3}\mathrm{}\right]}{\left(z1\right)}}`$
$`=\left[1+\left(z1\right)\right]\left[1{\displaystyle \frac{z1}{2}}+{\displaystyle \frac{\left(z1\right)^2}{3}}\right]\text{.}`$
Multiplying this out and combining terms gives
$$g\left(z\right)=1+\frac{1}{2}\left(z1\right)\frac{1}{6}\left(z1\right)^2+\mathrm{}\text{.}$$
The closed-form expression for $`g`$ is
$$g\left(z\right)=1+\underset{n=1}{\overset{\mathrm{}}{}}\frac{\left(1\right)^{n+1}}{n\left(n+1\right)}\left(z1\right)^n\text{.}$$
Meanwhile
$`e^{\mathrm{ad}X}e^{t\mathrm{ad}Y}I`$ $`=\left(I+\mathrm{ad}X+{\displaystyle \frac{\left(\mathrm{ad}X\right)^2}{2}}+\mathrm{}\right)\left(I+t\mathrm{ad}Y+{\displaystyle \frac{t^2\left(\mathrm{ad}Y\right)^2}{2}}+\mathrm{}\right)I`$
$`=\mathrm{ad}X+t\mathrm{ad}Y+t\mathrm{ad}X\mathrm{ad}Y+{\displaystyle \frac{\left(\mathrm{ad}X\right)^2}{2}}+{\displaystyle \frac{t^2\left(\mathrm{ad}Y\right)^2}{2}}+\mathrm{}\text{.}`$
The crucial observation here is that $`e^{\mathrm{ad}X}e^{t\mathrm{ad}Y}I`$ has no zero-order term, just first-order and higher in $`\mathrm{ad}X/\mathrm{ad}Y`$. Thus $`\left(e^{\mathrm{ad}X}e^{t\mathrm{ad}Y}I\right)^n`$ will contribute only terms of degree $`n`$ or higher in $`\mathrm{ad}X/\mathrm{ad}Y`$.
We have, then, up to degree two in $`\mathrm{ad}X/\mathrm{ad}Y`$
$`g\left(e^{\mathrm{ad}X}e^{t\mathrm{ad}Y}\right)`$ $`=I+{\displaystyle \frac{1}{2}}\left[\mathrm{ad}X+t\mathrm{ad}Y+t\mathrm{ad}X\mathrm{ad}Y+{\displaystyle \frac{\left(\mathrm{ad}X\right)^2}{2}}+{\displaystyle \frac{t^2\left(\mathrm{ad}Y\right)^2}{2}}+\mathrm{}\right]`$
$`{\displaystyle \frac{1}{6}}\left[\mathrm{ad}X+t\mathrm{ad}Y+\mathrm{}\right]^2`$
$`=I+{\displaystyle \frac{1}{2}}\mathrm{ad}X+{\displaystyle \frac{t}{2}}\mathrm{ad}Y+{\displaystyle \frac{t}{2}}\mathrm{ad}X\mathrm{ad}Y+{\displaystyle \frac{\left(\mathrm{ad}X\right)^2}{4}}+{\displaystyle \frac{t^2\left(\mathrm{ad}Y\right)^2}{4}}`$
$`{\displaystyle \frac{1}{6}}\left[\left(\mathrm{ad}X\right)^2+t^2\left(\mathrm{ad}Y\right)^2+t\mathrm{ad}X\mathrm{ad}Y+t\mathrm{ad}Y\mathrm{ad}X\right]`$
$`+\text{ higher-order terms.}`$
We now to apply $`g\left(e^{\mathrm{ad}X}e^{t\mathrm{ad}Y}\right)`$ to $`Y`$ and integrate. So (neglecting higher-order terms) by B-C-H, and noting that any term with $`adY`$ acting first is zero:
$`\mathrm{log}\left(e^Xe^Y\right)`$
$`=X+{\displaystyle _0^1}\left[Y+{\displaystyle \frac{1}{2}}[X,Y]+{\displaystyle \frac{1}{4}}[X,[X,Y]]{\displaystyle \frac{1}{6}}[X,[X,Y]]{\displaystyle \frac{t}{6}}[Y,[X,Y]]\right]𝑑t`$
$`=X+Y+{\displaystyle \frac{1}{2}}[X,Y]+\left({\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{6}}\right)[X,[X,Y]]{\displaystyle \frac{1}{6}}{\displaystyle _0^1}t𝑑t[Y,[X,Y]]\text{.}`$
Thus if we do the algebra we end up with
$`\mathrm{log}\left(e^Xe^Y\right)`$ $`=X+Y+{\displaystyle \frac{1}{2}}[X,Y]+{\displaystyle \frac{1}{12}}[X,[X,Y]]{\displaystyle \frac{1}{12}}[Y,[X,Y]]`$
$`+\text{ higher order terms.}`$
This is the expression in (4.2).
### 4.4. Subgroups and Subalgebras
Suppose that $`G`$ is a matrix Lie group, $`H`$ another matrix Lie group, and suppose that $`HG`$. Then certainly the Lie algebra $`𝔥`$ of $`H`$ will be a subalgebra of the Lie algebra $`𝔤`$ of $`G`$. Does this go the other way around? That is given a Lie group $`G`$ with Lie algebra $`𝔤`$, and a subalgebra $`𝔥`$ of $`𝔤`$, is there a matrix Lie group $`H`$ whose Lie algebra is $`𝔥`$?
In the case of the Heisenberg group, the answer is yes. This is easily seen using the fact that the exponential mapping is one-to-one and onto, together with the special form of the Baker-Campbell-Hausdorff formula. (See Exercise 6.)
Unfortunately, the answer in general is no. For example, let $`G=\mathrm{𝖦𝖫}(2;)`$ and let
$$𝔥=\left\{\left(\begin{array}{cc}it& 0\\ 0& ita\end{array}\right)\right|t\}\text{,}$$
where $`a`$ is irrational. If there is going to be a matrix Lie group $`H`$ with Lie algebra $`𝔥`$, then $`H`$ would contain the set
$$H_0=\left\{\left(\begin{array}{cc}e^{it}& 0\\ 0& e^{ita}\end{array}\right)\right|t\}\text{.}$$
To be a matrix Lie group, $`H`$ would have to be closed in $`\mathrm{𝖦𝖫}(2;)`$, and so it would contain the closure of $`H_0`$, which (see ) is the set
$$H_1=\left\{\left(\begin{array}{cc}e^{it}& 0\\ 0& e^{is}\end{array}\right)\right|s,t\}\text{.}$$
But then the Lie algebra of $`H`$ would have to contain the Lie algebra of $`H_1`$, which is two-dimensional!
Fortunately, all is not lost. We can still get a subgroup $`H`$ for each subalgebra $`𝔥`$, if we weaken the condition that $`H`$ be a matrix Lie group. In the above example, the subgroup we want is $`H_0`$, despite the fact that $`H_0`$ is not a matrix Lie group.
###### Definition 4.6.
If $`H`$ is any subgroup of $`\mathrm{𝖦𝖫}(n;)`$, define the Lie algebra $`𝔥`$ of $`H`$ to be the set of all matrices $`X`$ such that
$$e^{tX}H$$
for all real $`t.`$
###### Definition 4.7.
If $`G`$ is a matrix Lie group with Lie algebra $`𝔤`$, then $`H`$ is a connected Lie subgroup of $`G`$ if
i) $`H`$ is a subgroup of $`G`$
ii) $`H`$ is connected
iii) the Lie algebra $`𝔥`$ of $`H`$ is a subspace of $`𝔤`$
iv) Every element of $`H`$ can be written in the form $`e^{X_1}e^{X_2}\mathrm{}e^{X_n}`$, with $`X_1,\mathrm{},X_n𝔥`$.
###### Theorem 4.8.
If $`G`$ is a matrix Lie group with Lie algebra $`𝔤`$, and $`H`$ is a connected Lie subgroup of $`G`$, then the Lie algebra $`𝔥`$ of $`H`$ is a subalgebra of $`𝔤`$.
###### Proof.
Since by definition $`𝔥`$ is a subspace of $`𝔤`$, it remains only to show that $`𝔥`$ is closed under brackets. So assume $`X,Y𝔥`$. Then $`e^{tX}`$ and $`e^{sY}`$ are in $`H`$, and so (since $`H`$ is a subgroup) is the element
$$e^{tX}e^{sY}e^{tX}=\mathrm{exp}\left[s\left(e^{tX}Ye^{tX}\right)\right]\text{.}$$
This shows that $`e^{tX}Ye^{tX}`$ is in $`𝔥`$ for all $`t`$. But $`𝔥`$ is a subspace of $`𝔤`$, which is necessarily a closed subset of $`𝔤`$. Thus
$$[X,Y]=\frac{d}{dt}|_{t=0}e^{tX}Ye^{tX}=\underset{h0}{lim}\frac{\left(e^{hX}Ye^{hX}Y\right)}{h}$$
is in $`𝔥`$. (This argument is precisely the one we used to show that the Lie algebra of a matrix Lie group is a closed under brackets, once we had established that it is a subspace.) ∎
We are now ready to state the main theorem of this section, which is our second major application of the Baker-Campbell-Hausdorff formula.
###### Theorem 4.9.
Let $`G`$ be a matrix Lie group with Lie algebra $`𝔤`$. Let $`𝔥`$ be a Lie subalgebra of $`𝔤`$. Then there exists a unique connected Lie subgroup $`H`$ of $`G`$ such that the Lie algebra of $`H`$ is $`𝔥`$.
Given a matrix Lie group $`G`$ and a subalgebra $`𝔥`$ of $`𝔤`$, the associated connected Lie subgroup $`H`$ might be a matrix Lie group. This will happen precisely if $`H`$ is a closed subset of $`G`$. There are various conditions under which you can prove that $`H`$ is closed. For example, if $`G=\mathrm{𝖦𝖫}(n;)`$, and $`𝔥`$ is semisimple, then $`H`$ is automatically closed, and hence a matrix Lie group. (See Helgason, Chapter II, Exercises and Further Results, D.)
If only the Baker-Campbell-Hausdorff formula worked globally instead of only locally the proof of this theorem would be easy. If the B-C-H formula converged for all $`X,Y`$ we could just define $`H`$ to be the image of $`𝔥`$ under the exponential mapping. In that case B-C-H would show that this image is a subgroup, since then we would have $`e^{H_1}e^{H_2}=e^Z,`$ with $`Z=H_1+H_2+\frac{1}{2}[H_1,H_2]+\mathrm{}𝔥`$ provided that $`H_1,H_2𝔥.`$ Unfortunately, the B-C-H formula is not convergent in general, and in general the image of $`H`$ under the exponential mapping is not a subgroup.
###### Proof.
Not written at this time. ∎
### 4.5. Exercises
1. The center of a Lie algebra $`𝔤`$ is defined to be the set of all $`X𝔤`$ such that $`[X,Y]=0`$ for all $`Y𝔤`$. Now consider the Heisenberg group
$$H=\left\{\left(\begin{array}{ccc}1& a& b\\ 0& 1& c\\ 0& 0& 1\end{array}\right)\right|a,b,c\}$$
with Lie algebra
$$𝔥=\left\{\left(\begin{array}{ccc}0& \alpha & \beta \\ 0& 0& \gamma \\ 0& 0& 0\end{array}\right)\right|\alpha ,\beta ,\gamma \}\text{.}$$
Determine the center $`Z(𝔥)`$ of $`𝔥`$. For any $`X,Y𝔥`$, show that $`[X,Y]Z(𝔥)`$. This implies, in particular that both $`X`$ and $`Y`$ commute with their commutator $`[X,Y]`$.
Show by direct computation that for any $`X,Y𝔥`$,
(4.15)
$$e^Xe^Y=e^{X+Y+{\scriptscriptstyle \frac{1}{2}}[X,Y]}\text{.}$$
2. Let $`X`$ be a $`n\times n`$ complex matrix. Show that
$$\frac{Ie^X}{X}$$
is invertible if and only if $`X`$ has no eigenvalue of the form $`\lambda =2\pi in`$, with $`n`$ an non-zero integer.
Hint: When is $`\left(1e^z\right)/z`$ equal to zero?
Remark: This exercise, combined with the formula in Theorem 4.5, gives the following result (in the language of differentiable manifolds): The exponential mapping $`\mathrm{exp}:𝔤G`$ is a local diffeomorphism near $`X𝔤`$ if and only $`\mathrm{ad}X`$ has no eigenvalue of the form $`\lambda =2\pi in`$, with $`n`$ a non-zero integer.
3. Verify that the right side of the Baker-Campbell-Hausdorff formula (4.6) reduces to $`X+Y`$ in the case that $`X`$ and $`Y`$ commute.
Hint: Compute first $`e^{\mathrm{ad}X}e^{t\mathrm{ad}Y}(Y)`$ and $`\left(e^{\mathrm{ad}X}e^{t\mathrm{ad}Y}I\right)(Y)`$.
4. Compute $`\mathrm{log}\left(e^Xe^Y\right)`$ through third order in $`X/Y`$ by using the power series for the exponential and the logarithm. Show that you get the same answer as the Baker-Campbell-Hausdorff formula.
5. Using the techniques in Section 4.3, compute the series form of the Baker-Campbell-Hausdorff formula up through fourth-order brackets. (We have already computed up through third-order brackets.)
6. Let $`𝔞`$ be a subalgebra of the Lie algebra of the Heisenberg group. Show that $`\mathrm{exp}\left(𝔞\right)`$ is a connected Lie subgroup of the Heisenberg group. Show that in fact $`\mathrm{exp}\left(𝔞\right)`$ is a matrix Lie group.
7. Show that every connected Lie subgroup of $`\mathrm{𝖲𝖴}\left(2\right)`$ is closed. Show that this is not the case for $`\mathrm{𝖲𝖴}\left(3\right)`$.
## Chapter 5 Basic Representation Theory
### 5.1. Representations
###### Definition 5.1.
Let $`G`$ be a matrix Lie group. Then a finite-dimensional complex representation of $`G`$ is a Lie group homomorphism
$$\mathrm{\Pi }:G\mathrm{𝖦𝖫}(n;)$$
($`n1`$) or more generally a Lie group homomorphism
$$\mathrm{\Pi }:G\mathrm{𝖦𝖫}(V)$$
where $`V`$ is a finite-dimensional complex vector space (with dim$`(V)1`$). A finite-dimensional real representation of $`G`$ is a Lie group homomorphism $`\mathrm{\Pi }`$ of $`G`$ into $`\mathrm{𝖦𝖫}(n;)`$ or into $`\mathrm{𝖦𝖫}(V)`$, where $`V`$ is a finite-dimensional real vector space.
If $`𝔤`$ is a real or complex Lie algebra, then a finite-dimensional complex representation of $`𝔤`$ is a Lie algebra homomorphism $`\pi `$ of $`𝔤`$ into $`\mathrm{𝗀𝗅}(n;)`$ or into gl$`(\mathrm{V})`$, where $`V`$ is a finite-dimensional complex vector space. If $`𝔤`$ is a real Lie algebra, then a finite-dimensional real representation of $`𝔤`$ is a Lie algebra homomorphism $`\pi `$ of $`𝔤`$ into $`\mathrm{𝗀𝗅}(n;)`$ or into $`\mathrm{𝗀𝗅}(V)`$.
If $`\mathrm{\Pi }`$ or $`\pi `$ is a one-to-one homomorphism, then the representation is called faithful.
You should think of a representation as a (linear) action of a group or Lie algebra on a vector space. (Since, say, to every $`gG`$ there is associated an operator $`\mathrm{\Pi }(g)`$, which acts on the vector space $`V`$.) In fact, we will use terminology such as, “Let $`\mathrm{\Pi }`$ be a representation of $`G`$ acting on the space $`V`$.” Even if $`𝔤`$ is a real Lie algebra, we will consider mainly complex representations of $`𝔤`$. After making a few more definitions, we will discuss the question of why one should be interested in studying representations.
###### Definition 5.2.
Let $`\mathrm{\Pi }`$ be a finite-dimensional real or complex representation of a matrix Lie group $`G`$, acting on a space $`V`$. A subspace $`W`$ of $`V`$ is called invariant if $`\mathrm{\Pi }(A)wW`$ for all $`wW`$ and all $`AG`$. An invariant subspace $`W`$ is called non-trivial if $`W\{0\}`$ and $`WV`$. A representation with no non-trivial invariant subspaces is called irreducible.
The terms invariant, non-trivial, and irreducible are defined analogously for representations of Lie algebras.
###### Definition 5.3.
Let $`G`$ be a matrix Lie group, let $`\mathrm{\Pi }`$ be a representation of $`G`$ acting on the space $`V`$, and let $`\mathrm{\Sigma }`$ be a representation of $`G`$ acting on the space $`W`$. A linear map $`\varphi :VW`$ is called a morphism (or intertwining map) of representations if
$$\varphi (\mathrm{\Pi }(A)v)=\mathrm{\Sigma }(A)\varphi (v)$$
for all $`AG`$ and all $`vV`$. The analogous property defines morphisms of representations of a Lie algebra.
If $`\varphi `$ is a morphism of representations, and in addition $`\varphi `$ is invertible, then $`\varphi `$ is said to be an isomorphism of representations. If there exists an isomorphism between $`V`$ and $`W`$, then the representations are said to be isomorphic (or equivalent).
Two isomorphic representations should be regarded as being “the same” representation. A typical problem in representation theory is to determine, up to isomorphism, all the irreducible representations of a particular group or Lie algebra. In Section 5.4 we will determine all the finite-dimensional complex irreducible representations of the Lie algebra $`\mathrm{𝗌𝗎}(2)`$.
###### Proposition 5.4.
Let $`G`$ be a matrix Lie group with Lie algebra $`𝔤`$, and let $`\mathrm{\Pi }`$ be a (finite-dimensional real or complex) representation of $`G`$, acting on the space $`V`$. Then there is a unique representation $`\pi `$ of $`𝔤`$ acting on the same space such that
$$\mathrm{\Pi }(e^X)=e^{\pi (X)}$$
for all $`X𝔤`$. The representation $`\pi `$ can be computed as
$$\pi (X)=\frac{d}{dt}|_{t=0}\mathrm{\Pi }\left(e^{tX}\right)$$
and satisfies
$$\pi \left(AXA^1\right)=\mathrm{\Pi }(A)\pi (X)\mathrm{\Pi }(A)^1$$
for all $`X𝔤`$ and all $`AG`$.
###### Proof.
Theorem 3.18 in Chapter 3 states that for each Lie group homomorphism $`\varphi :GH`$ there is an associated Lie algebra homomorphism $`\stackrel{~}{\varphi }:𝔤𝔥`$. Take $`H=\mathrm{𝖦𝖫}(V)`$ and $`\varphi =\mathrm{\Pi }`$. Since the Lie algebra of $`\mathrm{𝖦𝖫}(V)`$ is $`\mathrm{𝗀𝗅}(V)`$ (since the exponential of any operator is invertible), the associated Lie algebra homomorphism $`\stackrel{~}{\varphi }=\pi `$ maps from $`𝔤`$ to $`\mathrm{𝗀𝗅}(V)`$, and so constitutes a representation of $`𝔤`$.
The properties of $`\pi `$ follow from the properties of $`\stackrel{~}{\varphi }`$ given in Theorem 6. ∎
###### Proposition 5.5.
Let $`𝔤`$ be a real Lie algebra, and $`𝔤_{}`$ its complexification. Then every finite-dimensional complex representation $`\pi `$ of $`𝔤`$ has a unique extension to a (complex-linear) representation of $`𝔤_{}`$, also denoted $`\pi `$. The representation of $`𝔤_{}`$ satisfies
$$\pi (X+iY)=\pi (X)+i\pi (Y)$$
for all $`X𝔤`$.
###### Proof.
This follows from Exercise 14 of Chapter 3. ∎
###### Definition 5.6.
Let $`G`$ be a matrix Lie group, let $``$ be a Hilbert space, and let $`U()`$ denote the group of unitary operators on $``$. Then a homomorphism $`\mathrm{\Pi }:GU()`$ is called a unitary representation of $`G`$ if $`\mathrm{\Pi }`$ satisfies the following continuity condition: If $`A_n,AG`$ and $`A_nA`$, then
$$\mathrm{\Pi }(A_n)v\mathrm{\Pi }(A)v$$
for all $`v`$. A unitary representation with no non-trivial closed invariant subspaces is called irreducible.
This continuity condition is called strong continuity. One could require the even stronger condition that $`\mathrm{\Pi }(A_n)\mathrm{\Pi }(A)0`$, but this turns out to be too stringent a requirement. (That is, most of the interesting representations of $`G`$ will not have this stronger continuity condition.) In practice, any homomorphism of $`G`$ into $`U()`$ you can write down explicitly will be strongly continuous.
One could try to define some analog of unitary representations for Lie algebras, but there are serious technical difficulties associated with getting the “right” definition.
### 5.2. Why Study Representations?
If a representation $`\mathrm{\Pi }`$ is a faithful representation of a matrix Lie group $`G`$, then $`\left\{\mathrm{\Pi }(A)\right|AG\}`$ is a group of matrices which is isomorphic to the original group $`G`$. Thus $`\mathrm{\Pi }`$ allows us to represent $`G`$ as a group of matrices. This is the motivation for the term representation. (Of course, we still call $`\mathrm{\Pi }`$ a representation even if it is not faithful.)
Despite the origin of the term, the point of representation theory is not (at least in this course) to represent a group as a group of matrices. After all, all of our groups are already matrix groups! While it might seem redundant to study representations of a group which is already represented as a group of matrices, this is precisely what we are going to do.
The reason for this is that a representation can be thought of (as we have already noted) as an action of our group on some vector space. Such actions (representations) arise naturally in many branches of both mathematics and physics, and it is important to understand them.
A typical example would be a differential equation in three-dimensional space which has rotational symmetry. If the equation has rotational symmetry, then the space of solutions will be invariant under rotations. Thus the space of solutions will constitute a representation of the rotation group $`\mathrm{𝖲𝖮}(3)`$. If you know what all of the representations of $`\mathrm{𝖲𝖮}(3)`$ are, this can help immensely in narrowing down what the space of solutions can be. (As we will see, $`\mathrm{𝖲𝖮}(3)`$ has lots of other representations besides the obvious one in which $`\mathrm{𝖲𝖮}(3)`$ acts on $`^3`$.)
In fact, one of the chief applications of representation theory is to exploit symmetry. If a system has symmetry, then the set of symmetries will form a group, and understanding the representations of the symmetry group allows you to use that symmetry to simplify the problem.
In addition, studying the representations of a group $`G`$ (or of a Lie algebra $`𝔤`$) can give information about the group (or Lie algebra) itself. For example, if $`G`$ is a finite group, then associated to $`G`$ is something called the group algebra. The structure of this group algebra can be described very nicely in terms of the irreducible representations of $`G`$.
In this course, we will be interested primarily in computing the finite-dimensional irreducible complex representations of matrix Lie groups. As we shall see, this problem can be reduced almost completely to the problem of computing the finite-dimensional irreducible complex representations of the associated Lie algebra. In this chapter, we will discuss the theory at an elementary level, and will consider in detail the example of $`\mathrm{𝖲𝖮}(3)`$ and $`\mathrm{𝖲𝖴}(2)`$. In Chapter 6, we will study the representations of $`\mathrm{𝖲𝖴}(3)`$, which is substantially more involved than that of $`\mathrm{𝖲𝖴}(2)`$, and give an overview of the representation theory of a very important class of Lie groups, namely, the semisimple ones.
### 5.3. Examples of Representations
#### 5.3.1. The Standard Representation
A matrix Lie group $`G`$ is by definition a subset of some $`\mathrm{𝖦𝖫}(n;)`$ or $`\mathrm{𝖦𝖫}(n;)`$. The inclusion map of $`G`$ into $`\mathrm{𝖦𝖫}(n)`$ (i.e., $`\mathrm{\Pi }(A)=A`$) is a representation of $`G`$, called the standard representation of $`G`$. Thus for example the standard representation of $`\mathrm{𝖲𝖮}(3)`$ is the one in which $`\mathrm{𝖲𝖮}(3)`$ acts in the usual way on $`^3`$. If $`G`$ is a subgroup of $`\mathrm{𝖦𝖫}(n;)`$ or $`\mathrm{𝖦𝖫}(n;)`$, then its Lie algebra $`𝔤`$ will be a subalgebra of $`\mathrm{𝗀𝗅}(n;)`$ or $`\mathrm{𝗀𝗅}(n;)`$. The inclusion of $`𝔤`$ into $`\mathrm{𝗀𝗅}(n;)`$ or $`\mathrm{𝗀𝗅}(n;)`$ is a representation of $`𝔤`$, called the standard representation.
#### 5.3.2. The Trivial Representation
Consider the one-dimensional complex vector space $``$. Given any matrix Lie group $`G`$, we can define the trivial representation of $`G`$, $`\mathrm{\Pi }:G\mathrm{𝖦𝖫}(1;)`$, by the formula
$$\mathrm{\Pi }(A)=I$$
for all $`AG`$. Of course, this is an irreducible representation, since $``$ has no non-trivial subspaces, let alone non-trivial invariant subspaces. If $`𝔤`$ is a Lie algebra, we can also define the trivial representation of $`𝔤`$, $`\pi :𝔤\mathrm{𝗀𝗅}(1;)`$, by
$$\pi (X)=0$$
for all $`X𝔤`$. This is an irreducible representation.
#### 5.3.3. The Adjoint Representation
Let $`G`$ be a matrix Lie group with Lie algebra $`𝔤`$. We have already defined the adjoint mapping
$$\mathrm{Ad}:G\mathrm{𝖦𝖫}(𝔤)$$
by the formula
$$\mathrm{𝖠𝖽}A(X)=AXA^1\text{.}$$
Recall that Ad is a Lie group homomorphism. Since Ad is a Lie group homomorphism into a group of invertible operators, we see that in fact Ad is a representation of $`G`$, acting on the space $`𝔤`$. Thus we can now give Ad its proper name, the adjoint representation of $`G`$. The adjoint representation is a real representation of $`G`$.
Similarly, if $`𝔤`$ is a Lie algebra, we have
$$\mathrm{ad}:𝔤\mathrm{𝗀𝗅}(𝔤)$$
defined by the formula
$$\mathrm{𝖺𝖽}X(Y)=[X,Y]\text{.}$$
We know that ad is a Lie algebra homomorphism (Chapter 3, Proposition 3.33), and is therefore a representation of $`𝔤`$, called the adjoint representation. In the case that $`𝔤`$ is the Lie algebra of some matrix Lie group $`G`$, we have already established (Chapter 3, Proposition 3.21 and Exercise 13) that Ad and ad are related as in Proposition 5.4.
Note that in the case of $`\mathrm{𝖲𝖮}(3)`$ the standard representation and the adjoint representation are both three dimensional real representations. In fact these two representations are equivalent (Exercise 4).
#### 5.3.4. Some Representations of $`\mathrm{𝖲𝖴}(2)`$
Consider the space $`V_m`$ of homogeneous polynomials in two complex variables with total degree $`m`$ ($`m0`$). That is, $`V_m`$ is the space of functions of the form
(5.1)
$$f(z_1,z_2)=a_0z_1^m+a_1z_1^{m1}z_2+a_2z_1^{m2}z_2^2\mathrm{}+a_mz_2^m$$
with $`z_1,z_2`$ and the $`a_i`$’s arbitrary complex constants. The space $`V_m`$ is an $`(m+1)`$-dimensional complex vector space.
Now by definition an element $`U`$ of $`\mathrm{𝖲𝖴}(2)`$ is a linear transformation of $`^2`$. Let $`z`$ denote the pair $`z=(z_1,z_2)`$ in $`^2`$. Then we may define a linear transformation $`\mathrm{\Pi }_m(U)`$ on the space $`V_m`$ by the formula
(5.2)
$$\left[\mathrm{\Pi }_m(U)f\right](z)=f(U^1z)\text{.}$$
Explicitly, if $`f`$ is as in (5.1), then
$$\left[\mathrm{\Pi }_m(U)f\right](z_1,z_2)=\underset{k=0}{\overset{m}{}}a_k\left(U_{11}^1z_1+U_{12}^1z_2\right)^{mk}\left(U_{21}^1z_1+U_{22}^1z_2\right)^k\text{.}$$
By expanding out the right side of this formula we see that $`\mathrm{\Pi }_m(U)f`$ is again a homogeneous polynomial of degree $`m`$. Thus $`\mathrm{\Pi }_m(U)`$ actually maps $`V_m`$ into $`V_m`$.
Now, compute
$`\mathrm{\Pi }_m\left(U_1\right)\left[\mathrm{\Pi }_m\left(U_2\right)f\right](z)`$ $`=\left[\mathrm{\Pi }_m\left(U_2\right)f\right](U_1^1z)=f\left(U_2^1U_1^1z\right)`$
$`=\mathrm{\Pi }_m\left(U_1U_2\right)f(z)\text{.}`$
Thus $`\mathrm{\Pi }_m`$ is a (finite-dimensional complex) representation of $`\mathrm{𝖲𝖴}(2)`$. (It is very easy to do the above computation incorrectly.) The inverse in definition (5.2) is necessary in order to make $`\mathrm{\Pi }_m`$ a representation. It turns out that each of the representations $`\mathrm{\Pi }_m`$ of $`\mathrm{𝖲𝖴}(2)`$ is irreducible, and that every finite-dimensional irreducible representation of $`\mathrm{𝖲𝖴}(2)`$ is equivalent to one (and only one) of the $`\mathrm{\Pi }_m`$’s. (Of course, no two of the $`\mathrm{\Pi }_m`$’s are equivalent, since they don’t even have the same dimension.)
Let us now compute the corresponding Lie algebra representation $`\pi _m`$. According to Proposition 5.4, $`\pi _m`$ can be computed as
$$\pi _m(X)=\frac{d}{dt}|_{t=0}\mathrm{\Pi }_m\left(e^{tX}\right)\text{.}$$
So
$$\left(\pi _m(X)f\right)(z)=\frac{d}{dt}|_{t=0}f\left(e^{tX}z\right)\text{.}$$
Now let $`z(t)`$ be the curve in $`^2`$ defined as $`z(t)=e^{tX}z`$, so that $`z(0)=z`$. Of course, $`z(t)`$ can be written as $`z(t)=(z_1(t),z_2(t))`$, with $`z_i(t)`$. By the chain rule,
$$\pi _m(X)f=\frac{f}{z_1}\frac{dz_1}{dt}|_{t=0}+\frac{f}{z_2}\frac{dz_2}{dt}|_{t=0}\text{.}$$
But $`dz/dt|_{t=0}=Xz`$, so we obtain the following formula for $`\pi _m(X)`$
(5.3)
$$\pi _m(X)f=\frac{f}{z_1}\left(X_{11}z_1+X_{12}z_2\right)\frac{f}{z_2}\left(X_{21}z_1+X_{22}z_2\right)\text{.}$$
Now, according to Proposition 5.5, every finite-dimensional complex representation of the Lie algebra $`\mathrm{𝗌𝗎}(2)`$ extends uniquely to a complex-linear representation of the complexification of $`\mathrm{𝗌𝗎}(2)`$. But the complexification of $`\mathrm{𝗌𝗎}(2)`$ is (isomorphic to) $`\mathrm{𝗌𝗅}(2;)`$ (Chapter 3, Proposition 3.36). To see that this is so, note that $`\mathrm{𝗌𝗅}(2;)`$ is the space of all $`2\times 2`$ complex matrices with trace zero. But if $`X`$ is in $`\mathrm{𝗌𝗅}(2;)`$, then
$$X=\frac{XX^{}}{2}+\frac{X+X^{}}{2}=\frac{XX^{}}{2}+i\frac{X+X^{}}{2i}$$
where both $`(XX^{})/2`$ and $`(X+X^{})/2i`$ are in $`\mathrm{𝗌𝗎}(2)`$. (Check!) It is easy to see that this decomposition is unique, so that every $`X\mathrm{𝗌𝗅}(2;)`$ can be written uniquely as $`X=X_1+iY_1`$ with $`X_1,Y_1\mathrm{𝗌𝗎}(2)`$. Thus $`\mathrm{𝗌𝗅}(2;)`$ is isomorphic as a vector space to $`\mathrm{𝗌𝗎}(2)_𝐂`$. But this is in fact an isomorphism of Lie algebras, since in both cases
$$[X_1+iY_1,X_2+iY_2]=[X_1,X_2][Y_1,Y_2]+i\left([X_1,Y_2]+[X_2,Y_1]\right)\text{.}$$
(See Exercise 5.)
So, the representation $`\pi _m`$ of $`\mathrm{𝗌𝗎}(2)`$ given by (5.3) extends to a representation of $`\mathrm{𝗌𝗅}(2;)`$, which we will also call $`\pi _m`$. I assert that in fact formula (5.3), still holds for $`X\mathrm{𝗌𝗅}(2;)`$. Why is this? Well, (5.3) is undoubtedly (complex) linear, and it agrees with the original $`\pi _m`$ for $`X\mathrm{𝗌𝗎}(2)`$. But there is only one complex linear extension of $`\pi _m`$ from $`\mathrm{𝗌𝗎}(2)`$ to $`\mathrm{𝗌𝗅}(2;)`$, so this must be it!
So, for example, consider the element
$$H=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)$$
in the Lie algebra $`\mathrm{𝗌𝗅}(2;)`$. Applying formula (5.3) gives
$$\left(\pi _m(H)f\right)(z)=\frac{f}{z_1}z_1+\frac{f}{z_2}z_2\text{.}$$
Thus we see that
(5.4)
$$\pi _m(H)=z_1\frac{}{z_1}+z_2\frac{}{z_2}\text{.}$$
Applying $`\pi _m(H)`$ to a basis element $`z_1^kz_2^{mk}`$ we get
$$\pi _m(H)z_1^kz_2^{mk}=kz_1^kz_2^{mk}+(mk)z_1^kz_2^{mk}=(m2k)z_1^kz_2^{mk}\text{.}$$
Thus $`z_1^kz_2^{mk}`$ is an eigenvector for $`\pi _m(H)`$ with eigenvalue $`(m2k)`$. In particular, $`\pi _m(H)`$ is diagonalizable.
Let $`X`$ and $`Y`$ be the elements
$$\begin{array}{cc}X=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right);& Y=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)\end{array}$$
in $`\mathrm{𝗌𝗅}(2;)`$. Then (5.3) tells us that
$$\begin{array}{cc}\pi _m(X)=z_2\frac{}{z_1};& \pi _m(Y)=z_1\frac{}{z_2}\end{array}$$
so that
$`\pi _m(X)z_1^kz_2^{mk}`$ $`=kz_1^{k1}z_2^{mk+1}`$
(5.5) $`\pi _m(Y)z_1^kz_2^{mk}`$ $`=(km)z_1^{k+1}z_2^{mk1}\text{.}`$
###### Proposition 5.7.
The representation $`\pi _m`$ is an irreducible representation of $`\mathrm{𝗌𝗅}(2;)`$.
###### Proof.
It suffices to show that every non-zero invariant subspace of $`V_m`$ is in fact equal to $`V_m`$. So let $`W`$ be such a space. Since $`W`$ is assumed non-zero, there is at least one non-zero element $`w`$ in $`W`$. Then $`w`$ can be written uniquely in the form
$$w=a_0z_1^m+a_1z_1^{m1}z_2+a_2z_1^{m2}z_2^2\mathrm{}+a_mz_2^m$$
with at least one of the $`a_k`$’s non-zero. Let $`k_0`$ be the largest value of $`k`$ for which $`a_k0`$, and consider
$$\pi _m(X)^{k_0}w\text{.}$$
Since (by (5.5)) each application of $`\pi _m(X)`$ lowers the power of $`z_1`$ by 1, $`\pi _m(X)^{k_0}`$ will kill all the terms in $`w`$ whose power of $`z_1`$ is less than $`k_0`$, that is, all except the $`a_{k_0}z_1^{k_0}z_2^{mk_0}`$ term. On the other hand, we compute easily that
$$\pi _m(X)^{k_0}\left(a_{k_0}z_1^{k_0}z_2^{mk_0}\right)=k_0!(1)^{k_0}a_{k_0}z_2^m\text{.}$$
We see, then, that $`\pi _m(X)^{k_0}w`$ is a non-zero multiple of $`z_2^m`$. Since $`W`$ is assumed invariant, $`W`$ must contain this multiple of $`z_2^m`$, and so also $`z_2^m`$ itself.
But now it follows from (5.5) that $`\pi _m(Y)^kz_2^m`$ is a non-zero multiple of $`z_1^kz_2^{mk}`$. Therefore $`W`$ must also contain $`z_1^kz_2^{mk}`$ for all $`0km`$. Since these elements form a basis for $`V_m`$, we see that in fact $`W=V_m`$, as desired. ∎
#### 5.3.5. Two Unitary Representations of $`\mathrm{𝖲𝖮}(3)`$
Let $`=L^2(^3,dx)`$. For each $`R\mathrm{𝖲𝖮}(3)`$, define an operator $`\mathrm{\Pi }_1(R)`$ on $``$ by the formula
$$\left[\mathrm{\Pi }_1(R)f\right](x)=f\left(R^1x\right)\text{.}$$
Since Lebesgue measure $`dx`$ is rotationally invariant, $`\mathrm{\Pi }_1(R)`$ is a unitary operator for each $`R\mathrm{𝖲𝖮}(3)`$. The calculation of the previous subsection shows that the map $`R\mathrm{\Pi }_1(R)`$ is a homomorphism of $`\mathrm{𝖲𝖮}(3)`$ into $`U()`$. This map is strongly continuous, and hence constitutes a unitary representation of $`\mathrm{𝖲𝖮}(3)`$.
Similarly, we may consider the unit sphere $`S^2^3`$, with the usual surface measure $`\mathrm{\Omega }`$. Of course, any $`R\mathrm{𝖲𝖮}(3)`$ maps $`S^2`$ into $`S^2`$. For each $`R`$ we can define $`\mathrm{\Pi }_2(R)`$ acting on $`L^2(S^2,d\mathrm{\Omega })`$ by
$$\left[\mathrm{\Pi }_2(R)f\right](x)=f\left(R^1x\right)\text{.}$$
Then $`\mathrm{\Pi }_2`$ is a unitary representation of $`\mathrm{𝖲𝖮}(3)`$.
Neither of the unitary representations $`\mathrm{\Pi }_1`$ and $`\mathrm{\Pi }_2`$ is irreducible. In the case of $`\mathrm{\Pi }_2`$, $`L^2(S^2,d\mathrm{\Omega })`$ has a very nice decomposition as the orthogonal direct sum of finite-dimensional invariant subspaces. This decomposition is the theory of “spherical harmonics,” which are well known in the physics (and mathematics) literature.
#### 5.3.6. A Unitary Representation of the Reals
Let $`=L^2(,dx)`$. For each $`a`$, define $`T_a:`$ by
$$\left(T_af\right)(x)=f(xa)\text{.}$$
Clearly $`T_a`$ is a unitary operator for each $`a`$, and clearly $`T_aT_b=T_{a+b}`$. The map $`aT_a`$ is strongly continuous, so $`T`$ is a unitary representation of $``$. This representation is not irreducible. The theory of the Fourier transform allows you to determine all the closed, invariant subspaces of $``$ (W. Rudin, Real and Complex Analysis, Theorem 9.17).
#### 5.3.7. The Unitary Representations of the Real Heisenberg Group
Consider the Heisenberg group
$$H=\left\{\left(\begin{array}{ccc}1& a& b\\ 0& 1& c\\ 0& 0& 1\end{array}\right)\right|a,b,c\}\text{.}$$
Now consider a real, non-zero constant, which for reasons of historical convention we will call $`\mathrm{}`$ (“aitch-bar”). Now for each $`\mathrm{}\backslash \{0\}`$, define a unitary operator $`\mathrm{\Pi }_{\mathrm{}}`$ on $`L^2(,dx)`$ by
(5.6)
$$\mathrm{\Pi }_{\mathrm{}}\left(\begin{array}{ccc}1& a& b\\ 0& 1& c\\ 0& 0& 1\end{array}\right)f=e^{i\mathrm{}b}e^{i\mathrm{}cx}f(xa)\text{.}$$
It is clear that the right side of (5.6) has the same norm as $`f`$, so $`\mathrm{\Pi }_{\mathrm{}}`$ is indeed unitary.
Now compute
$`\mathrm{\Pi }_{\mathrm{}}\left(\begin{array}{ccc}1& \stackrel{~}{a}& \stackrel{~}{b}\\ 0& 1& \stackrel{~}{c}\\ 0& 0& 1\end{array}\right)\mathrm{\Pi }_{\mathrm{}}\left(\begin{array}{ccc}1& a& b\\ 0& 1& c\\ 0& 0& 1\end{array}\right)f`$
$`=e^{i\mathrm{}\stackrel{~}{b}}e^{i\mathrm{}\stackrel{~}{c}x}e^{i\mathrm{}b}e^{i\mathrm{}c(x\stackrel{~}{a})}f(x\stackrel{~}{a}a)`$
$`=e^{i\mathrm{}(\stackrel{~}{b}+b+c\stackrel{~}{a})}e^{i\mathrm{}(\stackrel{~}{c}+c)x}f\left(x(\stackrel{~}{a}+a)\right)\text{.}`$
This shows that the map $`A\mathrm{\Pi }_{\mathrm{}}(A)`$ is a homomorphism of the Heisenberg group into $`U\left(L^2()\right)`$. This map is strongly continuous, and so $`\mathrm{\Pi }_{\mathrm{}}`$ is a unitary representation of $`H`$.
Note that a typical unitary operator $`\mathrm{\Pi }_{\mathrm{}}(A)`$ consists of first translating $`f`$, then multiplying $`f`$ by the function $`e^{i\mathrm{}cx}`$, and then multiplying $`f`$ by the constant $`e^{i\mathrm{}b}`$. Multiplying $`f`$ by the function $`e^{i\mathrm{}cx}`$ has the effect of translating the Fourier transform of $`f`$, or in physical language, “translating $`f`$ in momentum space.” Now, if $`U_1`$ is an ordinary translation and $`U_2`$ is a translation of the Fourier transform (i.e., $`U_2=`$ multiplication by some $`e^{i\mathrm{}cx}`$), then $`U_1`$ and $`U_2`$ will not commute, but $`U_1U_2U_1^1U_2^1`$ will be simply multiplication by a constant of absolute value one. Thus $`\left\{\mathrm{\Pi }_{\mathrm{}}(A)\right|AH\}`$ is the group of operators on $`L^2()`$ generated by ordinary translations and translations in Fourier space. It is this representation of the Heisenberg group which motivates its name. (See also Exercise 10.)
It follows fairly easily from standard Fourier transform theory (e.g., W. Rudin, Real and Complex Analysis, Theorem 9.17) that for each $`\mathrm{}\backslash \{0\}`$ the representation $`\mathrm{\Pi }_{\mathrm{}}`$ is irreducible. Furthermore, these are (up to equivalence) almost all of the irreducible unitary representations of $`H`$. The only remaining ones are the one-dimensional representations $`\mathrm{\Pi }_{\alpha ,\beta }`$
$$\mathrm{\Pi }_{\alpha ,\beta }\left(\begin{array}{ccc}1& a& b\\ 0& 1& c\\ 0& 0& 1\end{array}\right)=e^{i(\alpha a+\beta c)}I$$
with $`\alpha ,\beta `$. (The $`\mathrm{\Pi }_{\alpha ,\beta }`$’s are the irreducible unitary representations in which the center of $`H`$ acts trivially.) The fact that $`\mathrm{\Pi }_{\mathrm{}}`$’s and the $`\mathrm{\Pi }_{\alpha ,\beta }`$’s are all of the (strongly continuous) irreducible unitary representations of $`H`$ is closely related to the celebrated Stone-Von Neumann theorem in mathematical physics. See, for example, M. Reed and B. Simon, Methods of Modern Mathematical Physics, Vol. 3, Theorem XI.84. See also Exercise 11.
### 5.4. The Irreducible Representations of $`\mathrm{𝗌𝗎}(2)`$
In this section we will compute (up to equivalence) all the finite-dimensional irreducible complex representations of the Lie algebra $`\mathrm{𝗌𝗎}(2)`$. This computation is important for several reasons. In the first place, $`\mathrm{𝗌𝗎}(2)\mathrm{𝗌𝗈}(3)`$, and the representations of $`\mathrm{𝗌𝗈}(3)`$ are of physical significance. (The computation we will do here is found in every standard textbook on quantum mechanics, under the heading “angular momentum.”) In the second place, the representation theory of $`\mathrm{𝗌𝗎}(2)`$ is an illuminating example of how one uses commutation relations to determine the representations of a Lie algebra. In the third place, in determining the representations of general semisimple Lie algebras (Chapter 6), we will explicitly use the representation theory of $`\mathrm{𝗌𝗎}(2)`$.
Now, every finite-dimensional complex representation $`\pi `$ of $`\mathrm{𝗌𝗎}(2)`$ extends by Prop. 5.5 to a complex-linear representation (also called $`\pi `$) of the complexification of $`\mathrm{𝗌𝗎}(2)`$, namely $`\mathrm{𝗌𝗅}(2;)`$.
###### Proposition 5.8.
Let $`\pi `$ be a complex representation of $`\mathrm{𝗌𝗎}(2)`$, extended to a complex-linear representation of $`\mathrm{𝗌𝗅}(2;)`$. Then $`\pi `$ is irreducible as a representation of $`\mathrm{𝗌𝗎}(2)`$ if and only if it is irreducible as a representation of $`\mathrm{𝗌𝗅}(2;)`$.
###### Proof.
Let us make sure we are clear about what this means. Suppose that $`\pi `$ is a complex representation of the (real) Lie algebra $`\mathrm{𝗌𝗎}(2)`$, acting on the complex space $`V`$. Then saying that $`\pi `$ is irreducible means that there is no non-trivial invariant complex subspace $`WV`$. That is, even though $`\mathrm{𝗌𝗎}(2)`$ is a real Lie algebra, when considering complex representations we are interested only in complex invariant subspaces.
Now, suppose that $`\pi `$ is irreducible as a representation of $`\mathrm{𝗌𝗎}(2)`$. If $`W`$ is a (complex) subspace of $`V`$ which is invariant under $`\mathrm{𝗌𝗅}(2;)`$, then certainly $`W`$ is invariant under $`\mathrm{𝗌𝗎}(2)\mathrm{𝗌𝗅}(2;)`$. Therefore $`W=\{0\}`$ or $`W=V`$. Thus $`\pi `$ is irreducible as a representation of $`\mathrm{𝗌𝗅}(2;)`$.
On the other hand, suppose that $`\pi `$ is irreducible as a representation of $`\mathrm{𝗌𝗅}(2;)`$, and suppose that $`W`$ is a (complex) subspace of $`V`$ which is invariant under $`\mathrm{𝗌𝗎}(2)`$. Then $`W`$ will also be invariant under $`\pi (X+iY)=\pi (X)+i\pi (Y)`$, for all $`X,Y\mathrm{𝗌𝗎}(2)`$. Since every element of $`\mathrm{𝗌𝗅}(2;)`$ can be written as $`X+iY`$, we conclude that in fact $`W`$ is invariant under $`\mathrm{𝗌𝗅}(2;)`$. Thus $`W=\{0\}`$ or $`W=V`$, so $`\pi `$ is irreducible as a representation of $`\mathrm{𝗌𝗎}(2)`$. ∎
We see, then that studying the irreducible representations of $`\mathrm{𝗌𝗎}(2)`$ is equivalent to studying the irreducible representations of $`\mathrm{𝗌𝗅}(2;)`$. Passing to the complexified Lie algebra makes our computations easier.
We will use the following basis for $`\mathrm{𝗌𝗅}(2;)`$:
$$\begin{array}{ccc}H=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right);& X=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right);& Y=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)\end{array}$$
which have the commutation relations
$$\begin{array}{ccc}\hfill [H,X]& \hfill =& \hfill 2X\\ \hfill [H,Y]& \hfill =& \hfill 2Y\\ \hfill [X,Y]& \hfill =& \hfill H\end{array}\text{.}$$
If $`V`$ is a (finite-dimensional complex) vector space, and $`A,B,`$ and $`C`$ are operators on $`V`$ satisfying
$$\begin{array}{ccc}\hfill [A,B]& \hfill =& \hfill 2B\\ \hfill [A,C]& \hfill =& \hfill 2C\\ \hfill [B,C]& \hfill =& \hfill A\end{array}$$
then because of the skew-symmetry and bilinearity of brackets, the linear map $`\pi :\mathrm{𝗌𝗅}(2;)\mathrm{𝗀𝗅}(V)`$ satisfying
$$\begin{array}{ccc}\pi (H)=A;& \pi (X)=B;& \pi (Y)=C\end{array}$$
will be a representation of $`\mathrm{𝗌𝗅}(2;)`$.
###### Theorem 5.9.
For each integer $`m0`$, there is an irreducible representation of $`\mathrm{𝗌𝗅}(2;)`$ with dimension $`m+1`$. Any two irreducible representations of $`\mathrm{𝗌𝗅}(2;)`$ with the same dimension are equivalent. If $`\pi `$ is an irreducible representation of $`\mathrm{𝗌𝗅}(2;)`$ with dimension $`m+1`$, then $`\pi `$ is equivalent to the representation $`\pi _m`$ described in Section 5.3.
###### Proof.
Let $`\pi `$ be an irreducible representation of $`\mathrm{𝗌𝗅}(2;)`$ acting on a (finite-dimensional complex) space $`V`$. Our strategy is to diagonalize the operator $`\pi (H)`$. Of course, a priori, we don’t know that $`\pi (H)`$ is diagonalizable. However, because we are working over the (algebraically closed) field of complex numbers, $`\pi (H)`$ must have at least one eigenvector. ∎
###### Proof.
The following lemma is the key to the entire proof.
###### Lemma 5.10.
Let $`u`$ be an eigenvector of $`\pi (H)`$ with eigenvalue $`\alpha `$. Then
$$\pi (H)\pi (X)u=(\alpha +2)\pi (X)u\text{.}$$
Thus either $`\pi (X)u=0`$, or else $`\pi (X)u`$ is an eigenvector for $`\pi (H)`$ with eigenvalue $`\alpha +2`$. Similarly,
$$\pi (H)\pi (Y)u=(\alpha 2)\pi (Y)u$$
so that either $`\pi (Y)u=0`$, or else $`\pi (Y)u`$ is an eigenvector for $`\pi (H)`$ with eigenvalue $`\alpha 2`$.
###### Proof.
We call $`\pi (X)`$ the “raising operator,” because it has the effect of raising the eigenvalue of $`\pi (H)`$ by 2, and we call $`\pi (Y)`$ the “lowering operator.” We know that $`[\pi (H),\pi (X)]=\pi \left([H,X]\right)=2\pi (X)`$. Thus
$$\pi (H)\pi (X)\pi (X)\pi (H)=2\pi (X)$$
or
$$\pi (H)\pi (X)=\pi (X)\pi (H)+2\pi (X)\text{.}$$
Thus
$`\pi (H)\pi (X)u=\pi (X)\pi (H)u+2\pi (X)u`$
$`=\pi (X)\left(\alpha u\right)+2\pi (X)u`$
$`=(\alpha +2)\pi (X)u\text{.}`$
Similarly, $`[\pi (H),\pi (Y)]=2\pi (Y)`$, and so
$$\pi (H)\pi (Y)=\pi (Y)\pi (H)2\pi (Y)$$
so that
$`\pi (H)\pi (Y)u=\pi (Y)\pi (H)u2\pi (Y)u`$
$`=\pi (Y)\left(\alpha u\right)2\pi (Y)u`$
$`=(\alpha 2)\pi (Y)u\text{.}`$
This is what we wanted to show. ∎
As we have observed, $`\pi (H)`$ must have at least one eigenvector $`u`$ ($`u0`$), with some eigenvalue $`\alpha `$. By the lemma,
$$\pi (H)\pi (X)u=(\alpha +2)\pi (X)u$$
and more generally
$$\pi (H)\pi (X)^nu=(\alpha +2n)\pi (X)^nu\text{.}$$
This means that either $`\pi (X)^nu=0`$, or else $`\pi (X)^nu`$ is an eigenvector for $`\pi (H)`$ with eigenvalue $`(\alpha +2n)`$.
Now, an operator on a finite-dimensional space can have only finitely many distinct eigenvalues. Thus the $`\pi (X)^nu`$’s cannot all be different from zero. Thus there is some $`N0`$ such that
$$\pi (X)^Nu0$$
but
$$\pi (X)^{N+1}u=0\text{.}$$
Define $`u_0=\pi (X)^Nu`$ and $`\lambda =\alpha +2N`$. Then
(5.7) $`\pi (H)u_0=\lambda u_0`$
(5.8) $`\pi (X)u_0=0`$
Then define
$$u_k=\pi (Y)^ku_0$$
for $`k0`$. By the second part of the lemma, we have
(5.9)
$$\pi (H)u_k=\left(\lambda 2k\right)u_k\text{.}$$
Since, again, $`\pi (H)`$ can have only finitely many eigenvalues, the $`u_k`$’s cannot all be non-zero.
###### Lemma 5.11.
With the above notation,
$`\pi (X)u_k=\left[k\lambda k(k1)\right]u_{k1}\text{ }(k>0)`$
$`\pi \left(X\right)u_0=0.`$
###### Proof.
We proceed by induction on $`k`$. In the case $`k=1`$ we note that $`u_1=\pi (Y)u_0`$. Using the commutation relation $`[\pi (X),\pi (Y)]=\pi (H)`$ we have
$$\pi (X)u_1=\pi (X)\pi (Y)u_0=\left(\pi (Y)\pi (X)+\pi (H)\right)u_0\text{.}$$
But $`\pi (X)u_0=0`$, so we get
$$\pi (X)u_1=\lambda u_0$$
which is the lemma in the case $`k=1`$.
Now, by definition $`u_{k+1}=\pi (Y)u_k`$. Using (5.9) and induction we have
$`\pi (X)u_{k+1}=\pi (X)\pi (Y)u_k`$
$`=\left(\pi (Y)\pi (X)+\pi (H)\right)u_k`$
$`=\pi (Y)\left[k\lambda k(k1)\right]u_{k1}+(\lambda 2k)u_k`$
$`=\left[k\lambda k(k1)+(\lambda 2k)\right]u_k\text{.}`$
Simplifying the last expression give the Lemma. ∎
Since $`\pi (H)`$ can have only finitely many eigenvalues, the $`u_k`$’s cannot all be non-zero. There must therefore be an integer $`m0`$ such that
$$u_k=\pi (Y)^ku_00$$
for all $`km`$, but
$$u_{m+1}=\pi (Y)^{m+1}u_0=0\text{.}$$
Now if $`u_{m+1}=0`$, then certainly $`\pi (X)u_{m+1}=0`$. Then by Lemma 5.11,
$$0=\pi (X)u_{m+1}=\left[(m+1)\lambda m(m+1)\right]u_m=(m+1)(\lambda m)u_m\text{.}$$
But $`u_m0`$, and $`m+10`$ (since $`m0`$). Thus in order to have $`(m+1)(\lambda m)u_m`$ equal to zero, we must have $`\lambda =m`$.
We have made considerable progress. Given a finite-dimensional irreducible representation $`\pi `$ of $`\mathrm{𝗌𝗅}(2;)`$, acting on a space $`V`$, there exists an integer $`m0`$ and non-zero vectors $`u_0,\mathrm{}u_m`$ such that (putting $`\lambda `$ equal to $`m`$)
$`\pi (H)u_k=(m2k)u_k`$
$`\pi (Y)u_k=u_{k+1}(k<m)`$
$`\pi (Y)u_m=0`$
$`\pi (X)u_k=\left[kmk(k1)\right]u_{k1}\text{ }(k>0)`$
(5.10) $`\pi (X)u_0=0`$
The vectors $`u_0,\mathrm{}u_m`$ must be linearly independent, since they are eigenvectors of $`\pi (H)`$ with distinct eigenvalues. Moreover, the $`(m+1)`$-dimensional span of $`u_0,\mathrm{}u_m`$ is explicitly invariant under $`\pi (H)`$, $`\pi (X)`$, and $`\pi (Y)`$, and hence under $`\pi (Z)`$ for all $`Z\mathrm{𝗌𝗅}(2;)`$. Since $`\pi `$ is irreducible, this space must be all of $`V`$.
We have now shown that every irreducible representation of $`\mathrm{𝗌𝗅}(2;)`$ is of the form (5.10). It remains to show that everything of the form (5.10) is a representation, and that it is irreducible. That is, if we define $`\pi (H)`$, $`\pi (X)`$, and $`\pi (Y)`$ by (5.10) (where the $`u_k`$’s are basis elements for some $`(m+1)`$-dimensional vector space), then we want to show that they have the right commutation relations to form a representation of $`\mathrm{𝗌𝗅}(2;)`$, and that this representation is irreducible.
The computation of the commutation relations of $`\pi (H)`$, $`\pi (X)`$, and $`\pi (Y)`$ is straightforward, and is left as an exercise. Note that when dealing with $`\pi (Y)`$, you should treat separately the vectors $`u_k`$, $`k<m`$, and $`u_m`$. Irreducibility is also easy to check, by imitating the proof of Proposition 5.7. (See Exercise 6.)
We have now shown that there is an irreducible representation of $`\mathrm{𝗌𝗅}(2;)`$ in each dimension $`m+1`$, by explicitly writing down how $`H`$, $`X`$, and $`Y`$ should act (Equation 5.10) in a basis. But we have shown more than this. We also have shown that any $`(m+1)`$-dimensional irreducible representation of $`\mathrm{𝗌𝗅}(2;)`$ must be of the form (5.10). It follows that any two irreducible representations of $`\mathrm{𝗌𝗅}(2;)`$ of dimension $`(m+1)`$ must be equivalent. For if $`\pi _1`$ and $`\pi _2`$ are two irreducible representations of dimension $`(m+1)`$, acting on spaces $`V_1`$ and $`V_2`$, then $`V_1`$ has a basis $`u_0,\mathrm{}u_m`$ as in (5.10) and $`V_2`$ has a similar basis $`\stackrel{~}{u}_0,\mathrm{}\stackrel{~}{u}_m`$. But then the map $`\varphi :V_1V_2`$ which sends $`u_k`$ to $`\stackrel{~}{u}_k`$ will be an isomorphism of representations. (Think about it.)
In particular, the $`(m+1)`$-dimensional representation $`\pi _m`$ described in Section 5.3 must be equivalent to (5.10).This can be seen explicitly by introducing the following basis for $`V_m`$:
$$u_k=\left[\pi _m(Y)\right]^k(z_2^m)=(1)^k\frac{m!}{(mk)!}z_1^kz_2^{mk}(km)\text{.}$$
Then by definition $`\pi _m(Y)u_k=u_{k+1}`$ ($`k<m`$), and it is clear that $`\pi _m(Y)u_m=0`$. It is easy to see that $`\pi _m(H)u_k=(m2k)u_k`$. The only thing left to check is the behavior of $`\pi _m(X)`$. But direct computation shows that
$$\pi _m(X)u_k=k(mk+1)u_{k1}=\left[kmk(k1)\right]u_{k1}\text{.}$$
as required.
This completes the proof of Theorem 5.9. ∎
### 5.5. Direct Sums of Representations and Complete Reducibility
One way of generating representations is to take some representations you know and combine them in some fashion. We will consider two methods of generating new representations from old ones, namely direct sums and tensor products of representations. In this section we consider direct sums; in the next section we look at tensor products. (There is one other standard construction of this sort, namely the dual of a representation. See Exercise 14.)
###### Definition 5.12.
Let $`G`$ be a matrix Lie group, and let $`\mathrm{\Pi }_1,\mathrm{\Pi }_2,\mathrm{}\mathrm{\Pi }_n`$ be representations of $`G`$ acting on vector spaces $`V_1,V_2,\mathrm{}V_n`$. Then the direct sum of $`\mathrm{\Pi }_1,\mathrm{\Pi }_2,\mathrm{}\mathrm{\Pi }_n`$ is a representation $`\mathrm{\Pi }_1\mathrm{}\mathrm{\Pi }_n`$ of $`G`$ acting on the space $`V_1\mathrm{}V_n`$, defined by
$$\left[\mathrm{\Pi }_1\mathrm{}\mathrm{\Pi }_n(A)\right]\left(v_{1,}\mathrm{}v_n\right)=(\mathrm{\Pi }_1(A)v_1,\mathrm{},\mathrm{\Pi }_n(A)v_n)$$
for all $`AG`$.
Similarly, if $`𝔤`$ is a Lie algebra, and $`\pi _1,\pi _2,\mathrm{}\pi _n`$ are representations of $`𝔤`$ acting on $`V_1,V_2,\mathrm{}V_n`$, then we define the direct sum of $`\pi _1,\pi _2,\mathrm{}\pi _n`$, acting on $`V_1\mathrm{}V_n`$ by
$$\left[\pi _1\mathrm{}\pi _n(X)\right]\left(v_{1,}\mathrm{}v_n\right)=(\pi _1(X)v_1,\mathrm{},\pi _n(X)v_n)$$
for all $`X𝔤`$.
It is trivial to check that, say, $`\mathrm{\Pi }_1\mathrm{}\mathrm{\Pi }_n`$ is really a representation of $`G`$.
###### Definition 5.13.
A finite-dimensional representation of a group or Lie algebra, acting on a space $`V`$, is said to be completely reducible if the following property is satisfied: Given an invariant subspace $`WV`$, and a second invariant subspace $`UWV`$, there exists a third invariant subspace $`\stackrel{~}{U}W`$ such that $`U\stackrel{~}{U}=\left\{0\right\}`$ and $`U+\stackrel{~}{U}=W`$.
The following Proposition shows that complete reducibility is a nice property for a representation to have.
###### Proposition 5.14.
A finite-dimensional completely reducible representation of a group or Lie algebra is equivalent to a direct sum of (one or more) irreducible representations.
###### Proof.
The proof is by induction on the dimension of the space $`V`$. If $`dimV=1`$, then automatically the representation is irreducible, since then $`V`$ is has no non-trivial subspaces, let alone non-trivial invariant subspaces. Thus $`V`$ is a direct sum of irreducible representations, with just one summand, namely $`V`$ itself.
Suppose, then, that the Proposition holds for all representations with dimension strictly less than $`n`$, and that $`dimV=n`$. If $`V`$ is irreducible, then again we have a direct sum with only one summand, and we are done. If $`V`$ is not irreducible, then there exists a non-trivial invariant subspace $`UV`$. Taking $`W=V`$ in the definition of complete reducibility, we see that there is another invariant subspace $`\stackrel{~}{U}`$ with $`U\stackrel{~}{U}=\left\{0\right\}`$ and $`U+\stackrel{~}{U}=V`$. That is, $`VU\stackrel{~}{U}`$ as a vector space.
But since $`U`$ and $`\stackrel{~}{U}`$ are invariant, they can be viewed as representations in their own right. (That is, the action of our group or Lie algebra on $`U`$ or $`\stackrel{~}{U}`$ is a representation.) It is easy to see that in fact $`V`$ is isomorphic to $`U\stackrel{~}{U}`$ as a representation. Furthermore, it is easy to see that both $`U`$ and $`\stackrel{~}{U}`$ are completely reducible representations, since every invariant subspace $`W`$ of, say, $`U`$ is also an invariant subspace of $`V`$. But since $`U`$ is non-trivial (i.e., $`U\left\{0\right\}`$ and $`UV`$), we have $`dimU<dimV`$ and $`dim\stackrel{~}{U}<dimV`$. Thus by induction $`UU_1\mathrm{}U_n`$ (as representations), with the $`U_i`$’s irreducible, and $`\stackrel{~}{U}\stackrel{~}{U}_1\mathrm{}\stackrel{~}{U}_m`$, with the $`\stackrel{~}{U}_i`$’s irreducible, so that $`VU_1\mathrm{}U_n\stackrel{~}{U}_1\mathrm{}\stackrel{~}{U}_m`$. ∎
Certain groups and Lie algebras have the property that every (finite-dimensional) representation is completely reducible. This is a very nice property, because it implies (by the above Proposition) that every representation is equivalent to a direct sum of irreducible representations. (And, as it turns out, this decomposition is essentially unique.) Thus for such groups and Lie algebras, if you know (up to equivalence) what all the irreducible representations are, then you know (up to equivalence) what all the representations are.
Unfortunately, not every representation is irreducible. For example, the standard representation of the Heisenberg group is not completely reducible. (See Exercise 8.)
###### Proposition 5.15.
Let $`G`$ be a matrix Lie group. Let $`\mathrm{\Pi }`$ be a finite-dimensional unitary representation of $`G`$, acting on a finite-dimensional real or complex Hilbert space $`V`$. Then $`\mathrm{\Pi }`$ is completely reducible.
###### Proof.
So, we are assuming that our space $`V`$ is equipped with an inner product, and that $`\mathrm{\Pi }(A)`$ is unitary for each $`AG`$. Suppose that $`WV`$ is invariant, and that $`UWV`$ is also invariant. Define
$$\stackrel{~}{U}=U^{}W\text{.}$$
Then of course $`\stackrel{~}{U}U=\left\{0\right\}`$, and standard Hilbert space theory implies that $`\stackrel{~}{U}+U=W`$.
It remains only to show that $`\stackrel{~}{U}`$ is invariant. So suppose that $`vU^{}W`$. Since $`W`$ is assumed invariant, $`\mathrm{\Pi }(A)v`$ will be in $`W`$ for any $`AG`$. We need to show that $`\mathrm{\Pi }(A)v`$ is perpendicular to $`U`$. Well, since $`\mathrm{\Pi }(A^1)`$ is unitary, then for any $`uU`$
$$u,\mathrm{\Pi }(A)v=\mathrm{\Pi }(A^1)u,\mathrm{\Pi }(A^1)\mathrm{\Pi }(A)v=\mathrm{\Pi }(A^1)u,v\text{.}$$
But $`U`$ is assumed invariant, and so $`\mathrm{\Pi }(A^1)uU`$. But then since $`vU^{}`$, $`\mathrm{\Pi }(A^1)u,v=0`$. This means that
$$u,\mathrm{\Pi }(A)v=0$$
for all $`uU`$, i.e., $`\mathrm{\Pi }(A)vU^{}`$.
Thus $`\stackrel{~}{U}`$ is invariant, and we are done. ∎
###### Proposition 5.16.
If $`G`$ is a finite group, then every finite-dimensional real or complex representation of $`G`$ is completely reducible.
###### Proof.
Suppose that $`\mathrm{\Pi }`$ is a representation of $`G`$, acting on a space $`V`$. Choose an arbitrary inner product $``$ on $`V`$. Then define a new inner product $`_G`$ on $`V`$ by
$$v_1,v_2_G=\underset{gG}{}\mathrm{\Pi }(g)v_1,\mathrm{\Pi }(g)v_2\text{.}$$
It is very easy to check that indeed $`_G`$ is an inner product. Furthermore, if $`hG`$, then
$`\mathrm{\Pi }(h)v_1,\mathrm{\Pi }(h)v_2_G={\displaystyle \underset{gG}{}}\mathrm{\Pi }(g)\mathrm{\Pi }(h)v_1,\mathrm{\Pi }(g)\mathrm{\Pi }(h)v_2`$
$`={\displaystyle \underset{gG}{}}\mathrm{\Pi }(gh)v_1,\mathrm{\Pi }(gh)v_2\text{.}`$
But as $`g`$ ranges over $`G`$, so does $`gh`$. Thus in fact
$$\mathrm{\Pi }(h)v_1,\mathrm{\Pi }(h)v_2_G=v_1,v_2_G\text{.}$$
That is, $`\mathrm{\Pi }`$ is a unitary representation with respect to the inner product $`_G`$. Thus $`\mathrm{\Pi }`$ is completely reducible by Proposition 5.15. ∎
There is a variant of the above argument which can be used to prove the following result:
###### Proposition 5.17.
If $`G`$ is a compact matrix Lie group, then every finite-dimensional real or complex representation of $`G`$ is completely reducible.
###### Proof.
This proof requires the notion of Haar measure. A left Haar measure on a matrix Lie group $`G`$ is a non-zero measure $`\mu `$ on the Borel $`\sigma `$-algebra in $`G`$ with the following two properties: 1) it is locally finite, that is, every point in $`G`$ has a neighborhood with finite measure, and 2) it is left-translation invariant. Left-translation invariance means that $`\mu \left(gE\right)=\mu \left(E\right)`$ for all $`gG`$ and for all Borel sets $`EG`$, where
$$gE=\left\{ge\right|eE\}\text{.}$$
It is a fact which we cannot prove here that every matrix Lie group has a left Haar measure, and that this measure is unique up to multiplication by a constant. (One can analogously define right Haar measure, and a similar theorem holds for it. Left Haar measure and right Haar measure may or may not coincide; a group for which they do is called unimodular.)
Now, the key fact for our purpose is that left Haar measure is finite if and only if the group $`G`$ is compact. So if $`\mathrm{\Pi }`$ is a finite-dimensional representation of a compact group $`G`$ acting on a space $`V`$, then let $``$ be an arbitrary inner product on $`V`$, and define a new inner product $`_G`$ on $`V`$ by
$$v_1,v_2_G=_G\mathrm{\Pi }(g)v_1,\mathrm{\Pi }(g)v_2𝑑\mu \left(g\right)\text{,}$$
where $`\mu `$ is left Haar measure. Again, it is easy to check that $`_G`$ is an inner product. Furthermore, if $`hG`$, then by the left-invariance of $`\mu `$
$`\mathrm{\Pi }(h)v_1,\mathrm{\Pi }(h)v_2_G`$ $`={\displaystyle _G}\mathrm{\Pi }(g)\mathrm{\Pi }(h)v_1,\mathrm{\Pi }(g)\mathrm{\Pi }(h)v_2𝑑\mu \left(g\right)`$
$`={\displaystyle _G}\mathrm{\Pi }(gh)v_1,\mathrm{\Pi }(gh)v_2𝑑\mu \left(g\right)`$
$`=v_1,v_2_G\text{.}`$
So $`\mathrm{\Pi }`$ is a unitary representation with respect to $`_G`$, and thus completely reducible. Note that $`_G`$ is well-defined only because $`\mu `$ is finite. ∎
### 5.6. Tensor Products of Representations
Let $`U`$ and $`V`$ be finite-dimensional real or complex vector spaces. We wish to define the tensor product of $`U`$ and $`V`$, which is will be a new vector space $`UV`$ “built” out of $`U`$ and $`V`$. We will discuss the idea of this first, and then give the precise definition.
We wish to consider a formal “product” of an element $`u`$ of $`U`$ with an element $`v`$ of $`V`$, denoted $`uv`$. The space $`UV`$ is then the space of linear combinations of such products, i.e., the space of elements of the form
(5.11)
$$a_1u_1v_1+a_2u_2v_2+\mathrm{}+a_nu_nv_n\text{.}$$
Of course, if “$``$” is to be interpreted as a product, then it should be bilinear. That is, we should have
$`\left(u_1+au_2\right)v`$ $`=u_1v+au_2v`$
$`u\left(v_1+av_2\right)`$ $`=uv_1+auv_2\text{.}`$
We do not assume that the product is commutative. (In fact, the product in the other order, $`vu`$, is in a different space, namely, $`VU`$.)
Now, if $`e_1,e_2,\mathrm{},e_n`$ is a basis for $`U`$ and $`f_1,f_2,\mathrm{},f_m`$ is a basis for $`V`$, then using bilinearity it is easy to see that any element of the form (5.11) can be written as a linear combination of the elements $`e_if_{j\text{ }}`$. In fact, it seems reasonable to expect that $`\left\{e_if_j\right|0in,0jm\}`$ should be a basis for the space $`UV`$. This in fact turns out to be the case.
###### Definition 5.18.
If $`U`$ and $`V`$ are finite-dimensional real or complex vector spaces, then a tensor product of $`U`$ with $`V`$ is a vector space $`W`$, together with a bilinear map $`\varphi :U\times VW`$ with the following property: If $`\psi `$ is any bilinear map of $`U\times V`$ into a vector space $`X`$, then there exists a unique linear map $`\stackrel{~}{\psi }`$ of $`W`$ into $`X`$ such that the following diagram commutes:
$$\begin{array}{ccc}U\times V& \stackrel{\mathit{\varphi }}{}& W\\ \psi & & \stackrel{~}{\psi }\\ & X& \end{array}\text{.}$$
Note that the bilinear map $`\psi `$ from $`U\times V`$ into $`X`$ turns into the linear map $`\stackrel{~}{\psi }`$ of $`W`$ into $`X`$. This is one of the points of tensor products: bilinear maps on $`U\times V`$ turn into linear maps on $`W`$.
###### Theorem 5.19.
If $`U`$ and $`V`$ are any finite-dimensional real or complex vector spaces, then a tensor product $`(W,\varphi )`$ exists. Furthermore, $`(W,\varphi )`$ is unique up to canonical isomorphism. That is, if $`(W_1,\varphi _1)`$ and $`(W_2,\varphi _2)`$ are two tensor products, then there exists a unique vector space isomorphism $`\mathrm{\Phi }:W_1W_2`$ such that the following diagram commutes
$$\begin{array}{ccc}U\times V& \stackrel{\varphi _1}{}& W_1\\ \varphi _2& & \mathrm{\Phi }\\ & W_2& \end{array}\text{.}$$
Suppose that $`(W,\varphi )`$ is a tensor product, and that $`e_1,e_2,\mathrm{},e_n`$ is a basis for $`U`$ and $`f_1,f_2,\mathrm{},f_m`$ is a basis for $`V`$. Then $`\left\{\varphi (e_i,f_j)\right|0in,0jm\}`$ is a basis for $`W`$.
###### Proof.
Exercise 12. ∎
###### Notation 5.20.
Since the tensor product of $`U`$ and $`V`$ is essentially unique, we will let $`UV`$ denote an arbitrary tensor product space, and we will write $`uv`$ instead of $`\varphi (u,v)`$. In this notation, the Theorem says that $`\left\{e_if_j\right|0in,0jm\}`$ is a basis for $`UV`$, as expected. Note in particular that
$$dim\left(UV\right)=\left(dimU\right)\left(dimV\right)$$
(not $`dimU+dimV`$).
The defining property of $`UV`$ is called the universal property of tensor products. While it may seem that we are taking a simple idea and making it confusing, in fact there is a point to this universal property. Suppose we want to define a linear map $`T`$ from $`UV`$ into some other space. The most sensible way to define this is to define $`T`$ on elements of the form $`uv`$. (You might try defining it on a basis, but this forces you to worry about whether things depend on the choice of basis.) Now, every element of $`UV`$ is a linear combination of things of the form $`uv`$. However, this representation is far from unique. (Since, say, if $`u=u_1+u_2`$, then you can rewrite $`uv`$ as $`u_1v+u_2v`$.)
Thus if you try to define $`T`$ by what it does to elements of the form $`uv`$, you have to worry about whether $`T`$ is well-defined. This is where the universal property comes in. Suppose that $`\psi (u,v)`$ is some bilinear expression in $`u,v`$. Then the universal property says precisely that there is a unique linear map $`T`$ ($`=\stackrel{~}{\psi }`$) such that
$$T(uv)=\psi (u,v)\text{.}$$
(Think about it and make sure that you see that this is really what the universal property says.)
The conclusion is this: You can define a linear map $`T`$ on $`UV`$ by defining it on elements of the form $`uv`$, and this will be well-defined, provided that $`T(uv)`$ is bilinear in $`(u,v)`$. The following Proposition shows how to make use of this idea.
###### Proposition 5.21.
Let $`U`$ and $`V`$ be finite-dimensional real or complex vector spaces. Let $`A:UU`$ and $`B:VV`$ be linear operators. Then there exists a unique linear operator from $`UV`$ to $`UV`$, denoted $`AB`$, such that
$$AB(uv)=\left(Au\right)\left(Bv\right)$$
for all $`uU`$, $`vV`$.
If $`A_1,A_2`$ are linear operators on $`U`$ and $`B_1,B_2`$ are linear operators on $`V`$, then
$$\left(A_1B_1\right)\left(A_2B_2\right)=\left(A_1A_2\right)\left(B_1B_2\right)\text{.}$$
###### Proof.
Define a map $`\psi `$ from $`U\times V`$ into $`UV`$ by
$$\psi (u,v)=\left(Au\right)\left(Bv\right)\text{.}$$
Since $`A`$ and $`B`$ are linear, and since $``$ is bilinear, $`\psi `$ will be a bilinear map of $`U\times V`$ into $`UV`$. But then the universal property says that there is an associated linear map $`\stackrel{~}{\psi }:UVUV`$ such that
$$\stackrel{~}{\psi }(uv)=\psi (u,v)=\left(Au\right)\left(Bv\right)\text{.}$$
Then $`\stackrel{~}{\psi }`$ is the desired map $`AB`$.
Now, if $`A_1,A_2`$ are operators on $`U`$ and $`B_1,B_2`$ are operators on $`V`$, then compute that
$`\left(A_1B_1\right)\left(A_2B_2\right)(uv)=\left(A_1B_1\right)\left(A_2uB_2v\right)`$
$`=A_1A_2uB_1B_2v\text{.}`$
This shows that $`\left(A_1B_1\right)\left(A_2B_2\right)=\left(A_1A_2\right)\left(B_1B_2\right)`$ are equal on elements of the form $`uv`$. Since every element of $`UV`$ can be written as a linear combination of things of the form $`uv`$ (in fact of $`e_if_j`$), $`\left(A_1B_1\right)\left(A_2B_2\right)`$ and $`\left(A_1A_2\right)\left(B_1B_2\right)`$ must be equal on the whole space. ∎
We are now ready to define tensor products of representations. There are two different approaches to this, both of which are important. The first approach starts with a representation of a group $`G`$ acting on a space $`V`$ and a representation of another group $`H`$ acting on a space $`U,`$ and produces a representation of the product group $`G\times H`$ acting on the space $`UV`$. The second approach starts with two different representations of the same group $`G`$, acting on spaces $`U`$ and $`V`$, and produces a representation of $`G`$ acting on $`UV`$. Both of these approaches can be adapted to apply to Lie algebras.
###### Definition 5.22.
Let $`G`$ and $`H`$ be matrix Lie groups. Let $`\mathrm{\Pi }_1`$ be a representation of $`G`$ acting on a space $`U`$ and let $`\mathrm{\Pi }_2`$ be a representation of $`H`$ acting on a space $`V`$. The the tensor product of $`\mathrm{\Pi }_1`$ and $`\mathrm{\Pi }_2`$ is a representation $`\mathrm{\Pi }_1\mathrm{\Pi }_2`$ of $`G\times H`$ acting on $`UV`$ defined by
$$\mathrm{\Pi }_1\mathrm{\Pi }_2(A,B)=\mathrm{\Pi }_1(A)\mathrm{\Pi }_2(B)$$
for all $`AG`$ and $`BH.`$
Using the above Proposition, it is very easy to check that indeed $`\mathrm{\Pi }_1\mathrm{\Pi }_2`$ is a representation of $`G\times H`$.
Now, if $`G`$ and $`H`$ are matrix Lie groups, that is, $`G`$ is a closed subgroup of $`\mathrm{𝖦𝖫}(n;)`$ and $`H`$ is a closed subgroup of $`\mathrm{𝖦𝖫}(m;)`$, then $`G\times H`$ can be regarded in an obvious way as a closed subgroup of $`\mathrm{𝖦𝖫}(n+m;)`$. Thus the direct product of matrix Lie groups can be regarded as a matrix Lie group. It is easy to check that the Lie algebra of $`G\times H`$ is isomorphic to the direct sum of the Lie algebra of $`G`$ and the Lie algebra of $`H`$. See Exercise 13.
In light of Proposition 5.4, the representation $`\mathrm{\Pi }_1\mathrm{\Pi }_2`$ of $`G\times H`$ gives rise to a representation of the Lie algebra of $`G\times H`$, namely $`𝔤𝔥`$. The following Proposition shows that this representation of $`𝔤𝔥`$ is not what you might expect at first.
###### Proposition 5.23.
Let $`G`$ and $`H`$ be matrix Lie groups, let $`\mathrm{\Pi }_1`$, $`\mathrm{\Pi }_2`$ be representations of $`G,H`$ respectively, and consider the representation $`\mathrm{\Pi }_1\mathrm{\Pi }_2`$ of $`G\times H`$. Let $`\pi _1\pi _2`$ denote the associated representation of the Lie algebra of $`G\times H`$, namely $`𝔤𝔥`$. Then for all $`X𝔤`$ and $`Y𝔥`$
$$\pi _1\pi _2(X,Y)=\pi _1(X)I+I\pi _2(Y)\text{.}$$
###### Proof.
Suppose that $`u(t)`$ is a smooth curve in $`U`$ and $`v(t)`$ is a smooth curve in $`V`$. Then we verify the product rule in the usual way:
$`\underset{h0}{lim}{\displaystyle \frac{u(t+h)v(t+h)u(t)v(t)}{h}}`$
$`=\underset{h0}{lim}{\displaystyle \frac{u(t+h)v(t+h)u(t+h)v(t)}{h}}+{\displaystyle \frac{u(t+h)v(t)u(t)v(t)}{h}}`$
$`=\underset{h0}{lim}\left[u(t+h){\displaystyle \frac{\left(v(t+h)v\left(t\right)\right)}{h}}\right]+\underset{h0}{lim}\left[{\displaystyle \frac{\left(u(t+h)u\left(t\right)\right)}{h}}v(t)\right]\text{.}`$
Thus
$$\frac{d}{dt}\left(u(t)v(t)\right)=\frac{du}{dt}v(t)+u(t)\frac{dv}{dt}\text{.}$$
This being the case, we can compute $`\pi _1\pi _2(X,Y)`$:
$`\pi _1\pi _2(X,Y)(uv)={\displaystyle \frac{d}{dt}}|_{t=0}\mathrm{\Pi }_1\mathrm{\Pi }_2(e^{tX},e^{tY})(uv)`$
$`={\displaystyle \frac{d}{dt}}|_{t=0}\mathrm{\Pi }_1(e^{tX})u\mathrm{\Pi }_2(e^{tY})v`$
$`=\left({\displaystyle \frac{d}{dt}}|_{t=0}\mathrm{\Pi }_1(e^{tX})u\right)v+u\left({\displaystyle \frac{d}{dt}}|_{t=0}\mathrm{\Pi }_2(e^{tY})v\right)\text{.}`$
This shows that $`\pi _1\pi _2(X,Y)=\pi _1(X)I+I\pi _2(Y)`$ on elements of the form $`uv`$, and therefore on the whole space $`UV`$. ∎
###### Definition 5.24.
Let $`𝔤`$ and $`𝔥`$ be Lie algebras, and let $`\pi _1`$ and $`\pi _2`$ be representations of $`𝔤`$ and $`𝔥`$, acting on spaces $`U`$ and $`V`$. Then the tensor product of $`\pi _1`$ and $`\pi _2`$, denoted $`\pi _1\pi _2`$, is a representation of $`𝔤𝔥`$ acting on $`UV`$, given by
$$\pi _1\pi _2(X,Y)=\pi _1(X)I+I\pi _2(Y)$$
for all $`X𝔤`$ and $`Y𝔥.`$
It is easy to check that this indeed defines a representation of $`𝔤𝔥`$. Note that if we defined $`\pi _1\pi _2(X,Y)=\pi _1(X)\pi _2(Y)`$, this would not be a representation of $`𝔤𝔥`$, for this is not even a linear map. (E.g., we would then have $`\pi _1\pi _2(2X,2Y)=4\pi _1\pi _2(X,Y)`$!) Note also that the above definition applies even if $`\pi _1`$ and $`\pi _2`$ do not come from a representation of any matrix Lie group.
###### Definition 5.25.
Let $`G`$ be a matrix Lie group, and let $`\mathrm{\Pi }_1`$ and $`\mathrm{\Pi }_2`$ be representations of $`G`$, acting on spaces $`V_1`$ and $`V_2`$. Then the tensor product of $`\mathrm{\Pi }_1`$ and $`\mathrm{\Pi }_2`$ is a representation of $`G`$ acting on $`V_1V_2`$ defined by
$$\mathrm{\Pi }_1\mathrm{\Pi }_2(A)=\mathrm{\Pi }_1(A)\mathrm{\Pi }_2(A)$$
for all $`AG.`$
###### Proposition 5.26.
With the above notation, the associated representation of the Lie algebra $`𝔤`$ satisfies
$$\pi _1\pi _2(X)=\pi _1(X)I+I\pi _2(X)$$
for all $`X𝔤.`$
###### Proof.
Using the product rule,
$`\pi _1\pi _2(X)\left(uv\right)`$ $`={\displaystyle \frac{d}{dt}}|_{t=0}\mathrm{\Pi }_1\left(e^{tX}\right)u\mathrm{\Pi }_2\left(e^{tX}\right)v`$
$`=\pi _1\left(X\right)uv+v\pi _2\left(X\right)u\text{.}`$
This is what we wanted to show. ∎
###### Definition 5.27.
If $`𝔤`$ is a Lie algebra, and $`\pi _1`$ and $`\pi _2`$ are representations of $`𝔤`$ acting on spaces $`V_1`$ and $`V_2`$, then the tensor product of $`\pi _1`$ and $`\pi _2`$ is a representation of $`𝔤`$ acting on the space $`V_1V_2`$ defined by
$$\pi _1\pi _2(X)=\pi _1(X)I+I\pi _2(X)$$
for all $`X𝔤.`$
It is easy to check that $`\mathrm{\Pi }_1\mathrm{\Pi }_2`$ and $`\pi _1\pi _2`$ are actually representations of $`G`$ and $`𝔤`$, respectively. There is some ambiguity in the notation, say, $`\mathrm{\Pi }_1\mathrm{\Pi }_2`$. For even if $`\mathrm{\Pi }_1`$ and $`\mathrm{\Pi }_2`$ are both representations of the same group $`G`$, we could still regard $`\mathrm{\Pi }_1\mathrm{\Pi }_2`$ as a representation of $`G\times G`$, by taking $`H=G`$ in definition 5.22. We will rely on context to make clear whether we are thinking of $`\mathrm{\Pi }_1\mathrm{\Pi }_2`$ as a representation of $`G\times G`$ or as representation of $`G`$.
Suppose $`\mathrm{\Pi }_1`$ and $`\mathrm{\Pi }_2`$ are irreducible representations of a group $`G`$. If we regard $`\mathrm{\Pi }_1\mathrm{\Pi }_2`$ as a representation of $`G`$, it may no longer be irreducible. If it is not irreducible, one can attempt to decompose it as a direct sum of irreducible representations. This process is called Clebsch-Gordan theory. In the case of $`\mathrm{𝖲𝖴}(2)`$, this theory is relatively simple. (In the physics literature, the problem of analyzing tensor products of representations of $`\mathrm{𝖲𝖴}(2)`$ is called “addition of angular momentum.”) See Exercise 15.
### 5.7. Schur’s Lemma
Let $`\mathrm{\Pi }`$ and $`\mathrm{\Sigma }`$ be representations of a matrix Lie group $`G`$, acting on spaces $`V`$ and $`W`$. Recall that a morphism of representations is a linear map $`\varphi :VW`$ with the property that
$$\varphi \left(\mathrm{\Pi }(A)v\right)=\mathrm{\Sigma }(A)\left(\varphi (v)\right)$$
for all $`vV`$ and all $`AG`$. Schur’s Lemma is an extremely important result which tells us about morphisms of irreducible representations. Part of Schur’s Lemma applies to both real and complex representations, but part of it applies only to complex representations.
It is desirable to be able to state Schur’s lemma simultaneously for groups and Lie algebras. In order to do so, we need to indulge in a common abuse of notation. If, say, $`\mathrm{\Pi }`$ is a representation of $`G`$ acting on a space $`V`$, we will refer to $`V`$ as the representation, without explicit reference to $`\mathrm{\Pi }`$.
###### Theorem 5.28 (Schur’s Lemma).
1. Let $`V`$ and $`W`$ be irreducible real or complex representations of a group or Lie algebra, and let $`\varphi :VW`$ be a morphism. Then either $`\varphi =0`$ or $`\varphi `$ is an isomorphism.
2. Let $`V`$ be an irreducible complex representation of a group or Lie algebra, and let $`\varphi :VV`$ be a morphism of $`V`$ with itself. Then $`\varphi =\lambda I`$, for some $`\lambda `$.
3. Let $`V`$ and $`W`$ be irreducible complex representations of a group or Lie algebra, and let $`\varphi _1,\varphi _2:VW`$ be non-zero morphisms. Then $`\varphi _1=\lambda \varphi _2`$, for some $`\lambda `$.
###### Corollary 5.29.
Let $`\mathrm{\Pi }`$ be an irreducible complex representation of a matrix Lie group $`G`$. If $`A`$ is in the center of $`G`$, then $`\mathrm{\Pi }(A)=\lambda I`$. Similarly, if $`\pi `$ is an irreducible complex representation of a Lie algebra $`𝔤`$, and if $`X`$ is in the center of $`𝔤`$ (i.e., $`[X,Y]=0`$ for all $`Y𝔤`$), then $`\pi (X)=\lambda I`$.
###### Proof.
We prove the group case; the proof of the Lie algebra case is the same. If $`A`$ is in the center of $`G`$, then for all $`BG`$,
$$\mathrm{\Pi }(A)\mathrm{\Pi }(B)=\mathrm{\Pi }(AB)=\mathrm{\Pi }(BA)=\mathrm{\Pi }(B)\mathrm{\Pi }(A)\text{.}$$
But this says exactly that $`\mathrm{\Pi }(A)`$ is a morphism of $`\mathrm{\Pi }`$ with itself. So by Point 2 of Schur’s lemma, $`\mathrm{\Pi }(A)`$ is a multiple of the identity. ∎
###### Corollary 5.30.
An irreducible complex representation of a commutative group or Lie algebra is one-dimensional.
###### Proof.
Again, we prove only the group case. If $`G`$ is commutative, then the center of $`G`$ is all of $`G`$, so by the previous corollary $`\mathrm{\Pi }(A)`$ is a multiple of the identity for each $`AG`$. But this means that every subspace of $`V`$ is invariant! Thus the only way that $`V`$ can fail to have a non-trivial invariant subspace is for it not to have any non-trivial subspaces. This means that $`V`$ must be one-dimensional. (Recall that we do not allow $`V`$ to be zero-dimensional.) ∎
###### Proof.
As usual, we will prove just the group case; the proof of the Lie algebra case requires only the obvious notational changes.
Proof of 1. Saying that $`\varphi `$ is a morphism means $`\varphi (\mathrm{\Pi }(A)v)=\mathrm{\Sigma }(A)\left(\varphi (v)\right)`$ for all $`vV`$ and all $`AG`$. Now suppose that $`v\mathrm{ker}(\varphi )`$. Then
$$\varphi (\mathrm{\Pi }(A)v)=\mathrm{\Sigma }(A)\varphi (v)=0\text{.}$$
This shows that $`\mathrm{ker}\varphi `$ is an invariant subspace of $`V`$. Since $`V`$ is irreducible, we must have $`\mathrm{ker}\varphi =0`$ or $`\mathrm{ker}\varphi =V`$. Thus $`\varphi `$ is either one-to-one or zero.
Suppose $`\varphi `$ is one-to-one. Then the image of $`\varphi `$ is a non-zero subspace of $`W`$. On the other hand, the image of $`\varphi `$ is invariant, for if $`wW`$ is of the form $`\varphi (v)`$ for some $`vV`$, then
$$\mathrm{\Sigma }(A)w=\mathrm{\Sigma }(A)\varphi (v)=\varphi (\mathrm{\Pi }(A)v)\text{.}$$
Since $`W`$ is irreducible and image$`(V)`$ is non-zero and invariant, we must have image$`(V)=W`$. Thus $`\varphi `$ is either zero or one-to-one and onto.
Proof of 2. Suppose now that $`V`$ is an irreducible complex representation, and that $`\varphi :VV`$ is a morphism of $`V`$ to itself. This means that $`\varphi \mathrm{\Pi }(A)=\mathrm{\Pi }(A)\varphi `$ for all $`AG`$, i.e., that $`\varphi `$ commutes with all of the $`\mathrm{\Pi }(A)`$’s. Now, since we are over an algebraically complete field, $`\varphi `$ must have at least one eigenvalue $`\lambda `$. Let $`U`$ denote the eigenspace for $`\varphi `$ associated to the eigenvalue $`\lambda `$, and let $`uU`$. Then for each $`AG`$
$$\varphi \left(\mathrm{\Pi }(A)u\right)=\mathrm{\Pi }(A)\varphi (v)=\lambda \mathrm{\Pi }(A)u\text{.}$$
Thus applying $`\mathrm{\Pi }(A)`$ to an eigenvector of $`\varphi `$ with eigenvalue $`\lambda `$ yields another eigenvector of $`\varphi `$ with eigenvalue $`\lambda `$. That is, $`U`$ is invariant.
Since $`\lambda `$ is an eigenvalue, $`U0`$, and so we must have $`U=V`$. But this means that $`\varphi (v)=\lambda v`$ for all $`vV`$, i.e., that $`\varphi =\lambda I`$.
Proof of 3. If $`\varphi _20`$, then by (1) $`\varphi _2`$ is an isomorphism. Now look at $`\varphi _1\varphi _2^1`$. As is easily checked, the composition of two morphisms is a morphism, so $`\varphi _1\varphi _2^1`$ is a morphism of $`W`$ with itself. Thus by (2), $`\varphi _1\varphi _2^1=\lambda I`$, whence $`\varphi _1=\lambda \varphi _2`$. ∎
### 5.8. Group Versus Lie Algebra Representations
We know from Chapter 3 (Theorem 3.18) that every Lie group homomorphism gives rise to a Lie algebra homomorphism. In particular, this shows (Proposition 5.4) that every representation of a matrix Lie group gives rise to a representation of the associated Lie algebra. The goal of this section is to investigate the reverse process. That is, given a representation of the Lie algebra, under what circumstances is there an associated representation of the Lie group?
The climax of this section is Theorem 5.33, which states that if $`G`$ is a connected and simply connected matrix Lie group with Lie algebra $`𝔤`$, and if $`\pi `$ is a representation of $`𝔤`$, then there is a unique representation $`\mathrm{\Pi }`$ of $`G`$ such that $`\mathrm{\Pi }`$ and $`\pi `$ are related as in Proposition 5.4. Our proof of this theorem will make use of the Baker-Campbell-Hausdorff formula from Chapter 4. Before turning to this general theorem, we will examine two special cases, namely $`\mathrm{𝖲𝖮}(3)`$ and $`\mathrm{𝖲𝖴}(2)`$, for which we can work things out by hand. See Bröcker and tom Dieck, Chapter II, Section 5.
We have shown (Theorem 5.9) that every irreducible complex representation of $`\mathrm{𝗌𝗎}(2)`$ is equivalent to one of the representations $`\pi _m`$ described in Section 5.3. (Recall that the irreducible complex representations of $`\mathrm{𝗌𝗎}(2)`$ are in one-to-one correspondence with the irreducible representations of $`\mathrm{𝗌𝗅}(2;)`$.) Each of the representations $`\pi _m`$ of $`\mathrm{𝗌𝗎}(2)`$ was constructed from the corresponding representation $`\mathrm{\Pi }_m`$ of the group $`\mathrm{𝖲𝖴}(2)`$. Thus we see, by brute force computation, that every irreducible complex representation of $`\mathrm{𝗌𝗎}(2)`$ actually comes from a representation of the group $`\mathrm{𝖲𝖴}(2)`$! This is consistent with the fact that $`\mathrm{𝖲𝖴}(2)`$ is simply connected (Chapter 2, Prop. 2.12).
Let us now consider the situation for $`\mathrm{𝖲𝖮}(3)`$. (Which is not simply connected.) We know from Exercise 10 of Chapter 3 that the Lie algebras $`\mathrm{𝗌𝗎}(2)`$ and $`\mathrm{𝗌𝗈}(3)`$ are isomorphic. In particular, if we take the basis
$$\begin{array}{ccc}E_1=\frac{1}{2}\left(\begin{array}{cc}i& 0\\ 0& i\end{array}\right)& E_2=\frac{1}{2}\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)& E_3=\frac{1}{2}\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right)\end{array}$$
for $`\mathrm{𝗌𝗎}(2)`$ and the basis
$$\begin{array}{ccc}F_1=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 1\\ 0& 1& 0\end{array}\right)& F_2=\left(\begin{array}{ccc}0& 0& 1\\ 0& 0& 0\\ 1& 0& 0\end{array}\right)& F_3=\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right)\end{array}$$
then direct computation shows that $`[E_1,E_2]=E_3`$, $`[E_2,E_3]=E_1`$, $`[E_3,E_1]=E_2`$, and similarly with the $`E`$’s replaced by the $`F`$’s. Thus the map $`\varphi :\mathrm{𝗌𝗈}(3)\mathrm{𝗌𝗎}(2)`$ which takes $`F_i`$ to $`E_i`$ will be a Lie algebra isomorphism.
Since $`\mathrm{𝗌𝗎}(2)`$ and $`\mathrm{𝗌𝗈}(3)`$ are isomorphic Lie algebras, they must have “the same” representations. Specifically, if $`\pi `$ is a representation of $`\mathrm{𝗌𝗎}(2)`$, then $`\pi \varphi `$ will be a representation of $`\mathrm{𝗌𝗈}(3)`$, and every representation of $`\mathrm{𝗌𝗈}(3)`$ is of this form. In particular, the irreducible representations of $`\mathrm{𝗌𝗈}(3)`$ are precisely of the form $`\sigma _m=\pi _m\varphi `$. We wish to determine, for a particular $`m`$, whether there is a representation $`\mathrm{\Sigma }_m`$ of the group $`\mathrm{𝖲𝖮}(3)`$ such that $`\sigma _m`$ and $`\mathrm{\Sigma }_m`$ are related as in Proposition 5.4.
###### Proposition 5.31.
Let $`\sigma _m=\pi _m\varphi `$ be the irreducible complex representations of the Lie algebra $`\mathrm{𝗌𝗈}(3)`$ ($`m0`$). If $`m`$ is even, then there is a representation $`\mathrm{\Sigma }_m`$ of the group $`\mathrm{𝖲𝖮}(3)`$ such that $`\sigma _m`$ and $`\mathrm{\Sigma }_m`$ are related as in Proposition 5.4. If $`m`$ is odd, then there is no such representation of $`\mathrm{𝖲𝖮}(3)`$.
Note that the condition that $`m`$ be even is equivalent to the condition that $`dimV_m=m+1`$ be odd. Thus it is the odd-dimensional representations of the Lie algebra $`\mathrm{𝗌𝗈}(3)`$ which come from group representations.
In the physics literature, the representations of $`\mathrm{𝗌𝗎}(2)/\mathrm{𝗌𝗈}(3)`$ are labeled by the parameter $`l=m/2`$. In terms of this notation, a representation of $`\mathrm{𝗌𝗈}(3)`$ comes from a representation of $`\mathrm{𝖲𝖮}(3)`$ if and only if $`l`$ is an integer. The representations with $`l`$ an integer are called “integer spin”; the others are called “half-integer spin.”
##### Proof
###### Proof.
Case 1: m odd. In this case, we want to prove that there is no representation $`\mathrm{\Sigma }_m`$ such that $`\sigma _m`$ and $`\mathrm{\Sigma }_m`$ are related as in Proposition 5.4. (We have already considered the case $`m=1`$ in Exercise 7.) Suppose, to the contrary, that there is such a $`\mathrm{\Sigma }_m`$. Then Proposition 5.4 says that
$$\mathrm{\Sigma }_m(e^X)=e^{\sigma _m(X)}$$
for all $`X\mathrm{𝗌𝗈}(3)`$. In particular, take $`X=2\pi F_1`$. Then, computing as in Chapter 3, Section 3.2 we see that
$$e^{2\pi F_1}=\left(\begin{array}{ccc}1& 0& 0\\ 0& \mathrm{cos}2\pi & \mathrm{sin}2\pi \\ 0& \mathrm{sin}2\pi & \mathrm{cos}2\pi \end{array}\right)=I\text{.}$$
Thus on the one hand $`\mathrm{\Sigma }_m\left(e^{2\pi F_1}\right)=\mathrm{\Sigma }_m(I)=I`$, while on the other hand $`\mathrm{\Sigma }_m\left(e^{2\pi F_1}\right)=e^{2\pi \sigma _m(F_1)}`$.
Let us compute $`e^{2\pi \sigma _m(F_1)}`$. By definition, $`\sigma _m(F_1)=\pi _m(\varphi (F_1))=\pi _m(E_1)`$. But, $`E_1=\frac{i}{2}H`$, where as usual
$$H=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\text{.}$$
We know that there is a basis $`u_0,u_1,\mathrm{},u_m`$ for $`V_m`$ such that $`u_k`$ is an eigenvector for $`\pi _m(H)`$ with eigenvalue $`m2k`$. This means that $`u_k`$ is also an eigenvector for $`\sigma _m(F_1)=\frac{i}{2}\pi _m(H)`$, with eigenvalue $`\frac{i}{2}(m2k)`$. Thus in the basis $`\left\{u_k\right\}`$ we have
$$\sigma _m(F_1)=\left(\begin{array}{cccc}\frac{i}{2}m& & & \\ & \frac{i}{2}(m2)& & \\ & & \mathrm{}& \\ & & & \frac{i}{2}(m)\end{array}\right)\text{.}$$
But we are assuming the $`m`$ is odd! This means that $`m2k`$ is an odd integer. Thus $`e^{2\pi \frac{i}{2}(m2k)}=1`$, and in the basis $`\left\{u_k\right\}`$
$$e^{2\pi \sigma _m(F_1)}=\left(\begin{array}{cccc}e^{2\pi \frac{i}{2}m}& & & \\ & e^{2\pi \frac{i}{2}(m2)}& & \\ & & \mathrm{}& \\ & & & e^{2\pi \frac{i}{2}(m)}\end{array}\right)=I\text{.}$$
Thus on the one hand, $`\mathrm{\Sigma }_m\left(e^{2\pi F_1}\right)=\mathrm{\Sigma }_m(I)=I`$, while on the other hand $`\mathrm{\Sigma }_m\left(e^{2\pi F_1}\right)=e^{2\pi \sigma _m(F_1)}=I`$. This is a contradiction, so there can be no such group representation $`\mathrm{\Sigma }_m`$.
Case 2: m is even. We will use the following:
###### Lemma 5.32.
There exists a Lie group homomorphism $`\mathrm{\Phi }:\mathrm{𝖲𝖴}(2)\mathrm{𝖲𝖮}(3)`$ such that
1) $`\mathrm{\Phi }`$ maps $`\mathrm{𝖲𝖴}(2)`$ onto $`\mathrm{𝖲𝖮}(3),`$
2) $`\mathrm{ker}\mathrm{\Phi }=\{I,I\},`$ and
3) the associated Lie algebra homomorphism $`\stackrel{~}{\mathrm{\Phi }}:\mathrm{𝗌𝗎}(2)\mathrm{𝗌𝗈}(3)`$ is an isomorphism which takes $`E_i`$ to $`F_i`$. That is, $`\stackrel{~}{\mathrm{\Phi }}=\varphi ^1`$.
###### Proof.
Exercise 17. ∎
Now consider the representations $`\mathrm{\Pi }_m`$ of $`\mathrm{𝖲𝖴}(2)`$. I claim that if $`m`$ is even, then $`\mathrm{\Pi }_m(I)=I`$. To see this, note that
$$e^{2\pi E_1}=\mathrm{exp}\left(\begin{array}{cc}\pi i& 0\\ 0& \pi i\end{array}\right)=I\text{.}$$
Thus $`\mathrm{\Pi }_m(I)=\mathrm{\Pi }_m(e^{2\pi E_1})=e^{\pi _m(2\pi E_1)}`$. But as in Case 1,
$$e^{\pi _m(2\pi E_1)}=\left(\begin{array}{cccc}e^{2\pi \frac{i}{2}m}& & & \\ & e^{2\pi \frac{i}{2}(m2)}& & \\ & & \mathrm{}& \\ & & & e^{2\pi \frac{i}{2}(m)}\end{array}\right)\text{.}$$
Only, this time, $`m`$ is even, and so $`\frac{i}{2}(m2k)`$ is an integer, so that $`\mathrm{\Pi }_m(I)=e^{\pi _m(2\pi E_1)}=I`$.
Since $`\mathrm{\Pi }_m(I)=I`$, $`\mathrm{\Pi }_m(U)=\mathrm{\Pi }_m(U)`$ for all $`U\mathrm{𝖲𝖴}(2)`$. According to Lemma 5.32, for each $`R\mathrm{𝖲𝖮}(3)`$, there is a unique pair of elements $`\{U,U\}`$ such that $`\mathrm{\Phi }(U)=\mathrm{\Phi }(U)=R`$. Since $`\mathrm{\Pi }_m(U)=\mathrm{\Pi }_m(U)`$, it makes sense to define
$$\mathrm{\Sigma }_m(R)=\mathrm{\Pi }_m(U)\text{.}$$
It is easy to see that $`\mathrm{\Sigma }_m`$ is a Lie group homomorphism (hence, a representation). By construction, we have
(5.12)
$$\mathrm{\Pi }_m=\mathrm{\Sigma }_m\mathrm{\Phi }\text{.}$$
Now, if $`\stackrel{~}{\mathrm{\Sigma }}_m`$ denotes the Lie algebra representation associated to $`\mathrm{\Sigma }_m`$, then it follows from (5.12) that
$$\pi _m=\stackrel{~}{\mathrm{\Sigma }}_m\stackrel{~}{\mathrm{\Phi }}\text{.}$$
But the Lie algebra homomorphism $`\stackrel{~}{\mathrm{\Phi }}`$ takes $`E_i`$ to $`F_i`$, that is, $`\stackrel{~}{\mathrm{\Phi }}=\varphi ^1`$. So $`\pi _m=\stackrel{~}{\mathrm{\Sigma }}_m\varphi ^1`$, or $`\stackrel{~}{\mathrm{\Sigma }}_m=\pi _m\varphi `$. Thus $`\stackrel{~}{\mathrm{\Sigma }}_m=\sigma _m`$, which is what we want to show. ∎
It is now time to state the main theorem.
###### Theorem 5.33.
1. Let $`G,H`$ be a matrix Lie groups, let $`\varphi _1,\varphi _2:GH`$ be Lie group homomorphisms, and let $`\stackrel{~}{\varphi }_1,\stackrel{~}{\varphi }_2:𝔤𝔥`$ be the associated Lie algebra homomorphisms. If $`G`$ is connected and $`\stackrel{~}{\varphi }_1=\stackrel{~}{\varphi }_2`$, then $`\varphi _1=\varphi _2`$.
2. Let $`G,H`$ be a matrix Lie groups with Lie algebras $`𝔤`$ and $`𝔥`$. Let $`\stackrel{~}{\varphi }:𝔤𝔥`$ be a Lie algebra homomorphism. If $`G`$ is connected and simply connected, then there exists a unique Lie group homomorphism $`\varphi :GH`$ such that $`\varphi `$ and $`\stackrel{~}{\varphi }`$ are related as in Theorem 3.18 of Chapter 3.
This has the following corollaries.
###### Corollary 5.34.
Suppose $`G`$ and $`H`$ are connected, simply connected matrix Lie groups with Lie algebras $`𝔤`$ and $`𝔥`$. If $`𝔤𝔥`$ then $`GH`$.
###### Proof.
Let $`\stackrel{~}{\varphi }:𝔤𝔥`$ be a Lie algebra isomorphism. By Theorem 5.33, there exists an associated Lie group homomorphism $`\varphi :GH`$. Since $`\stackrel{~}{\varphi }^1:𝔥𝔤`$ is also a Lie algebra homomorphism, there is a corresponding Lie group homomorphism $`\psi :HG`$. We want to show that $`\varphi `$ and $`\psi `$ are inverses of each other.
Well, $`\stackrel{~}{\varphi \psi }=\stackrel{~}{\varphi }\stackrel{~}{\psi }=I_𝔥`$, so by the Point 1 of the Theorem, $`\varphi \psi =I_H`$. Similarly, $`\psi \varphi =I_G`$. ∎
###### Corollary 5.35.
1. Let $`G`$ be a connected matrix Lie group, let $`\mathrm{\Pi }_1`$ and $`\mathrm{\Pi }_2`$ be representations of $`G`$, and let $`\pi _1`$ and $`\pi _2`$ be the associated Lie algebra representations. If $`\pi _1`$ and $`\pi _2`$ are equivalent, then $`\mathrm{\Pi }_1`$ and $`\mathrm{\Pi }_2`$ are equivalent.
2. Let $`G`$ be connected and simply connected. If $`\pi `$ is a representation of $`𝔤`$, then there exists a representation $`\mathrm{\Pi }`$ of $`G`$, acting on the same space, such that $`\mathrm{\Pi }`$ and $`\pi `$ are related as in Proposition 5.4.
###### Proof.
For (1), let $`\mathrm{\Pi }_1`$ act on $`V`$ and $`\mathrm{\Pi }_2`$ on $`W`$. We assume that the associated Lie algebra representations are equivalent, i.e., that there exists an invertible linear map $`\varphi :VW`$ such that
$$\varphi \left(\pi _1(X)v\right)=\pi _2(X)\varphi (v)$$
for all $`X𝔤`$ and all $`vV`$. This is the same as saying that $`\varphi \pi _1(X)=\pi _2(X)\varphi `$, or equivalently that $`\varphi \pi _1(X)\varphi ^1=\pi _2(X)`$ (for all $`X𝔤`$).
Now define a map $`\mathrm{\Sigma }_2:G\mathrm{𝖦𝖫}(W)`$ by the formula
$$\mathrm{\Sigma }_2(A)=\varphi \mathrm{\Pi }_1(A)\varphi ^1\text{.}$$
It is trivial to check that $`\mathrm{\Sigma }_2`$ is a homomorphism. Furthermore, differentiation shows that the associated Lie algebra homomorphism is
$$\sigma _2(X)=\varphi \pi _1(X)\varphi ^1=\pi _2(X)$$
for all $`X`$. Then by (1) in the Theorem, we must also have $`\mathrm{\Sigma }_2=\mathrm{\Pi }_2`$, i.e.,
$$\varphi \mathrm{\Pi }_1(A)\varphi ^1=\mathrm{\Pi }_2(A)$$
for all $`AG`$. But this shows that $`\mathrm{\Pi }_1`$ and $`\mathrm{\Pi }_2`$ are equivalent.
Point (2) of the Corollary follows immediately from Point (2) of the Theorem, by taking $`H=\mathrm{𝖦𝖫}(V)`$. $`\mathrm{}`$
We now proceed with the proof of Theorem 5.33.
###### Proof.
Step 1: Verify Point (1) of the Theorem.
Since $`G`$ is connected, Corollary 3.26 of Chapter 3 tells us that every element $`A`$ of $`G`$ is a finite product of the form $`A=\mathrm{exp}X_1\mathrm{exp}X_2\mathrm{}\mathrm{exp}X_n`$, with $`X_i𝔤`$. But then if $`\stackrel{~}{\varphi }_1=\stackrel{~}{\varphi }_2`$, we have
$$\varphi _1\left(e^{X_1}\mathrm{}e^{X_n}\right)=e^{\stackrel{~}{\varphi }_1(X_1)}\mathrm{}e^{\stackrel{~}{\varphi }_1(X_n)}=e^{\stackrel{~}{\varphi }_2(X_1)}\mathrm{}e^{\stackrel{~}{\varphi }_2(X_n)}=\varphi _2\left(e^{X_1}\mathrm{}e^{X_n}\right)\text{.}$$
So we now need only prove Point (2).
Step 2: Define $`\varphi `$ in a neighborhood of the identity.
Proposition 3.23 of Chapter 3 says that the exponential mapping for $`G`$ has a local inverse which maps a neighborhood $`V`$ of the identity into the Lie algebra $`𝔤`$. On this neighborhood $`V`$ we can define $`\varphi :VH`$ by
$$\varphi (A)=\mathrm{exp}\left\{\stackrel{~}{\varphi }(\mathrm{log}A)\right\}\text{.}$$
That is
$$\varphi =\mathrm{exp}\stackrel{~}{\varphi }\mathrm{log}\text{.}$$
(Note that if there is to be a homomorphism $`\varphi `$ as in Theorem 3.18 of Chapter 3, then on $`V`$, $`\varphi `$ must be $`\mathrm{exp}\stackrel{~}{\varphi }\mathrm{log}`$.)
It follows from Corollary 4.4 to the Baker-Campbell-Hausdorff formula that this $`\varphi `$ is a “local homomorphism.” That is, if $`A`$ and $`B`$ are in $`V`$, and if $`AB`$ happens to be in $`V`$ as well, then $`\varphi (AB)=\varphi (A)\varphi (B)`$. (See the discussion at the beginning of Chapter 4.)
Step 3: Define $`\varphi `$ along a path.
Recall that when we say $`G`$ is connected, we really mean that $`G`$ is path-connected. Thus for any $`AG`$, there exists a path $`A(t)G`$ with $`A(0)=I`$ and $`A(1)=A`$. A compactness argument shows that there exists numbers $`0=t_0<t_1<t_2\mathrm{}<t_n=1`$ such that
(5.13)
$$A(s)A(t_i)^1V$$
for all $`s`$ between $`t_i`$ and $`t_{i+1}`$.
In particular, for $`i=0`$, we have $`A(s)V`$ for $`0st_1`$. Thus we can define $`\varphi \left(A(s)\right)`$ by Step 2 for $`s[0,t_1]`$. Now, for $`s[t_1,t_2]`$ we have by (5.13) $`A(s)A(t_1)^1V`$. Moving the $`A(t_1)`$ to the other side, this means that for $`s[t_1,t_2]`$ we can write
$$A(s)=\left[A(s)A(t_1)^1\right]A(t_1)\text{.}$$
with $`A(s)A(t_1)^1V`$. If $`\varphi `$ is to be a homomorphism, we must have
(5.14)
$$\varphi \left(A(s)\right)=\varphi \left(\left[A(s)A(t_1)^1\right]A(t_1)\right)=\varphi \left(A(s)A(t_1)^1\right)\varphi \left(A(t_1)\right)\text{.}$$
But $`\varphi \left(A(t_1)\right)`$ has already been defined, and we can define $`\varphi \left(A(s)A(t_1)^1\right)`$ by Step 2. In this way we can use (5.14) to define $`\varphi \left(A(s)\right)`$ for $`s[t_1,t_2]`$.
Proceeding on in the same way, we can define $`\varphi \left(A(s)\right)`$ successively on each interval $`[t_i,t_{i+1}]`$ until eventually we have defined $`\varphi \left(A(s)\right)`$ on the whole time interval $`[0,1]`$. This in particular serves to define $`\varphi \left(A(1)\right)=\varphi (A)`$.
Step 4: Prove independence of path.
In Step 3, we “defined” $`\varphi (A)`$ by defining $`\varphi `$ along a path joining the identity to $`A`$. For this to make sense as a definition of $`\varphi (A)`$ we have to prove that the answer is independent of the choice of path, and also, for a particular path, independent of the choice of partition $`(t_0,t_1,\mathrm{}t_n)`$.
To establish independence of partition, we first show that passing from a particular partition to a refinement of that partition doesn’t change the answer. (A refinement of a partition is one which contains all the points of the original partition, plus some other ones.) This is proved by means of the Baker-Campbell-Hausdorff formula. For example, suppose we insert an extra partition point $`s`$ between $`t_0`$ and $`t_1`$. Under the old partition we have
(5.15)
$$\varphi \left(A(t_1)\right)=\mathrm{exp}\stackrel{~}{\varphi }\mathrm{log}\left(A(t_1)\right)\text{.}$$
Under the new partition we write
$$A(t_1)=\left[A(t_1)A(s)^1\right]A(s)$$
so that
(5.16)
$$\varphi \left(A(t_1)\right)=\mathrm{exp}\stackrel{~}{\varphi }\mathrm{log}\left(A(t_1)A(s)^1\right)\mathrm{exp}\stackrel{~}{\varphi }\mathrm{log}\left(A(s)\right)\text{.}$$
But (as noted in Step 2), Corollary 4.4 of the Baker-Campbell-Hausdorff formula (Chapter 4, Section 4.2) implies that for $`A`$ and $`B`$ sufficiently near the identity
$$\mathrm{exp}\stackrel{~}{\varphi }\mathrm{log}(AB)=\left[\mathrm{exp}\stackrel{~}{\varphi }\mathrm{log}(A)\right]\left[\mathrm{exp}\stackrel{~}{\varphi }\mathrm{log}(B)\right]\text{.}$$
Thus the right sides of (5.15) and (5.16) are equal. Once we know that passing to a refinement doesn’t change the answer, we have independence of partition. For any two partitions of $`[0,1]`$ have a common refinement, namely, the union of the two.
Once we know independence of partition, we need to prove independence of path. It is at this point that we use the fact that $`G`$ is simply connected. In particular, because of simple connectedness, any two paths $`A_1(t)`$ and $`A_2(t)`$ joining the identity to $`A`$ will be homotopic with endpoints fixed. (This is a standard topological fact.) Using this, we want to prove that Step 3 gives the same answer for $`A_1`$ and $`A_2`$.
Our strategy is to deform $`A_1`$ into $`A_2`$ in a series of steps, where during each step we only change the path in a small time interval $`(t,t+ϵ)`$, keeping everything fixed on $`[0,t]`$ and on $`[t+ϵ,1]`$. Since we have independence of partition, we can take $`t`$ and $`t+ϵ`$ to be partition points. Since the time interval is small, we can assume there are no partition points between $`t`$ and $`t+ϵ`$. Then we have
$$\varphi \left(A(t+ϵ)\right)=\varphi \left(A(t+ϵ)A(t)^1\right)\varphi \left(A(t)\right)$$
where $`\varphi \left(A(t+ϵ)A(t)^1\right)`$ is defined as in Step 2.
But notice that our value for $`\varphi \left(A(t+ϵ)\right)`$ depends only on $`A\left(t\right)`$ and $`A\left(t+ϵ\right)`$, not on how we get from $`A\left(t\right)`$ to $`A\left(t+ϵ\right)`$! Thus the value $`\varphi \left(A(t+ϵ)\right)`$ doesn’t change as we deform the path. But if $`\varphi \left(A(t+ϵ)\right)`$ doesn’t change as we deform the path, neither does $`\varphi \left(A(1)\right)`$, since the path isn’t changing on $`[t+ϵ,1]`$.
Since $`A_1`$ and $`A_2`$ are homotopic with endpoints fixed, it is possible (by a standard topological argument) to deform $`A_1`$ into $`A_2`$ in a series of small steps as above.
Step 5: Prove that $`\varphi `$ is a homomorphism, and is properly related to $`\stackrel{~}{\varphi }`$.
Now that we have independence of path (and partition), we can give a simpler description of how to compute $`\varphi `$. Given any group element $`A`$, $`A`$ can be written in the form
$$A=C_nC_{n1}\mathrm{}C_1$$
with each $`C_i`$ in $`V`$. (This follows from the (path-)connectedness of $`G`$.) We can then choose a path $`A(t)`$ which starts at the identity, then goes to $`C_1`$, then to $`C_2C_1`$, and so on to $`C_nC_{n1}\mathrm{}C_1=A`$. We can choose a partition so that $`A(t_i)=C_iC_{i1}\mathrm{}C_1`$. By the way we have defined things
$$\varphi (A)=\varphi \left(A(1)A(t_{n1})^1\right)\varphi \left(A(t_{n1})A(t_{n2})^1\right)\mathrm{}\varphi \left(A(t_1)A(0)\right)\text{.}$$
But
$$A(t_i)A(t_{i1})^1=\left(C_iC_{i1}\mathrm{}C_1\right)\left(C_{i1}\mathrm{}C_1\right)^1=C_i$$
so
$$\varphi (A)=\varphi (C_n)\varphi (C_{n1})\mathrm{}\varphi (C_1)\text{.}$$
Now suppose that $`A`$ and $`B`$ are two elements of $`G`$ and we wish to compute $`\varphi (AB)`$. Well, write
$`A=C_nC_{n1}\mathrm{}C_1`$
$`B=D_nD_{n1}\mathrm{}D_1\text{.}`$
Then
$`\varphi \left(AB\right)=\varphi \left(C_nC_{n1}\mathrm{}C_1D_nD_{n1}\mathrm{}D_1\right)`$
$`=\left[\varphi (C_n)\mathrm{}\varphi (C_1)\right]\left[\varphi (D_n)\mathrm{}\varphi (D_1)\right]`$
$`=\varphi (A)\varphi (B)\text{.}`$
We see then that $`\varphi `$ is a homomorphism. It remains only to verify that $`\varphi `$ has the proper relationship to $`\stackrel{~}{\varphi }`$. But since $`\varphi `$ is defined near the identity to be $`\varphi =\mathrm{exp}\stackrel{~}{\varphi }\mathrm{log}`$, we see that
$$\frac{d}{dt}|_{t=0}\varphi \left(e^{tX}\right)=\frac{d}{dt}|_{t=0}e^{t\stackrel{~}{\varphi }(X)}=\stackrel{~}{\varphi }(X)\text{.}$$
Thus $`\stackrel{~}{\varphi }`$ is the Lie algebra homomorphism associated to the Lie group homomorphism $`\varphi `$.
This completes the proof of Theorem 5.33. ∎
### 5.9. Covering Groups
It is at this point that we pay the price for our decision to consider only matrix Lie groups. For the universal covering group of a matrix Lie group (defined below) is always a Lie group, but not always a matrix Lie group. For example, the universal covering group of $`\mathrm{𝖲𝖫}(n;)`$ ($`n2`$) is a Lie group, but not a matrix Lie group. (See Exercise 20.)
The notion of a universal cover allows us to determine, in the case of a non-simply connected group, which representations of the Lie algebra correspond to representations of the group. See Theorem 5.41 below.
###### Definition 5.36.
Let $`G`$ be a connected matrix Lie group. A universal covering group of $`G`$ (or just universal cover) is a connected, simply connected Lie group $`\stackrel{~}{G}`$, together with a Lie group homomorphism $`\varphi :\stackrel{~}{G}G`$ (called the projection map) with the following properties:
1. $`\varphi `$ maps $`\stackrel{~}{G}`$ onto $`G`$.
2. There is a neighborhood $`U`$ of $`I`$ in $`\stackrel{~}{G}`$ which maps homeomorphically under $`\varphi `$ onto a neighborhood $`V`$ of $`I`$ in $`G`$.
###### Proposition 5.37.
If $`G`$ is any connected matrix Lie group, then a universal covering group $`\stackrel{~}{G}`$ of $`G`$ exists and is unique up to canonical isomorphism.
We will not prove this theorem, but the idea of proof is as follows. We assume that $`G`$ is a matrix Lie group, hence a Lie group (that is, a manifold). As a manifold, $`G`$ has a topological universal cover $`\stackrel{~}{G}`$ which is a connected, simply connected manifold. The universal cover comes with a “projection map” $`\varphi :\stackrel{~}{G}G`$ which is a local homeomorphism. Now, since $`G`$ is not only a manifold but also a group, $`\stackrel{~}{G}`$ also becomes a group, and the projection map $`\varphi `$ becomes a homomorphism.
###### Proposition 5.38.
Let $`G`$ be a connected matrix Lie group, $`\stackrel{~}{G}`$ its universal cover, and $`\varphi `$ the projection map from $`\stackrel{~}{G}`$ to $`G`$. Suppose that $`\stackrel{~}{G}`$ is a matrix Lie group with Lie algebra $`\stackrel{~}{𝔤}`$. Then the associated Lie algebra map
$$\stackrel{~}{\varphi }:\stackrel{~}{𝔤}𝔤$$
is an isomorphism.
In light of this Proposition, we often say that $`G`$ and $`\stackrel{~}{G}`$ have the same Lie algebra.
The above Proposition is true even if $`\stackrel{~}{G}`$ is not a matrix Lie group. But to make sense out of the Proposition in that case, we need the definition of the Lie algebra of a general Lie group, which we have not defined.
###### Proof.
Exercise 18. ∎
#### 5.9.1. Examples
The universal cover of $`S^1`$ is $``$, and the projection map is the map $`xe^{ix}`$. The universal cover of $`\mathrm{𝖲𝖮}(3)`$ is $`\mathrm{𝖲𝖴}(2)`$, and the projection map is the homomorphism described in Lemma 5.32.
More generally, we can consider $`\mathrm{𝖲𝖮}(n)`$ for $`n3`$. As it turns out, for $`n3`$ the universal cover of $`\mathrm{𝖲𝖮}(n)`$ is a double cover. (That is, the projection map $`\varphi `$ is two-to-one.) The universal cover of $`\mathrm{𝖲𝖮}(n)`$ is called $`\mathrm{𝖲𝗉𝗂𝗇}(n)`$, and may be constructed as a certain group of invertible elements in the Clifford algebra over $`^n`$. See Bröcker and tom Dieck, Chapter I, Section 6, especially Propositions I.6.17 and I.6.19. In particular, $`\mathrm{𝖲𝗉𝗂𝗇}(n)`$ is a matrix Lie group.
The case $`n=4`$ is quite special. It turns out that the universal cover of $`\mathrm{𝖲𝖮}(4)`$ (i.e., $`\mathrm{𝖲𝗉𝗂𝗇}(4)`$) is isomorphic to $`\mathrm{𝖲𝖴}(2)\times \mathrm{𝖲𝖴}(2)`$. This is best seen by regarding $`^4`$ as the quaternion algebra.
###### Theorem 5.39.
Let $`G`$ be a matrix Lie group, and suppose that $`\stackrel{~}{G}`$ is also a matrix Lie group. Identify the Lie algebra of $`\stackrel{~}{G}`$ with the Lie algebra $`𝔤`$ of $`G`$ as in Proposition 5.38. Suppose that $`H`$ is a matrix Lie group with Lie algebra $`𝔥`$, and that $`\stackrel{~}{\varphi }:𝔤𝔥`$ is a homomorphism. Then there exists a unique Lie group homomorphism $`\varphi :\stackrel{~}{G}H`$ such that $`\varphi `$ and $`\stackrel{~}{\varphi }`$ are related as in Theorem 3.18 of Chapter 3.
###### Proof.
$`\stackrel{~}{G}`$ is simply connected. ∎
###### Corollary 5.40.
Let $`G`$ and $`\stackrel{~}{G}`$ be as in Theorem 5.39, and let $`\pi `$ be a representation of $`𝔤`$. Then there exists a unique representation $`\stackrel{~}{\mathrm{\Pi }}`$ of $`\stackrel{~}{G}`$ such that
$$\pi (X)=\frac{d}{dt}|_{t=0}\stackrel{~}{\mathrm{\Pi }}\left(e^{tX}\right)$$
for all $`X𝔤`$.
###### Theorem 5.41.
Let $`G`$ and $`\stackrel{~}{G}`$ be as in Theorem 5.39, and let $`\varphi :\stackrel{~}{G}G`$. Now let $`\pi `$ be a representation of $`𝔤`$, and $`\stackrel{~}{\mathrm{\Pi }}`$ the associated representation of $`\stackrel{~}{G}`$, as in the Corollary. Then there exists a representation $`\mathrm{\Pi }`$ of $`G`$ corresponding to $`\pi `$ if and only if
$$\mathrm{ker}\stackrel{~}{\mathrm{\Pi }}\mathrm{ker}\varphi \text{.}$$
###### Proof.
Exercise 19. ∎
### 5.10. Exercises
1. Let $`G`$ be a matrix Lie group, and $`𝔤`$ its Lie algebra. Let $`\mathrm{\Pi }_1`$ and $`\mathrm{\Pi }_2`$ be representations of $`G`$, and let $`\pi _1`$ and $`\pi _2`$ be the associated representations of $`𝔤`$ (Proposition 5.4). Show that if $`\mathrm{\Pi }_1`$ and $`\mathrm{\Pi }_2`$ are equivalent representations of $`G`$, then $`\pi _1`$ and $`\pi _2`$ are equivalent representations of $`𝔤`$. Show that if $`G`$ is connected, and if $`\pi _1`$ and $`\pi _2`$ are equivalent representations of $`𝔤`$, then $`\mathrm{\Pi }_1`$ and $`\mathrm{\Pi }_2`$ are equivalent representations of $`G`$.
Hint: Use Corollary 3.26 of Chapter 3.
2. Let $`G`$ be a connected matrix Lie group with Lie algebra $`𝔤`$. Let $`\mathrm{\Pi }`$ be a representation of $`G`$ acting on a space $`V`$, and let $`\pi `$ be the associated Lie algebra representation. Show that a subspace $`WV`$ is invariant for $`\mathrm{\Pi }`$ if and only if it is invariant for $`\pi `$. Show that $`\mathrm{\Pi }`$ is irreducible if and only if $`\pi `$ is irreducible.
3. Suppose that $`\mathrm{\Pi }`$ is a finite-dimensional unitary representation of a matrix Lie group $`G`$. (That is, $`V`$ is a finite-dimensional Hilbert space, and $`\mathrm{\Pi }`$ is a continuous homomorphism of $`G`$ into $`U(V)`$.) Let $`\pi `$ be the associated representation of the Lie algebra $`𝔤`$. Show that for each $`X𝔤`$, $`\pi (X)^{}=\pi (X)`$.
4. Show explicitly that the adjoint representation and the standard representation are equivalent representations of the Lie algebra $`\mathrm{𝗌𝗈}(3)`$. Show that the adjoint and standard representations of the group $`\mathrm{𝖲𝖮}(3)`$ are equivalent.
5. Consider the elements $`E_1`$, $`E_2`$, and $`E_3`$ in $`\mathrm{𝗌𝗎}(2)`$ defined in Exercise 9 of Chapter 3. These elements form a basis for the real vector space $`\mathrm{𝗌𝗎}(2)`$. Show directly that $`E_1`$, $`E_2`$, and $`E_3`$ form a basis for the complex vector space $`\mathrm{𝗌𝗅}(2;)`$.
6. Define a vector space with basis $`u_0,u_1\mathrm{}u_m`$. Now define operators $`\pi (H)`$, $`\pi (X)`$, and $`\pi (Y)`$ by formula (5.10). Verify by direct computation that the operators defined by (5.10) satisfy the commutation relations $`[\pi (H),\pi (X)]=2\pi (X)`$, $`[\pi (H),\pi (Y)]=2\pi (Y)`$, and $`[\pi (X),\pi (Y)]=\pi (H)`$. (Thus $`\pi (H)`$, $`\pi (X)`$, and $`\pi (Y)`$ define a representation of $`\mathrm{𝗌𝗅}(2;)`$.) Show that this representation is irreducible.
Hint: It suffices to show, for example, that $`[\pi (H),\pi (X)]=2\pi (X)`$ on each basis element. When dealing with $`\pi (Y)`$, don’t forget to treat separately the case of $`u_k`$, $`k<m`$, and the case of $`u_m`$.
7. We can define a two-dimensional representation of $`\mathrm{𝗌𝗈}(3)`$ as follows:
$$\pi \left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 1\\ 0& 1& 0\end{array}\right)=\frac{1}{2}\left(\begin{array}{cc}i& 0\\ 0& i\end{array}\right);$$
$$\pi \left(\begin{array}{ccc}0& 0& 1\\ 0& 0& 0\\ 1& 0& 0\end{array}\right)=\frac{1}{2}\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right);$$
$$\pi \left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right)=\frac{1}{2}\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right)\text{.}$$
(You may assume that this actually gives a representation.) Show that there is no group representation $`\mathrm{\Pi }`$ of $`\mathrm{𝖲𝖮}(3)`$ such that $`\mathrm{\Pi }`$ and $`\pi `$ are related as in Proposition 5.4.
Hint: If $`X\mathrm{𝗌𝗈}(3)`$ is such that $`e^X=I`$, and $`\mathrm{\Pi }`$ is any representation of $`\mathrm{𝖲𝖮}(3)`$, then $`\mathrm{\Pi }(e^X)=\mathrm{\Pi }(I)=I`$.
Remark: In the physics literature, this non-representation of $`\mathrm{𝖲𝖮}(3)`$ is called “spin $`\frac{1}{2}`$.”
8. Consider the standard representation of the Heisenberg group, acting on $`^3`$. Determine all subspaces of $`^3`$ which are invariant under the action of the Heisenberg group. Is this representation completely reducible?
9. Give an example of a representation of the commutative group $``$ which is not completely reducible.
10. Consider the unitary representations $`\mathrm{\Pi }_{\mathrm{}}`$ of the real Heisenberg group. Assume that there is some sort of associated representation $`\pi _{\mathrm{}}`$ of the Lie algebra, which should be given by
$$\pi _{\mathrm{}}(X)f=\frac{d}{dt}|_{t=0}\mathrm{\Pi }_{\mathrm{}}\left(e^{tX}\right)f$$
(We have not proved any theorem of this sort for infinite-dimensional unitary representations.)
Computing in a purely formal manner (that is, ignoring all technical issues) compute
$$\begin{array}{ccc}\pi _{\mathrm{}}\left(\begin{array}{ccc}0& 1& 0\\ 0& 0& 0\\ 0& 0& 0\end{array}\right);& \pi _{\mathrm{}}\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 1\\ 0& 0& 0\end{array}\right);& \pi _{\mathrm{}}\left(\begin{array}{ccc}0& 0& 1\\ 0& 0& 0\\ 0& 0& 0\end{array}\right)\end{array}\text{.}$$
Verify (still formally) that these operators have the right commutation relations to generate a representation of the Lie algebra of the real Heisenberg group. (That is, verify that on this basis, $`\pi _{\mathrm{}}[X,Y]=[\pi _{\mathrm{}}(X),\pi _{\mathrm{}}(Y)]`$.)
Why is this computation not rigorous?
11. Consider the Heisenberg group over the field $`_p`$ of integers mod $`p`$, with $`p`$ prime, namely
$$H_p=\left\{\left(\begin{array}{ccc}1\hfill & a\hfill & b\hfill \\ 0\hfill & 1\hfill & c\hfill \\ 0\hfill & 0\hfill & 1\hfill \end{array}\right)\right|a,b,c_p\}\text{.}$$
This is a subgroup of the group $`\mathrm{𝖦𝖫}(3;_p)`$, and has $`p^3`$ elements.
Let $`V_p`$ denote the space of complex-valued functions on $`_p`$, which is a $`p`$-dimensional complex vector space. For each non-zero $`n_p`$, define a representation of $`H_p`$ by the formula
$$\left(\mathrm{\Pi }_nf\right)\left(x\right)=e^{i2\pi nb/p}e^{i2\pi ncx/p}f\left(xa\right)\text{}x_p\text{.}$$
(These representations are analogous to the unitary representations of the real Heisenberg group, with the quantity $`2\pi n/p`$ playing the role of $`\mathrm{}`$.)
a) Show that for each $`n`$, $`\mathrm{\Pi }_n`$ is actually a representation of $`H_p`$, and that it is irreducible.
b) Determine (up to equivalence) all the one-dimensional representations of $`H_p`$.
c) Show that every irreducible representation of $`H_p`$ is either one-dimensional or equivalent to one of the $`\mathrm{\Pi }_n`$’s.
12. Prove Theorem 5.19.
Hints: For existence, choose bases $`\left\{e_i\right\}`$ and $`\left\{f_j\right\}`$ for $`U`$ and $`V`$. Then define a space $`W`$ which has as a basis $`\left\{w_{ij}\right|0in,0jm\}`$. Define $`\varphi (e_i,f_j)=w_{ij}`$ and extend by bilinearity. For uniqueness, use the universal property.
13. Let $`𝔤`$ and $`𝔥`$ be Lie algebras, and consider the vector space $`𝔤𝔥`$. Show that the following operation makes $`𝔤𝔥`$ into a Lie algebra
$$[(X_1,Y_1),(X_2,Y_2)]=([X_1,X_2],[Y_1,Y_2])\text{.}$$
Now let $`G`$ and $`H`$ be matrix Lie groups, with Lie algebras $`𝔤`$ and $`𝔥`$. Show that $`G\times H`$ can be regarded as a matrix Lie group in an obvious way, and that the Lie algebra of $`G\times H`$ is isomorphic to $`𝔤𝔥`$.
14. Suppose that $`\pi `$ is a representation of a Lie algebra $`𝔤`$ acting on a finite-dimensional vector space $`V`$. Let $`V^{}`$ denote as usual the dual space of $`V`$, that is, the space of linear functionals on $`V`$. If $`A`$ is a linear operator on $`V`$, let $`A^{tr}`$ denote the dual or transpose operator on $`V^\text{ }`$,
$$\left(A^{tr}\varphi \right)\left(v\right)=\varphi \left(Av\right)$$
for $`\varphi V^{}`$, $`vV`$. Define a representation $`\pi ^{}`$ of $`𝔤`$ on $`V^{}`$ by the formula
$$\pi ^{}\left(X\right)=\pi \left(X^{tr}\right)\text{.}$$
a) Show that $`\pi ^{}`$ is really a representation of $`𝔤`$.
b) Show that $`\left(\pi ^{}\right)^{}`$ is isomorphic to $`\pi `$.
c) Show that $`\pi ^{}`$ is irreducible if and only if $`\pi `$ is.
d) What is the analogous construction of the dual representation for representations of groups?
15. Recall the spaces $`V_m`$ introduced in Section 5.3, viewed as representations of the Lie algebra $`\mathrm{𝗌𝗅}(2;)`$. In particular, consider the space $`V_1`$ (which has dimension 2).
a) Regard $`V_1V_1`$ as a representation of $`\mathrm{𝗌𝗅}(2;)`$, as in Definition 5.27. Show that this representation is not irreducible.
b) Now view $`V_1V_1`$ as a representation of $`\mathrm{𝗌𝗅}(2;)\mathrm{𝗌𝗅}(2;)`$, as in Definition 5.24. Show that this representation is irreducible.
c) More generally, show that $`V_mV_n`$ is irreducible as a representation of $`\mathrm{𝗌𝗅}(2;)\mathrm{𝗌𝗅}(2;)`$, but reducible (except if one of $`n`$ or $`m`$ is zero) as a representation of $`\mathrm{𝗌𝗅}(2;)`$.
16. Show explicitly that $`\mathrm{exp}:\mathrm{𝗌𝗈}(3)\mathrm{𝖲𝖮}(3)`$ is onto.
Hint: Using the fact that $`\mathrm{𝖲𝖮}(3)\mathrm{𝖲𝖴}(3)`$, show that the eigenvalues of $`R\mathrm{𝖲𝖮}(3)`$ must be of one of the three following forms: $`(1,1,1)`$, $`(1,1,1)`$, or $`(1,e^{i\theta },e^{i\theta })`$. In particular, $`R`$ must have an eigenvalue equal to one. Now show that in a suitable orthonormal basis, $`R`$ is of the form
$$R=\left(\begin{array}{ccc}1& 0& 0\\ 0& \mathrm{cos}\theta & \mathrm{sin}\theta \\ 0& \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)\text{.}$$
17. Proof of Lemma 5.32.
Let $`\{E_1,E_2,E_3\}`$ be the usual basis for $`\mathrm{𝗌𝗎}(2)`$, and $`\{F_1,F_2,F_3\}`$ be the basis for $`\mathrm{𝗌𝗈}(3)`$ introduced in Section 5.8. Identify $`\mathrm{𝗌𝗎}(2)`$ with $`^3`$ by identifying the basis $`\{E_1,E_2,E_3\}`$ with the standard basis for $`^3`$. Consider $`\mathrm{ad}E_1`$, $`\mathrm{ad}E_2`$, and $`\mathrm{ad}E_3`$ as operators on $`\mathrm{𝗌𝗎}(2)`$, hence on $`^3`$. Show that $`\mathrm{ad}E_i=F_i`$, for $`i=1,2,3`$. In particular, ad is a Lie algebra isomorphism of $`\mathrm{𝗌𝗎}(2)`$ onto $`\mathrm{𝗌𝗈}(3)`$.
Now consider $`\mathrm{Ad}:\mathrm{𝖲𝖴}(2)\mathrm{𝖦𝖫}\left(\mathrm{𝖲𝖴}(2)\right)=\mathrm{𝖦𝖫}(3;)`$. Show that the image of Ad is precisely $`\mathrm{𝖲𝖮}(3)`$. Show that the kernel of Ad is $`\{I,I\}`$.
Show that $`\mathrm{Ad}:\mathrm{𝖲𝖴}(2)\mathrm{𝖲𝖮}(3)`$ is the homomorphism $`\mathrm{\Phi }`$ required by Lemma 5.32.
18. Proof of Proposition 5.38.
Suppose that $`G`$ and $`\stackrel{~}{G}`$ are matrix Lie groups. Suppose that $`\varphi :\stackrel{~}{G}G`$ is a Lie group homomorphism such that $`\varphi `$ maps some neighborhood $`U`$ of $`I`$ in $`\stackrel{~}{G}`$ homeomorphically onto a neighborhood $`V`$ of $`I`$ in $`G`$. Prove that the associated Lie algebra map $`\stackrel{~}{\varphi }:\stackrel{~}{𝔤}𝔤`$ is an isomorphism.
Hints: Suppose that $`\stackrel{~}{\varphi }`$ were not one-to-one. Show, then, that there exists a sequence of points $`A_n`$ in $`\stackrel{~}{G}`$ with $`A_nI`$, $`A_nI`$ and $`\varphi (A_n)=I`$, giving a contradiction.
To show that $`\stackrel{~}{\varphi }`$ is onto, use Step 1 of the proof of Theorem 5.33 to show that on a sufficiently small neighborhood of zero in $`\stackrel{~}{𝔤}`$,
$$\stackrel{~}{\varphi }=\mathrm{log}\varphi \mathrm{exp}\text{.}$$
Use this to show that the image of $`\stackrel{~}{\varphi }`$ contains a neighborhood of zero in $`𝔤`$. Now use linearity to show that the image of $`\stackrel{~}{\varphi }`$ is all of $`𝔤`$.
19. Proof of Theorem 5.41.
First suppose that $`\mathrm{ker}\stackrel{~}{\mathrm{\Pi }}\mathrm{ker}\varphi `$. Then construct $`\mathrm{\Pi }`$ as in the proof of Proposition 5.31.
Now suppose that there is a representation $`\mathrm{\Pi }`$ of $`G`$ for which the associated Lie algebra representation is $`\pi `$. We want to show, then, that $`\mathrm{ker}\stackrel{~}{\mathrm{\Pi }}\mathrm{ker}\varphi `$. Well, define a new representation $`\mathrm{\Sigma }`$ of $`\stackrel{~}{G}`$ by
$$\mathrm{\Sigma }=\mathrm{\Pi }\varphi \text{.}$$
Show that the associated Lie algebra homomorphism $`\sigma `$ is equal to $`\pi `$, so that, by Point (1) of Theorem 5.33, $`\stackrel{~}{\mathrm{\Pi }}=\mathrm{\Sigma }`$. What can you say about the kernel of $`\mathrm{\Sigma }`$?
20. Fix an integer $`n2`$.
a) Show that every (finite-dimensional complex) representation of the Lie algebra $`\mathrm{𝗌𝗅}(n;)`$ gives rise to a representation of the group $`\mathrm{𝖲𝖫}(n;)`$, even though $`\mathrm{𝖲𝖫}(n;)`$ is not simply connected. (You may use the fact that $`\mathrm{𝖲𝖫}(n;)`$ is simply connected.)
b) Show that the universal cover of $`\mathrm{𝖲𝖫}(n;)`$ is not isomorphic to any matrix Lie group. (You may use the fact that $`\mathrm{𝖲𝖫}(n;)`$ is not simply connected.)
21. Let $`G`$ be a matrix Lie group with Lie algebra $`𝔤`$, let $`𝔥`$ be a subalgebra of $`𝔤`$, and let $`H`$ be the unique connected Lie subgroup of $`G`$ with Lie algebra $`𝔥`$. Suppose that there exists a compact simply connected matrix Lie group $`K`$ such that the Lie algebra of $`K`$ is isomorphic to $`𝔥`$. Show that $`H`$ is closed. Is $`H`$ necessarily isomorphic to $`K`$?
## Chapter 6 The Representations of $`\mathrm{𝖲𝖴}(3)`$, and Beyond
### 6.1. Preliminaries
There is a theory of the representations of semisimple groups/Lie algebras which includes as a special case the representation theory of $`\mathrm{𝖲𝖴}(3)`$. However, I feel that it is worthwhile to examine the case of $`\mathrm{𝖲𝖴}(3)`$ separately. I feel this way partly because $`\mathrm{𝖲𝖴}(3)`$ is an important group in physics, but chiefly because the general semisimple theory is difficult to digest. Considering a non-trivial example makes it much clearer what is going on. In fact, all of the elements of the general theory are present already in the case of $`\mathrm{𝖲𝖴}(3)`$, so we do not lose too much by considering at first just this case.
The main result of this chapter is Theorem 1, which states that an irreducible finite-dimensional representation of $`\mathrm{𝖲𝖴}(3)`$ can be classified in terms of its “highest weight.” This is analogous to labeling the irreducible representations $`V_m`$ of $`\mathrm{𝖲𝖴}(2)/\mathrm{𝗌𝗅}(2;)`$ by the highest eigenvalue of $`\pi _m(H)`$. (The highest eigenvalue of $`\pi _m(H)`$ in $`V_m`$ is precisely $`m`$.) We will then discuss, without proofs, what the corresponding results are for general semisimple Lie algebras.
The group $`\mathrm{𝖲𝖴}(3)`$ is connected and simply connected (Bröcker and tom Dieck), so by Corollary 1 of Chapter 5, the finite-dimensional representations of $`\mathrm{𝖲𝖴}(3)`$ are in one-to-one correspondence with the finite-dimensional representations of the Lie algebra $`\mathrm{𝗌𝗎}(3)`$. Meanwhile, the complex representations of $`\mathrm{𝗌𝗎}(3)`$ are in one-to-one correspondence with the complex-linear representations of the complexified Lie algebra $`\mathrm{𝗌𝗎}(3)_{}`$. But $`\mathrm{𝗌𝗎}(3)_{}\mathrm{𝗌𝗅}(3;)`$, as is easily verified. Moreover, since $`\mathrm{𝖲𝖴}(3)`$ is connected, it follows that a subspace $`WV`$ is invariant under the action of $`\mathrm{𝖲𝖴}(3)`$ if and only if it is invariant under the action of $`\mathrm{𝗌𝗅}(3;)`$. Thus we have the following:
###### Proposition 6.1.
There is a one-to-one correspondence between the finite-dimensional complex representations $`\mathrm{\Pi }`$ of $`\mathrm{𝖲𝖴}(3)`$ and the finite-dimensional complex-linear representations $`\pi `$ of $`\mathrm{𝗌𝗅}(3;)`$. This correspondence is determined by the property that
$$\mathrm{\Pi }\left(e^X\right)=e^{\pi (X)}$$
for all $`X\mathrm{𝗌𝗎}(3)\mathrm{𝗌𝗅}(3;)`$.
The representation $`\mathrm{\Pi }`$ is irreducible if and only the representation $`\pi `$ is irreducible. Moreover, a subspace $`WV`$ is invariant for $`\mathrm{\Pi }`$ if and only if it is invariant for $`\pi `$.
Since $`\mathrm{𝖲𝖴}(3)`$ is compact, Proposition 5.17 of Chapter 5 tells us that all the finite-dimensional representations of $`\mathrm{𝖲𝖴}(3)`$ are completely reducible. The above proposition then implies that all the finite-dimensional representations of $`\mathrm{𝗌𝗅}(3;)`$ are completely reducible.
Moreover, we can apply the same reasoning to the group $`\mathrm{𝖲𝖴}(2)`$, its Lie algebra $`\mathrm{𝗌𝗎}(2)`$, and its complexified Lie algebra $`\mathrm{𝗌𝗅}(2;)`$. Since $`\mathrm{𝖲𝖴}(2)`$ is simply connected, there is a one-to-one correspondence between the complex representations of $`\mathrm{𝖲𝖴}(2)`$ and the representations of the complexified Lie algebra $`\mathrm{𝗌𝗅}(2;)`$. Since $`\mathrm{𝖲𝖴}(2)`$ is compact, all of the representations of $`\mathrm{𝖲𝖴}(2)`$–and therefore also of $`\mathrm{𝗌𝗅}(2;)`$–are completely reducible. Thus we have established the following.
###### Proposition 6.2.
Every finite-dimensional (complex-linear) representation of $`\mathrm{𝗌𝗅}(2;)`$ or $`\mathrm{𝗌𝗅}(3;)`$ is completely reducible. In particular, every finite-dimensional representation of $`\mathrm{𝗌𝗅}(2;)`$ or $`\mathrm{𝗌𝗅}(3;)`$ decomposes as a direct sum of irreducible invariant subspaces.
We will use the following basis for $`\mathrm{𝗌𝗅}(3;)`$:
$$\begin{array}{ccc}H_1=\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 0\end{array}\right)& H_2=\left(\begin{array}{ccc}0& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right)& \\ & & \\ X_1=\left(\begin{array}{ccc}0& 1& 0\\ 0& 0& 0\\ 0& 0& 0\end{array}\right)& X_2=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 1\\ 0& 0& 0\end{array}\right)& X_3=\left(\begin{array}{ccc}0& 0& 1\\ 0& 0& 0\\ 0& 0& 0\end{array}\right)\\ & & \\ Y_1=\left(\begin{array}{ccc}0& 0& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right)& Y_2=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 0\\ 0& 1& 0\end{array}\right)& Y_3=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 0\\ 1& 0& 0\end{array}\right)\text{.}\end{array}$$
Note that the span of $`\{H_1,X_1,Y_1\}`$ is a subalgebra of $`\mathrm{𝗌𝗅}(3;)`$ which is isomorphic to $`\mathrm{𝗌𝗅}(2;)`$, by ignoring the third row and the third column. Similarly, the span of $`\{H_2,X_2,Y_2\}`$ is a subalgebra isomorphic to $`\mathrm{𝗌𝗅}(2;)`$, by ignoring the first row and first column. Thus we have the following commutation relations
$$\begin{array}{ccccccc}\hfill [H_1,X_1]& \hfill =& \hfill 2X_1& & \hfill [H_2,X_2]& \hfill =& \hfill 2X_2\\ \hfill [H_1,Y_1]& \hfill =& \hfill 2Y_1& & \hfill [H_2,Y_2]& \hfill =& \hfill 2Y_2\\ \hfill [X_1,Y_1]& \hfill =& \hfill H_1& & \hfill [X_2,Y_2]& \hfill =& \hfill H_2\text{.}\end{array}$$
We now list all of the commutation relations among the basis elements which involve at least one of $`H_1`$ and $`H_2`$. (This includes some repetitions of the commutation relations above.)
(6.1)
$$\begin{array}{ccccccc}\hfill [H_1,H_2]& \hfill =& \hfill 0& & & & \\ & & & & & & \\ \hfill [H_1,X_1]& \hfill =& \hfill 2X_1& & \hfill [H_1,Y_1]& \hfill =& \hfill 2Y_1\\ \hfill [H_2,X_1]& \hfill =& \hfill X_1& & \hfill [H_2,Y_1]& \hfill =& \hfill Y_1\\ & & & & & & \\ \hfill [H_1,X_2]& \hfill =& \hfill X_2& & \hfill [H_1,Y_2]& \hfill =& \hfill Y_2\\ \hfill [H_2,X_2]& \hfill =& \hfill 2X_2& & \hfill [H_2,Y_2]& \hfill =& \hfill 2Y_2\\ & & & & & & \\ \hfill [H_1,X_3]& \hfill =& \hfill X_3& & \hfill [H_1,Y_3]& \hfill =& \hfill Y_3\\ \hfill [H_2,X_3]& \hfill =& \hfill X_3& & \hfill [H_2,Y_3]& \hfill =& \hfill Y_3\end{array}$$
We now list all of the remaining commutation relations.
$$\begin{array}{ccccccc}[X_1,Y_1]& =& H_1& & & & \\ [X_2,Y_2]& =& H_2& & & & \\ [X_3,Y_3]& =& H_1+H_2& & & & \end{array}$$
$$\begin{array}{ccccccc}[X_1,X_2]& =& X_3& & [Y_1,Y_2]& =& Y_3\\ [X_1,Y_2]& =& 0& & [X_2,Y_1]& =& 0\\ & & & & & & \\ [X_1,X_3]& =& 0& & [Y_1,Y_3]& =& 0\\ [X_2,X_3]& =& 0& & [Y_2,Y_3]& =& 0\\ & & & & & & \\ [X_2,Y_3]& =& Y_1& & [X_3,Y_2]& =& X_1\\ [X_1,Y_3]& =& Y_2& & [X_3,Y_1]& =& X_2\end{array}$$
Note that there is a kind of symmetry between the $`X_i`$’s and the $`Y_i`$’s. If a relation in the first column involves an $`X_i`$ and/or a $`Y_j`$, the corresponding relation in the second column will involve a $`Y_i`$ and/or an $`X_j`$. (E.g., we have the relation $`[H_1,X_2]=X_2`$ in the first column, and the relation $`[H_2,Y_2]=Y_2`$ in the second column.) See Exercise 1.
All of the analysis we will do for the representations of $`\mathrm{𝗌𝗅}(3;)`$ will be in terms of the above basis. From now on, all representations of $`\mathrm{𝗌𝗅}(3;)`$ will be assumed to be finite-dimensional and complex-linear.
### 6.2. Weights and Roots
Our basic strategy in classifying the representations of $`\mathrm{𝗌𝗅}(3;)`$ is to simultaneously diagonalize $`\pi (H_1)`$ and $`\pi (H_2)`$. Since $`H_1`$ and $`H_2`$ commute, $`\pi (H_1)`$ and $`\pi (H_2)`$ will also commute, and so there is at least a chance that $`\pi (H_1)`$ and $`\pi (H_2)`$ can be simultaneously diagonalized.
###### Definition 6.3.
If $`(\pi ,V)`$ is a representation of $`\mathrm{𝗌𝗅}(3;)`$, then an ordered pair $`\mu =(\mu _1,\mu _2)^2`$ is called a weight for $`\pi `$ if there exists $`v0`$ in $`V`$ such that
$`\pi (H_1)v=\mu _1v`$
(6.2) $`\pi (H_2)v=\mu _2v\text{.}`$
The vector $`v`$ is called a weight vector corresponding to the weight $`\mu `$. If $`\mu =(\mu _1,\mu _2)`$ is a weight, then the space of all vectors $`v`$ satisfying (6.2) is the weight space corresponding to the weight $`\mu `$.
Thus a weight is simply a pair of simultaneous eigenvalues for $`\pi (H_1)`$ and $`\pi (H_2)`$.
###### Proposition 6.4.
Every representation of $`\mathrm{𝗌𝗅}(3;)`$ has at least one weight.
###### Proof.
Since we are working over the complex numbers, $`\pi (H_1)`$ has at least one eigenvalue $`\mu _1`$. Let $`WV`$ be the eigenspace for $`\pi (H_1)`$ with eigenvalue $`\mu _1`$. I assert that $`W`$ is invariant under $`\pi (H_2)`$. To see this consider $`wW`$, and compute
$`\pi (H_1)\left(\pi (H_2)w\right)`$ $`=\pi (H_2)\pi (H_1)w`$
$`=\pi (H_2)\left(\mu _1w\right)=\mu _1\pi (H_2)w\text{.}`$
This shows that $`\pi (H_2)w`$ is either zero or an eigenvector for $`\pi (H_1)`$ with eigenvalue $`\mu _1`$; thus $`W`$ is invariant.
Thus $`\pi (H_2)`$ can be viewed as an operator on $`W`$. Again, since we are over $``$, the restriction of $`\pi (H_2)`$ to $`W`$ must have at least one eigenvector $`w`$ with eigenvalue $`\mu _2`$. But then $`w`$ is a simultaneous eigenvector for $`\pi (H_1)`$ and $`\pi (H_2)`$ with eigenvalues $`\mu _1`$ and $`\mu _2`$. ∎
Now, every representation $`\pi `$ of $`\mathrm{𝗌𝗅}(3;)`$ can be viewed, by restriction, as a representation of the subalgebra $`\{H_1,X_1,Y_1\}\mathrm{𝗌𝗅}(2;)`$. Note that, even if $`\pi `$ is irreducible as a representation of $`\mathrm{𝗌𝗅}(3;)`$, there is no reason to expect that it will still be irreducible as a representation of the subalgebra $`\{H_1,X_1,Y_1\}`$. Nevertheless, $`\pi `$ restricted to $`\{H_1,X_1,Y_1\}`$ must be some finite-dimensional representation of $`\mathrm{𝗌𝗅}(2;)`$. The same reasoning applies to the restriction of $`\pi `$ to the subalgebra $`\{H_2,X_2,Y_2\}`$, which is also isomorphic to $`\mathrm{𝗌𝗅}(2;)`$.
###### Proposition 6.5.
Let $`(\pi ,V)`$ be any finite-dimensional complex-linear representation of $`\mathrm{𝗌𝗅}(2;)=\{H,X,Y\}`$. Then all the eigenvalues of $`\pi (H)`$ are integers.
###### Proof.
By Proposition 6.2, $`V`$ decomposes as a direct sum of irreducible invariant subspaces $`V_i`$. Each $`V_i`$ must be one of the irreducible representations of $`\mathrm{𝗌𝗅}(2;)`$, which we have classified. In particular, in each $`V_i`$, $`\pi (H)`$ can be diagonalized, and the eigenvalues of $`\pi (H)`$ are integers. Thus $`\pi (H)`$ can be diagonalized on the whole space $`V`$, and all of the eigenvalues are integers. ∎
###### Corollary 6.6.
If $`\pi `$ is a representation of $`\mathrm{𝗌𝗅}(3;)`$, then all of the weights of $`\pi `$ are of the form
$$\mu =(m_1,m_2)$$
with $`m_1`$ and $`m_2`$ integers.
###### Proof.
Apply Proposition 6.5 to the restriction of $`\pi `$ to $`\{H_1,X_1,Y_1\}`$, and to the restriction of $`\pi `$ to $`\{H_2,X_2,Y_2\}`$. ∎
Our strategy now is to begin with one simultaneous eigenvector for $`\pi (H_1)`$ and $`\pi (H_2)`$, and then to apply $`\pi (X_i)`$ or $`\pi (Y_i)`$, and see what the effect is. The following definition is relevant in this context. (See Lemma 6.8 below.)
###### Definition 6.7.
An ordered pair $`\alpha =(\alpha _1,\alpha _2)^2`$ is called a root if
1. $`\alpha _1`$ and $`\alpha _2`$ are not both zero, and
2. there exists $`Z\mathrm{𝗌𝗅}(3;)`$ such that
$`[H_1,Z]=\alpha _1Z`$
$`[H_2,Z]=\alpha _2Z\text{.}`$
The element $`Z`$ is called a root vector corresponding to the root $`\alpha `$.
That is, a root is a non-zero weight for the adjoint representation. The commutation relations (6.1) tell us what the roots for $`\mathrm{𝗌𝗅}(3;)`$ are. There are six roots.
(6.3)
$$\begin{array}{cc}\alpha & 𝐙\\ (2,1)& X_1\\ (1,2)& X_2\\ (1,1)& X_3\\ (2,1)& Y_1\\ (1,2)& Y_2\\ (1,1)& Y_3\end{array}$$
It is convenient to single out the two roots corresponding to $`X_1`$ and $`X_2`$ and give them special names:
$`\alpha ^{(1)}`$ $`=(2,1)`$
(6.4) $`\alpha ^{(2)}`$ $`=(1,2)\text{.}`$
The roots $`\alpha ^{(1)}`$ and $`\alpha ^{(2)}`$ are called the simple roots. They have the property that all of the roots can be expressed as linear combinations of $`\alpha ^{(1)}`$ and $`\alpha ^{(2)}`$ with integer coefficients, and these coefficients are either all greater than or equal to zero or all less than or equal to zero. This is verified by direct computation:
$$\begin{array}{ccc}(2,1)& =& \alpha ^{(1)}\\ (1,2)& =& \alpha ^{(2)}\\ (1,1)& =& \alpha ^{(1)}+\alpha ^{(2)}\\ (2,1)& =& \alpha ^{(1)}\\ (1,2)& =& \alpha ^{(2)}\\ (1,1)& =& \alpha ^{(1)}\alpha ^{(2)}\text{.}\end{array}$$
The significance of the roots for the representation theory of $`\mathrm{𝗌𝗅}(3;)`$ is contained in the following Lemma. Although its proof is very easy, this Lemma plays a crucial role in the classification of the representations of $`\mathrm{𝗌𝗅}(3;)`$. Note that this Lemma is the analog of Lemma 5.10 of Chapter 5, which was the key to the classification of the representations of $`\mathrm{𝗌𝗅}(2;)`$.
###### Lemma 6.8.
Let $`\alpha =(\alpha _1,\alpha _2)`$ be a root, and $`Z_\alpha 0`$ a corresponding root vector in $`\mathrm{𝗌𝗅}(3;𝐂)`$. Let $`\pi `$ be a representation of $`\mathrm{𝗌𝗅}(3;)`$, $`\mu =(m_1,m_2)`$ a weight for $`\pi `$, and $`v0`$ a corresponding weight vector. Then
$`\pi (H_1)\pi (Z_\alpha )v=(m_1+\alpha _1)\pi (Z_\alpha )v`$
$`\pi (H_2)\pi (Z_\alpha )v=(m_2+\alpha _2)\pi (Z_\alpha )v\text{.}`$
Thus either $`\pi (Z_\alpha )v=0`$ or else $`\pi (Z_\alpha )v`$ is a new weight vector with weight
$$\mu +\alpha =(m_1+\alpha _1,m_2+\alpha _2)\text{.}$$
###### Proof.
The definition of a root tells us that we have the commutation relation $`\left[H_1Z_\alpha \right]=\alpha _1Z_\alpha `$. Thus
$`\pi (H_1)\pi (Z_\alpha )v=\left(\pi (Z_\alpha )\pi (H_1)+\alpha _1\pi (Z_a)\right)v`$
$`=\pi (Z_\alpha )(m_1v)+\alpha _1\pi (Z_\alpha )v`$
$`=(m_1+\alpha _1)\pi (Z_\alpha )v\text{.}`$
A similar argument allows us to compute $`\pi (H_2)\pi (Z_\alpha )v`$. ∎
### 6.3. Highest Weights and the Classification Theorem
We see then that if we have a representation with a weight $`\mu =(m_1,m_2)`$, then by applying the root vectors $`X_1,X_2,X_3,Y_1,Y_2,Y_3`$ we can get some new weights of the form $`\mu +\alpha `$, where $`\alpha `$ is the root. Of course, some of the weight vectors may simply give zero. In fact, since our representation is finite-dimensional, there can be only finitely many weights, so we must get zero quite often. By analogy to the classification of the representations of $`\mathrm{𝗌𝗅}(2;𝐂)`$, we would like to single out in each representation a “highest” weight, and then work from there. The following definition gives the “right” notion of highest.
###### Definition 6.9.
Let $`\alpha ^{(1)}=(2,1)`$ and $`\alpha ^{(2)}=(1,2)`$ be the roots introduced in (6.4). Let $`\mu _1`$ and $`\mu _2`$ be two weights. Then $`\mu _1`$ is higher than $`\mu _2`$ (or equivalently, $`\mu _2`$ is lower than $`\mu _1`$) if $`\mu _1\mu _2`$ can be written in the form
$$\mu _1\mu _2=a\alpha ^{(1)}+b\alpha ^{(2)}$$
with $`a0`$ and $`b0`$. This relationship is written as $`\mu _1\mu _2`$ or $`\mu _2\mu _1`$.
If $`\pi `$ is a representation of $`\mathrm{𝗌𝗅}(3;)`$, then a weight $`\mu _0`$ for $`\pi `$ is said to be a highest weight if for all weights $`\mu `$ of $`\pi `$, $`\mu \mu _0`$.
Note that the relation of “higher” is only a partial ordering. That is, one can easily have $`\mu _1`$ and $`\mu _2`$ such that $`\mu _1`$ is neither higher nor lower than $`\mu _2`$. For example, $`\alpha ^{(1)}\alpha ^{(2)}`$ is neither higher nor lower than $`0`$. This in particular means that a finite set of weights need not have a highest element. (E.g., the set $`\{0,\alpha ^{(1)}\alpha ^{(2)}\}`$ has no highest element.)
We are now ready to state the main theorem regarding the irreducible representations of $`\mathrm{𝗌𝗅}(3;)`$.
###### Theorem 6.10.
1. Every irreducible representation $`\pi `$ of $`\mathrm{𝗌𝗅}(3;)`$ is the direct sum of its weight spaces. That is, $`\pi (H_1)`$ and $`\pi (H_2)`$ are simultaneously diagonalizable.
2. Every irreducible representation of $`\mathrm{𝗌𝗅}(3;)`$ has a unique highest weight $`\mu _0`$, and two equivalent irreducible representations have the same highest weight.
3. Two irreducible representations of $`\mathrm{𝗌𝗅}(3;)`$ with the same highest weight are equivalent.
4. If $`\pi `$ is an irreducible representation of $`\mathrm{𝗌𝗅}(3;)`$, then the highest weight $`\mu _0`$ of $`\pi `$ is of the form
$$\mu _0=(m_1,m_2)$$
with $`m_1`$ and $`m_2`$ non-negative integers.
5. Conversely, if $`m_1`$ and $`m_2`$ are non-negative integers, then there exists a unique irreducible representation $`\pi `$ of $`\mathrm{𝗌𝗅}(3;)`$ with highest weight $`\mu _0=(m_1,m_2)`$.
Note the parallels between this result and the classification of the irreducible representations of $`\mathrm{𝗌𝗅}(2;)`$: In each irreducible representation of $`\mathrm{𝗌𝗅}(2;)`$, $`\pi (H)`$ is diagonalizable, and there is a largest eigenvalue of $`\pi (H)`$. Two irreducible representations of $`\mathrm{𝗌𝗅}(2;)`$ with the same largest eigenvalue are equivalent. The highest eigenvalue is always a non-negative integer, and conversely, for every non-negative integer $`m`$, there is an irreducible representation with highest eigenvalue $`m`$.
However, note that in the classification of the representations of $`\mathrm{𝗌𝗅}(3;)`$ the notion of “highest” does not mean what we might have thought it should mean. For example, the weight $`(1,1)`$ is higher than the weights $`(1,2)`$ and $`(2,1)`$. (In fact, $`(1,1)`$ is the highest weight for the adjoint representation, which is irreducible.)
It is possible to obtain much more information about the irreducible representations besides the highest weight. For example, we have the following formula for the dimension of the representation with highest weight $`(m_1,m_2)`$.
###### Theorem 6.11.
The dimension of the irreducible representation with highest weight $`(m_1,m_2)`$ is
$$\frac{1}{2}(m_1+1)(m_2+1)(m_1+m_2+2)\text{.}$$
We will not prove this formula. It is a consequence of the “Weyl character formula.” See Humphreys, Section 24.3. Humphreys refers to $`\mathrm{𝗌𝗅}(3;)`$ as $`A_2`$.
### 6.4. Proof of the Classification Theorem
It will take us some time to prove Theorem 1. The proof will consist of a series of Propositions.
###### Proposition 6.12.
In every irreducible representation $`(\pi ,V)`$ of $`\mathrm{𝗌𝗅}(3;)`$, $`\pi (H_1)`$ and $`\pi (H_2)`$ can be simultaneously diagonalized. That is, $`V`$ is the direct sum of its weight spaces.
###### Proof.
Let $`W`$ be the direct sum of the weight spaces in $`V`$. Equivalently, $`W`$ is the space of all vectors $`wV`$ such that $`w`$ can be written as a linear combination of simultaneous eigenvectors for $`\pi (H_1)`$ and $`\pi (H_2)`$. Since (Proposition 6.4) $`\pi `$ always has at least one weight, $`W\left\{0\right\}`$.
On the other hand, Lemma 6.8 tells us that if $`Z_\alpha `$ is a root vector corresponding to the root $`\alpha `$, then $`\pi (Z_\alpha )`$ maps the weight space corresponding to $`\mu `$ into the weight space corresponding to $`\mu +\alpha `$. Thus $`W`$ is invariant under the action of all of the root vectors, namely, under the action $`X_1,X_2,X_3,Y_1,Y_2,`$ and $`Y_3`$. Since $`W`$ is certainly invariant under the action of $`H_1`$ and $`H_2`$, $`W`$ is invariant. Thus by irreducibility, $`W=V`$. ∎
###### Definition 6.13.
A representation $`(\pi ,V)`$ of $`\mathrm{𝗌𝗅}(3;)`$ is said to be a highest weight cyclic representation with weight $`\mathrm{\mu }_0=(\mathrm{m}_1,\mathrm{m}_2)`$ if there exists $`v0`$ in $`V`$ such that
1. $`v`$ is a weight vector with weight $`\mu _0`$.
2. $`\pi (X_1)v=\pi (X_2)v=0`$.
3. The smallest invariant subspace of $`V`$ containing $`v`$ is all of $`V`$.
The vector $`v`$ is called a cyclic vector for $`\pi `$.
###### Proposition 6.14.
Let $`(\pi ,V)`$ be a highest weight cyclic representation of $`\mathrm{𝗌𝗅}(3;)`$ with weight $`\mu _0`$. Then
1. $`\pi `$ has highest weight $`\mu _0`$.
2. The weight space corresponding to the highest weight $`\mu _0`$ is one-dimensional.
##### Proof
###### Proof.
Let $`v`$ be as in the definition. Consider the subspace $`W`$ of $`V`$ spanned by elements of the form
(6.5)
$$w=\pi (Y_{i_1})\pi (Y_{i_2})\mathrm{}\pi (Y_{i_n})v$$
with each $`i_l=1,2`$, and $`n0`$. (If $`n=0`$, it is understood that $`w`$in (6.5) is equal to $`v`$.) I assert that $`W`$ is invariant. To see this, it suffices to check that $`W`$ is invariant under each of the basis elements.
By definition, $`W`$ is invariant under $`\pi (Y_1)`$ and $`\pi (Y_2)`$. It is thus also invariant under $`\pi (Y_3)=[\pi (Y_1),\pi (Y_2)]`$.
Now, Lemma 6.8 tells us that applying a root vector $`Z_\alpha \mathrm{𝗌𝗅}(3;)`$ to a weight vector $`v`$ with weight $`\mu `$ gives either zero, or else a new weight vector with weight $`\mu +\alpha `$. Now, by assumption, $`v`$ is a weight vector with weight $`\mu _0`$. Furthermore, $`Y_1`$ and $`Y_2`$ are root vectors with roots $`\alpha ^{(1)}=(2,1)`$ and $`\alpha ^{(2)}=(1,2)`$, respectively. (See Equation (6.3).) Thus each application of $`\pi (Y_1)`$ or $`\pi (Y_2)`$ subtracts $`\alpha ^{(1)}`$ or $`\alpha ^{(2)}`$ from the weight. In particular, each non-zero element of the form (6.5) is a simultaneous eigenvector for $`\pi (H_1)`$ and $`\pi (H_2)`$. Thus $`W`$ is invariant under $`\pi (H_1)`$ and $`\pi (H_2)`$.
To show that $`W`$ is invariant under $`\pi (X_1)`$ and $`\pi (X_2)`$, we argue by induction on $`n`$. For $`n=0`$, we have $`\pi (X_1)v=\pi (X_2)v=0W`$. Now consider applying $`\pi (X_1)`$ or $`\pi (X_2)`$ to a vector of the form (6.5). Recall the commutation relations involving an $`X_1`$ or $`X_2`$ and a $`Y_1`$ or $`Y_2`$:
$$\begin{array}{ccccccc}[X_1,Y_1]& =& H_1& & [X_1,Y_2]& =& 0\\ [X_2,Y_1]& =& 0& & [X_2,Y_2]& =& H_2\text{.}\end{array}$$
Thus (for $`i`$ and $`j`$ equal to 1 or 2) $`\pi (X_i)\pi (Y_j)=\pi (Y_j)\pi (X_i)+\pi (H_{ij})`$, where $`H_{ij}`$ is either $`H_1`$ or $`H_2`$ or zero. Hence (for $`i`$ equal to 1 or 2)
$`\pi (X_i)\pi (Y_{i_1})\pi (Y_{i_2})\mathrm{}\pi (Y_{i_n})v`$
$`=\pi (Y_{i_1})\pi (X_i)\pi (Y_{i_2})\mathrm{}\pi (Y_{i_n})v+\pi (H_{ij})\pi (Y_{i_2})\mathrm{}\pi (Y_{i_n})v\text{.}`$
But $`\pi (X_i)\pi (Y_{i_2})\mathrm{}\pi (Y_{i_n})v`$ is in $`W`$ by induction, and $`\pi (H_{ij})\pi (Y_{i_2})\mathrm{}\pi (Y_{i_n})v`$ is in $`W`$ since $`W`$ is invariant under $`\pi (H_1)`$ and $`\pi (H_2)`$.
Finally, $`W`$ is invariant under $`\pi (X_3)`$ since $`\pi (X_3)=[\pi (X_1),\pi (X_2)]`$. Thus $`W`$ is invariant. Since by definition $`W`$ contains $`v`$, we must have $`W=V`$.
Since $`Y_1`$ is a root vector with root $`\alpha ^{(1)}`$ and $`Y_2`$ is a root vector with root $`\alpha ^{(2)}`$, Lemma 6.8 tells us that each element of the form (6.5) is either zero or a weight vector with weight $`\mu _0\alpha ^{(i_1)}\mathrm{}\alpha ^{(i_n)}`$. Thus $`V=W`$ is spanned by $`v`$ together with weight vectors with weights lower than $`\mu _0`$. Thus $`\mu _0`$ is the highest weight for $`V`$.
Furthermore,every element of $`W`$ can be written as a multiple of $`v`$ plus a linear combination of weight vectors with weights lower than $`\mu _0`$. Thus the weight space corresponding to $`\mu _0`$ is spanned by $`v`$; that is, the weight space corresponding to $`\mu _0`$ is one-dimensional. ∎
###### Proposition 6.15.
Every irreducible representation of $`\mathrm{𝗌𝗅}(3;)`$ is a highest weight cyclic representation, with a unique highest weight $`\mu _0`$.
###### Proof.
Uniqueness is immediate, since by the previous Proposition, $`\mu _0`$ is the highest weight, and two distinct weights cannot both be highest.
We have already shown that every irreducible representation is the direct sum of its weight spaces. Since the representation is finite-dimensional, there can be only finitely many weights. It follows that there must exist a weight $`\mu _0`$ such that there is no weight $`\mu \mu _0`$ with $`\mu \mu _0`$. This says that there is no weight higher than $`\mu _0`$ (which is not the same as saying the $`\mu _0`$ is highest). But if there is no weight higher than $`\mu _0`$, then for any non-zero weight vector $`v`$ with weight $`\mu _0`$, we must have
$$\pi (X_1)v=\pi (X_2)v=0\text{.}$$
(For otherwise, say, $`\pi (X_1)v`$ will be a weight vector with weight $`\mu _0+\alpha ^{(1)}\mu _0`$.)
Since $`\pi `$ is assumed irreducible, the smallest invariant subspace containing $`v`$ must be the whole space; therefore the representation is highest weight cyclic. $`\mathrm{}`$
###### Proposition 6.16.
Every highest weight cyclic representation of $`\mathrm{𝗌𝗅}(3;)`$ is irreducible.
###### Proof.
Let $`(\pi ,V)`$ be a highest weight cyclic representation with highest weight $`\mu _0`$ and cyclic vector $`v`$. By complete reducibility (Proposition 6.2), $`V`$ decomposes as a direct sum of irreducible representations
(6.6)
$$V\underset{i}{}V_i\text{.}$$
By Proposition 6.12, each of the $`V_i`$’s is the direct sum of its weight spaces. Thus since the weight $`\mu _0`$ occurs in $`V`$, it must occur in some $`V_i`$. On the other hand, Proposition 6.14 says that the weight space corresponding to $`\mu _0`$ is one-dimensional, that is, $`v`$ is (up to a constant) the only vector in $`V`$ with weight $`\mu _0`$. Thus $`V_i`$ must contain $`v`$. But then that $`V_i`$ is an invariant subspace containing $`v`$, so $`V_i=V`$. Thus there is only one term in the sum (6.6), and $`V`$ is irreducible. ∎
###### Proposition 6.17.
Two irreducible representations of $`\mathrm{𝗌𝗅}(3;)`$ with the same highest weight are equivalent.
###### Proof.
We now know that a representation is irreducible if and only if it is highest weight cyclic. Suppose that $`(\pi ,V)`$ and $`(\sigma ,W)`$ are two such representations with the same highest weight $`\mu _0`$. Let $`v`$ and $`w`$ be the cyclic vectors for $`V`$ and $`W`$, respectively. Now consider the representation $`VW`$, and let $`U`$ be smallest invariant subspace of $`VW`$ which contains the vector $`(v,w)`$.
By definition, $`U`$ is a highest weight cyclic representation, therefore irreducible by Proposition. 6.16. Consider the two “projection” maps $`P_1:VWV`$, $`P_1(v,w)=v`$ and $`P_2:VWW`$, $`P_1(v,w)=w`$. It is easy to check that $`P_1`$ and $`P_2`$ are morphisms of representations. Therefore the restrictions of $`P_1`$ and $`P_2`$ to $`UVW`$ will also be morphisms.
Clearly neither $`P_1|_U`$ nor $`P_2|_U`$ is the zero map (since both are non-zero on $`(v,w)`$). Moreover, $`U`$, $`V`$, and $`W`$ are all irreducible. Therefore, by Schur’s Lemma, $`P_1|_U`$ is an isomorphism of $`U`$ with $`V`$, and $`P_2|_U`$ is an isomorphism of $`U`$ with $`W`$. Thus $`VUW`$. ∎
###### Proposition 6.18.
If $`\pi `$ is an irreducible representation of $`\mathrm{𝗌𝗅}(3;)`$, then the highest weight of $`\pi `$ is of the form
$$\mu =(m_1,m_2)$$
with $`m_1`$ and $`m_2`$ non-negative integers.
###### Proof.
We already know that all of the weights of $`\pi `$ are of the form $`(m_1,m_2)`$, with $`m_1`$ and $`m_2`$ integers. We must show that if $`\mu _0=(m_1,m_2)`$ is the highest weight, then $`m_1`$ and $`m_2`$ are both non-negative. For this, we again use what we know about the representations of $`\mathrm{𝗌𝗅}(2;)`$. The following result can be obtained from the proof of the classification of the irreducible representations of $`\mathrm{𝗌𝗅}(2;)`$.
Let $`(\pi ,V)`$ be any finite-dimensional representation of $`\mathrm{𝗌𝗅}(2;)`$. Let $`v`$ be an eigenvector for $`\pi (H)`$ with eigenvalue $`\lambda `$. If $`\pi (X)v=0`$, then $`\lambda `$ is a non-negative integer.
Now, if $`\pi `$ is an irreducible representation of $`\mathrm{𝗌𝗅}(3;𝐂)`$ with highest weight $`\mu _0=(m_1,m_2)`$, and if $`v0`$ is a weight vector with weight $`\mu _0`$, then we must have $`\pi (X_1)v=\pi (X_2)v=0`$. (Otherwise, $`\mu _0`$ wouldn’t be highest.) Thus applying the above result to the restrictions of $`\pi `$ to $`\{H_1,X_1,Y_1\}`$ and to $`\{H_2,X_2,Y_2\}`$ shows that $`m_1`$ and $`m_2`$ must be non-negative. ∎
###### Proposition 6.19.
If $`m_1`$ and $`m_2`$ are non-negative integers, then there exists an irreducible representation of $`\mathrm{𝗌𝗅}(3;)`$ with highest weight $`\mu =(m_1,m_2)`$.
###### Proof.
Note that the trivial representation is an irreducible representation with highest weight $`(0,0)`$. So we need only construct representations with at least one of $`m_1`$ and $`m_2`$ positive.
First, we construct two irreducible representations with highest weights $`(1,0)`$ and $`(0,1)`$. (These are the so-called fundamental representations.) The standard representation of $`\mathrm{𝗌𝗅}(3;)`$ is an irreducible representation with highest weight $`(1,0)`$, as is easily checked. To construct an irreducible representation with weight $`(0,1)`$ we modify the standard representation. Specifically, we define
(6.7)
$$\pi (Z)=Z^{tr}$$
for all $`Z\mathrm{𝗌𝗅}(3;)`$. Using the fact that $`\left(AB\right)^{tr}=B^{tr}A^{tr}`$, it is easy to check that
$$[Z_1,Z_2]^{tr}=[Z_1^{tr},Z_2^{tr}]$$
so that $`\pi `$ is really a representation. (This is isomorphic to the dual of the standard representation, as defined in Exercise 14 of Chapter 5.) It is easy to see that $`\pi `$ is an irreducible representation with highest weight $`(0,1)`$.
Let $`(\pi _1,V_1)`$ denote $`^3`$ acted on by the standard representation, and let $`v_1`$ denote a weight vector corresponding to the highest weight $`(1,0)`$. (So, $`v_1=(1,0,0)`$.) Let $`(\pi _2,V_2)`$ denote $`^3`$ acted on by the representation (6.7), and let $`v_2`$ denote a weight vector for the highest weight $`(0,1)`$. (So, $`v_2=(0,0,1)`$.) Now consider the representation
$$V_1V_1\mathrm{}V_1V_2V_2\mathrm{}V_2$$
where $`V_1`$ occurs $`m_1`$ times, and $`V_2`$ occurs $`m_2`$ times. Note that the action of $`\mathrm{𝗌𝗅}(3;)`$ on this space is
$`Z\left(\pi _1(Z)I\mathrm{}I\right)`$
(6.8) $`+\left(I\pi _1(Z)I\mathrm{}I\right)+\mathrm{}+\left(I\mathrm{}I\pi _2(Z)\right)\text{.}`$
Let $`\pi _{m_1,m_2}`$ denote this representation.
Consider the vector
$$v_{m_1,m_2}=v_1v_1\mathrm{}v_1v_2v_2\mathrm{}v_2\text{.}$$
Then applying (6.8) shows that
$`\pi _{m_1,m_2}(H_1)v_{m_1,m_2}=m_1v_{m_1,m_2}`$
$`\pi _{m_1,m_2}(H_2)v_{m_1,m_2}=m_2v_{m_1,m_2}`$
$`\pi _{m_1,m_2}(X_1)v_{m_1,m_2}=0`$
(6.9) $`\pi _{m_1,m_2}(X_2)v_{m_1,m_2}=0\text{.}`$
Now, the representation $`\pi _{m_1,m_2}`$ is not irreducible (unless $`(m_1,m_2)=(1,0)`$ or $`(0,1)`$). However, if we let $`W`$ denote the smallest invariant subspace containing the vector $`v_{m_1,m_2}`$, then in light of (6.9), $`W`$ will be highest weight cyclic with highest weight $`(m_1,m_2)`$. Therefore by Proposition 6.16, $`W`$ is irreducible with highest weight $`(m_1,m_2)`$.
Thus $`W`$ is the representation we want. ∎
We have now completed the proof of Theorem 1.
### 6.5. An Example: Highest Weight $`(1,1)`$
To obtain the irreducible representation with highest weight $`(1,1)`$ we are supposed to take the tensor product of the irreducible representations with highest weights $`(1,0)`$ and $`(0,1)`$, and then extract a certain invariant subspace. Let us establish some notation for the representations $`(1,0)`$ and $`(0,1)`$. In the standard representation, the weight vectors for
$$\begin{array}{cc}H_1=\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 0\end{array}\right);& H_2=\left(\begin{array}{ccc}0& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right);\end{array}$$
are the standard basis elements for $`^3`$, namely, $`e_1`$, $`e_2`$, and $`e_3`$. The corresponding weights are $`(1,0)`$, $`(1,1)`$, and $`(0,1)`$. The highest weight is $`(1,0)`$.
Recall that
$$\begin{array}{cc}Y_1=\left(\begin{array}{ccc}0& 0& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right);& Y_2=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 0\\ 0& 1& 0\end{array}\right)\text{.}\end{array}$$
Thus
(6.10)
$$\begin{array}{ccccccc}Y_1(e_1)& =& e_2& & Y_2(e_1)& =& 0\\ Y_1(e_2)& =& 0& & Y_2(e_2)& =& e_3\\ Y_1(e_3)& =& 0& & Y_2(e_3)& =& 0\text{.}\end{array}$$
Now, the representation with highest weight $`(0,1)`$ is the representation $`\pi (Z)=Z^{tr}`$, for $`Z\mathrm{𝗌𝗅}(3;)`$. Let us define
$$\overline{Z}=Z^{tr}$$
for all $`Z\mathrm{𝗌𝗅}(3;)`$. Thus $`\pi (Z)=\overline{Z}`$. Note that
$$\begin{array}{cc}\overline{H_1}=\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 0\end{array}\right);& \overline{H_2}=\left(\begin{array}{ccc}0& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right)\text{.}\end{array}$$
The weight vectors are again $`e_1`$, $`e_2`$, and $`e_3`$, with weights $`(1,0)`$, $`(1,1)`$, and $`(0,1)`$. The highest weight is $`(0,1)`$.
Define new basis elements
$$\begin{array}{ccc}f_1& =& e_3\\ f_2& =& e_2\\ f_3& =& e_1\text{.}\end{array}$$
Then since
$$\begin{array}{cc}\overline{Y_1}=\left(\begin{array}{ccc}0& 1& 0\\ 0& 0& 0\\ 0& 0& 0\end{array}\right);& \overline{Y_2}=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 1\\ 0& 0& 0\end{array}\right);\end{array}$$
we have
(6.11)
$$\begin{array}{ccccccc}\overline{Y_1}(f_1)& =& 0& & \overline{Y_2}(f_1)& =& f_2\\ \overline{Y_1}(f_2)& =& f_3& & \overline{Y_2}(f_2)& =& 0\\ \overline{Y_1}(f_3)& =& 0& & \overline{Y_2}(f_3)& =& 0\text{.}\end{array}$$
Note that the highest weight vector is $`f_1=e_3`$.
So, to obtain an irreducible representation with highest weight $`(1,1)`$ we are supposed to take the tensor product of the representations with highest weights $`(1,0)`$ and $`(0,1)`$, and then take the smallest invariant subspace containing the vector $`e_1f_1`$. In light of the proof of Proposition 6.14, this smallest invariant subspace is obtained by starting with $`e_1f_1`$ and applying all possible combinations of $`Y_1`$ and $`Y_2`$.
Recall that if $`\pi _1`$ and $`\pi _2`$ are two representations of the Lie algebra $`\mathrm{𝗌𝗅}(3;)`$, then
$`\left(\pi _1\pi _2\right)(Y_1)`$ $`=\pi _1(Y_1)I+I\pi _2(Y_1)`$
$`\left(\pi _1\pi _2\right)(Y_2)`$ $`=\pi _1(Y_2)I+I\pi _2(Y_2)\text{.}`$
In our case we want $`\pi _1(Y_i)=Y_i`$ and $`\pi _2(Y_i)=\overline{Y_i}`$. Thus
$`\left(\pi _1\pi _2\right)(Y_1)`$ $`=Y_1I+I\overline{Y_1}`$
$`\left(\pi _1\pi _2\right)(Y_2)`$ $`=Y_2I+I\overline{Y_2}\text{.}`$
The actions of $`Y_i`$ and $`\overline{Y_i}`$ are described in (6.10) and (6.11).
Note that $`\pi _1\pi _2`$ is not an irreducible representation. The representation $`\pi _1\pi _2`$ has dimension 9, whereas the smallest invariant subspace containing $`e_1f_1`$ has, as it turns out, dimension 8.
So, it remains only to begin with $`e_1f_1`$, apply $`Y_1`$ and $`Y_2`$ repeatedly until we get zero, and then figure out what dependence relations exist among the vectors we get. These computations are done on a supplementary page. Note that the weight $`(0,0)`$ has multiplicity two. This is because, starting with $`e_1f_1`$, applying $`Y_1`$ and then $`Y_2`$ gives something different than applying $`Y_2`$ and then $`Y_1`$.
### 6.6. The Weyl Group
The set of weights of an arbitrary irreducible representation of $`\mathrm{𝗌𝗅}(3;)`$ has a certain symmetry associated to it. This symmetry is in terms of something called the “Weyl group.” (My treatment of the Weyl group follows Bröcker and tom Dieck, Chap. IV, 1.3.) We consider the following subgroup of $`\mathrm{𝖲𝖴}(3)`$:
$$W=\left\{\begin{array}{ccc}w_0=\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right);& w_1=\left(\begin{array}{ccc}0& 0& 1\\ 1& 0& 0\\ 0& 1& 0\end{array}\right);& w_2=\left(\begin{array}{ccc}0& 1& 0\\ 0& 0& 1\\ 1& 0& 0\end{array}\right)\\ w_3=\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 1\end{array}\right);& w_4=\left(\begin{array}{ccc}0& 0& 1\\ 0& 1& 0\\ 1& 0& 0\end{array}\right);& w_5=\left(\begin{array}{ccc}1& 0& 0\\ 0& 0& 1\\ 0& 1& 0\end{array}\right)\end{array}\right\}\text{.}$$
These are simply the matrices which permute the standard basis elements of $`^3`$, with an adjustment of overall sign when necessary to make the determinant equal one.
Now, for any $`A\mathrm{𝖲𝖴}(3)`$, we have the associated map $`\mathrm{Ad}A:\mathrm{𝗌𝗎}(3)\mathrm{𝗌𝗎}(3)`$, where
$$\mathrm{Ad}A(X)=AXA^1\text{.}$$
Now, since each element of $`\mathrm{𝗌𝗅}(3;)`$ is of the form $`Z=X+iY`$ with $`X,Y\mathrm{𝗌𝗎}(3)`$, it follows that $`\mathrm{𝗌𝗅}(3;)`$ is invariant under the map $`ZAZA^1`$. That is, we can think of $`\mathrm{Ad}A`$ as a map of $`\mathrm{𝗌𝗅}(3;)`$ to itself.
The reason for selecting the above group is the following: If $`wW`$, then $`\mathrm{Ad}w(H_1)`$ and $`\mathrm{Ad}w(H_2)`$ are linear combinations of $`H_1`$ and $`H_2`$. That is, each $`\mathrm{Ad}w`$ preserves the space spanned by $`H_1`$ and $`H_2`$. (There are other elements of $`\mathrm{𝖲𝖴}(3)`$ with this property, notably, the diagonal elements. However, these actually commute with $`H_1`$ and $`H_2`$. Thus the adjoint action of these elements on the span of $`H_1`$ and $`H_2`$ is trivial and therefore uninteresting. See Exercise 3.)
Now, for each $`wW`$ and each irreducible representation $`\pi `$ of $`\mathrm{𝗌𝗅}(3;)`$, let’s define a new representation $`\pi _w`$ by the formula
$$\pi _w(X)=\pi \left(\mathrm{Ad}w^1(X)\right)=\pi (w^1Xw)\text{.}$$
Since $`\mathrm{Ad}w^1`$ is a Lie algebra automorphism, $`\pi _w`$ will in fact be a representation of $`\mathrm{𝗌𝗅}(3;)`$.
Recall that since $`\mathrm{𝖲𝖴}(3)`$ is simply connected, then for each representation $`\pi `$ of $`\mathrm{𝗌𝗅}(3;)`$ there is an associated representation $`\mathrm{\Pi }`$ of $`\mathrm{𝖲𝖴}(3)`$ (acting on the same space) such that
$$\mathrm{\Pi }\left(e^X\right)=e^{\pi (X)}$$
for all $`X\mathrm{𝗌𝗎}(3)\mathrm{𝗌𝗅}(3;)`$. The representation $`\mathrm{\Pi }`$ has the property that
(6.12)
$$\pi (AXA^1)=\mathrm{\Pi }(A)\pi (X)\mathrm{\Pi }(A)^1$$
for all $`X\mathrm{𝗌𝗎}(3)`$. Again since every element of $`\mathrm{𝗌𝗅}(3;)`$ is of the form $`X+iY`$ with $`X,Y\mathrm{𝗌𝗎}(3)`$, it follows that (6.12) holds also for $`X\mathrm{𝗌𝗅}(3;)`$.
In particular, taking $`A=w^1W`$ we have
(6.13)
$$\pi _w(X)=\pi (w^1Xw)=\mathrm{\Pi }(w)^1\pi (X)\mathrm{\Pi }(w)$$
for all $`X\mathrm{𝗌𝗅}(3;)`$.
###### Proposition 6.20.
For each representation $`\pi `$ of $`\mathrm{𝗌𝗅}(3;)`$ and for each $`wW`$, the representation $`\pi _w`$ is equivalent to the representation $`\pi `$.
###### Proof.
We need a map $`\varphi :VV`$ with the property that
$$\varphi \left(\pi _w(X)v\right)=\pi (X)\varphi (v)$$
for all $`vV`$. This is the same as saying that $`\varphi \pi _w(X)=\pi (X)\varphi `$, or equivalently that $`\pi _w(X)=\varphi ^1\pi (X)\varphi `$. But in light of (6.13), we can take $`\varphi =\mathrm{\Pi }(w)`$. ∎
Although $`\pi `$ and $`\pi _w`$ are equivalent, they are not equal. That is, in general $`\pi (X)\pi _w(X)`$. You should think of $`\pi `$ and $`\pi _w`$ as differing by a change of basis on $`V`$, where the change-of-basis matrix is $`\mathrm{\Pi }(w)`$. Two representations that differ just by a change of basis are automatically equivalent.
###### Corollary 6.21.
Let $`\pi `$ be a representation of $`\mathrm{𝗌𝗅}(3;)`$ and $`wW`$. Then a pair $`\mu =(m_1,m_2)`$ is a weight for $`\pi `$ if and only if it is a weight for $`\pi _w`$. The multiplicity of $`\mu `$ as a weight of $`\pi `$ is the same as the multiplicity of $`\mu `$ as a weight for $`\pi _w`$.
###### Proof.
Equivalent representations must have the same weights and the same multiplicities. ∎
Let us now compute explicitly the action of $`\mathrm{Ad}w^1`$ on the span of $`H_1`$ and $`H_2`$, for each $`wW`$. This is a straightforward computation.
(6.14)
$$\begin{array}{ccccccc}w_0^1H_1w_0& =& H_1& & w_3^1H_1w_3& =& H_1\\ w_0^1H_2w_0& =& H_2& & w_3^1H_2w_3& =& H_1+H_2\\ & & & & & & \\ w_1^1H_1w_1& =& H_1H_2& & w_4^1H_1w_4& =& H_2\\ w_1^1H_2w_1& =& H_1& & w_4^1H_2w_4& =& H_1\\ & & & & & & \\ w_2^1H_1w_2& =& H_2& & w_5^1H_1w_5& =& H_1+H_2\\ w_2^1H_2w_2& =& H_1H_2& & w_5^1H_2w_5& =& H_2\text{.}\end{array}$$
We can now see the significance of the Weyl group. Let $`\pi `$ be a representation of $`\mathrm{𝗌𝗅}(3;)`$, $`\mu =(m_1,m_2)`$ a weight, and $`v0`$ a weight vector with weight $`\mu `$. Then, for example,
$`\pi _{w_1}(H_1)v`$ $`=\pi (w_1^1H_1w_1)v=\pi (H_1H_2)v=(m_1m_2)v`$
$`\pi _{w_1}(H_2)v`$ $`=\pi (w_1^1H_2w_1)v=\pi (H_1)v=m_1v\text{.}`$
Thus $`v`$ is a weight vector for $`\pi _w`$ with weight $`(m_1m_2,m_1)`$. But by Corollary 6.21, the weights of $`\pi `$ and of $`\pi _w`$ are the same!
> Conclusion: If $`\mu =(m_1,m_2)`$ is a weight for $`\pi `$, so is $`(m_1m_2,m_1)`$. The multiplicities of $`(m_1,m_2)`$ and $`(m_1m_2,m_1)`$ are the same.
Of course, a similar argument applies to each of the other elements of the Weyl group. Specifically, if $`\mu `$ is a weight for some representation $`\pi `$, and $`w`$ is an element of $`W`$, then there will be some new weight which must also be a weight of $`\pi `$. We will denote this new weight $`w\mu `$. For example, if $`\mu =(m_1,m_2)`$, then $`w_1\mu =(m_1m_2,m_1)`$. (We define $`w\mu `$ so that if $`v`$ is a weight vector for $`\pi `$ with weight $`\mu `$, then $`v`$ will be a weight for $`\pi _w`$ with weight $`w\mu `$.) From (6.14) we can read off what $`w\mu `$ is for each $`w`$.
(6.15)
$$\begin{array}{ccccccc}w_0(m_1,m_2)& =& (m_1,m_2)& & w_3(m_1,m_2)& =& (m_1,m_1+m_2)\\ w_1(m_1,m_2)& =& (m_1m_2,m_1)& & w_4(m_1,m_2)& =& (m_2,m_1)\\ w_2(m_1,m_2)& =& (m_2,m_1m_2)& & w_5(m_1,m_2)& =& (m_1+m_2,m_2)\end{array}$$
It is straightforward to check that
(6.16)
$$w_i(w_j\mu )=(w_iw_j)\mu \text{.}$$
We have now proved the following.
###### Theorem 6.22.
If $`\mu =(m_1,m_2)`$ is a weight and $`w`$ is an element of the Weyl group, let $`w\mu `$ be defined by (6.15). If $`\pi `$ is a finite-dimensional representation of $`\mathrm{𝗌𝗅}(3;)`$, then $`\mu `$ is a weight for $`\pi `$ if and only if $`w\mu `$ is a weight for $`\pi `$. The multiplicity of $`\mu `$ is the same as the multiplicity of $`w\mu `$.
If we think of the weights $`\mu =(m_1,m_2)`$ as sitting inside $`^2`$, then we can think of (6.15) as a finite group of linear transformations of $`^2`$. (The fact that this is a group of transformations follows form (6.16).) Since this is a finite group of transformations, it is possible to choose an inner product on $`^2`$ such that the action of $`W`$ is orthogonal. (As in the proof of Proposition 5.16 in Chapter 5.) In fact, there is (up to a constant) exactly one such inner product. In this inner product, the action (6.15) of the Weyl group is generated by a $`120^{}`$ rotation and a reflection about the $`y`$-axis. Equivalently, the Weyl group is the symmetry group of an equilateral triangle centered at the origin with one vertex on the $`y`$-axis.
### 6.7. Complex Semisimple Lie Algebras
This section gives a brief synopsis of the structure theory and representation theory of complex semisimple Lie algebras. The moral of the story is that all such Lie algebras look and feel a lot like $`\mathrm{𝗌𝗅}(3;)`$. This section will not contain any (non-trivial) proofs.
If $`𝔤`$ is a Lie algebra, a subspace $`I𝔤`$ is said to be an ideal if $`[X,Y]I`$ for all $`X𝔤`$ and all $`YI`$. A Lie algebra $`𝔤`$ is a said to be simple if $`dim𝔤2`$ and $`𝔤`$ has no ideals other than $`\left\{0\right\}`$ and $`𝔤`$. A Lie algebra $`𝔤`$ is said to be semisimple if $`𝔤`$ can be written as the direct sum of simple Lie algebras.
In this section we consider semisimple Lie algebras over the complex numbers. Examples of complex semisimple Lie algebras include $`\mathrm{𝗌𝗅}(n;)`$, $`\mathrm{𝗌𝗈}(n;)`$ ($`n3`$), and $`\mathrm{𝗌𝗉}(n;)`$. All of these are actually simple, except for $`\mathrm{𝗌𝗈}(4;)`$ which is isomorphic to $`\mathrm{𝗌𝗅}(2;)\mathrm{𝗌𝗅}(2;)`$.
###### Definition 6.23.
Let $`𝔤`$ be a complex semisimple Lie algebra. A subspace $`𝔥`$ of $`𝔤`$ is said to be a Cartan subalgebra if
1. $`𝔥`$ is abelian. That is, $`[H_1,H_2]=0`$ for all $`H_1,H_2𝔥`$.
2. $`𝔥`$ is maximal abelian. That is, if $`X𝔤`$ satisfies $`[H,X]=0`$ for all $`H𝔥`$, then $`X𝔥`$.
3. For all $`H𝔥`$, $`\mathrm{ad}H:𝔤𝔤`$ is diagonalizable.
Since all the $`H`$’s commute, so do the $`\mathrm{ad}H`$’s. (I.e., $`[\mathrm{ad}H_1,\mathrm{ad}H_2]=\mathrm{ad}[H_1,H_2]=0`$.) By assumption, each $`\mathrm{ad}H`$ is diagonalizable, and they commute, therefore the $`\mathrm{ad}H`$’s are simultaneously diagonalizable. (Using a standard linear algebra fact.) Let $`𝔥^{}`$ denote the dual of $`𝔥`$, namely, the space of linear functionals on $`𝔥`$.
###### Definition 6.24.
If $`𝔤`$ is a complex semisimple Lie algebra and $`𝔥`$ a Cartan subalgebra, then an element $`\alpha `$ of $`𝔥^{}`$ is said to be a root (for $`𝔤`$ with respect to $`𝔥`$) if $`\alpha `$ is non-zero and there exists $`Z0`$ in $`𝔤`$ such that
(6.17)
$$[H,Z]=\alpha (H)Z$$
for all $`H𝔥`$. (Thus a root is a non-zero set of simultaneous eigenvalues for the $`\mathrm{ad}H`$’s.)
The vector $`Z`$ is called a root vector corresponding to the root $`\alpha `$, and the space of all $`Z𝔤`$ satisfying (6.17) is the root space corresponding to $`\alpha `$. This space is denoted $`𝔤^\alpha `$.
The set of all roots will be denoted $`\mathrm{\Delta }`$.
Note that if $`𝔤=\mathrm{𝗌𝗅}(3;)`$, then one Cartan subalgebra is the space spanned by $`H_1`$ and $`H_2`$. The roots (with respect to this Cartan subalgebra) have been calculated in (6.3).
###### Theorem 6.25.
If $`𝔤`$ is a complex semisimple Lie algebra, then a Cartan subalgebra $`𝔥`$ exists. If $`𝔥_1`$ and $`𝔥_2`$ are two Cartan subalgebras, then there is an automorphism of $`𝔤`$ which takes $`𝔥_1`$ to $`𝔥_2`$. In particular, any two Cartan subalgebras have the same dimension.
From now on, $`𝔤`$ will denote a complex semisimple Lie algebra, and $`𝔥`$ a fixed Cartan subalgebra in $`𝔤`$.
###### Definition 6.26.
The rank of a complex semisimple Lie algebra is the dimension of a Cartan subalgebra.
For example, the rank of $`\mathrm{𝗌𝗅}(n;)`$ is $`n1`$. One Cartan subalgebra in $`\mathrm{𝗌𝗅}(n;)`$ is the space of diagonal matrices with trace zero. (Note that in the case $`n=3`$ the space of diagonal matrices with trace zero is precisely the span of $`H_1`$ and $`H_2`$.) Both $`\mathrm{𝗌𝗈}(2n;)`$ and $`\mathrm{𝗌𝗈}(2n+1;)`$ have rank $`n`$.
###### Definition 6.27.
Let $`(\pi ,V)`$ be a finite-dimensional, complex-linear representation of $`𝔤`$. Then $`\mu 𝔥^{}`$ is called a weight for $`\pi `$ if there exists $`v0`$ in $`V`$ such that
$$\pi (H)v=\mu (H)v$$
for all $`H𝔥`$. The vector $`v`$ is called a weight vector for the weight $`\mu `$.
Note that the roots are precisely the non-zero weights for the adjoint representation.
###### Lemma 6.28.
Let $`\alpha `$ be a root and $`Z`$ a corresponding root vector. Let $`\mu `$ be a weight for a representation $`\pi `$ and $`v`$ a corresponding weight vector. Then either $`\pi (Z)v=0`$ or else $`\pi (Z)v`$ is a weight vector with weight $`\mu +\alpha `$.
###### Proof.
Same as for $`\mathrm{𝗌𝗅}(3;𝐂)`$. ∎
###### Definition 6.29.
A set of roots $`\{\alpha _1,\mathrm{}\alpha _l\}`$ is called a simple system (or basis) if
1. $`\{\alpha _1,\mathrm{}\alpha _l\}`$ is a vector space basis for $`𝔥^{}`$.
2. Every root $`\alpha \mathrm{\Delta }`$ can be written in the form
$$\alpha =n_1\alpha _1+n_2\alpha _2+\mathrm{}+n_l\alpha _l$$
with each $`n_i`$ an integer, and such that the $`n_i`$’s are either all non-negative or all non-positive.
A root $`\alpha `$ is said to be positive (with respect to the given simple system) if the $`n_i`$’s are non-negative; otherwise $`\alpha `$ is negative.
If $`𝔤=\mathrm{𝗌𝗅}(3;)`$ and $`𝔥=\{H_1,H_2\}`$, then one simple system of roots is $`\{\alpha ^{(1)},\alpha ^{(2)}\}=\{(2,1),(1,2)\}`$ (with the corresponding root vectors being $`X_1`$ and $`X_2`$). The positive roots are $`\{(2,1),(1,2),(1,1)\}`$. The negative roots are $`\{(2,1),(1,2),(1,1)\}`$.
###### Definition 6.30.
Let $`\{\alpha _1,\mathrm{}\alpha _l\}`$ be a simple system of roots and let $`\mu _1`$ and $`\mu _2`$ be two weights. Then $`\mu _1`$ is higher than $`\mu _2`$ (or $`\mu _2`$ is lower than $`\mu _1`$) if $`\mu _1\mu _2`$ can be written as
$$\mu _1\mu _2=a_1\alpha _1+a_2\alpha _2+\mathrm{}+a_l\alpha _l$$
with $`a_i0.`$ This relation is denoted $`\mu _1\mu _2`$ or $`\mu _2\mu _1`$.
A weight $`\mu _0`$ for a representation $`\pi `$ is highest if all the weights $`\mu `$ of $`\pi `$ satisfy $`\mu \mu _0`$.
The following deep theorem captures much of the structure theory of semisimple Lie algebras.
###### Theorem 6.31.
Let $`𝔤`$ be a complex semisimple Lie algebra, $`𝔥`$ a Cartan subalgebra, and $`\mathrm{\Delta }`$ the set of roots. Then
1. For each root $`\alpha \mathrm{\Delta }`$, the corresponding root space $`𝔤^\alpha `$ is one-dimensional.
2. If $`\alpha `$ is a root, then so is $`\alpha `$.
3. A simple system of roots $`\{\alpha _1,\mathrm{}\alpha _l\}`$ exists.
We now need to identify the correct set of weights to be highest weights of irreducible representations.
###### Theorem 6.32.
Let $`\{\alpha _1,\mathrm{}\alpha _l\}`$ denote a simple system of roots, $`X_i`$ an element of the root space $`𝔤^{\alpha _i}`$ and $`Y_i`$ an element of the root space $`𝔤^{\alpha _i}`$. Define
$$H_i=[X_i,Y_i]\text{.}$$
Then it is possible to choose $`X_i`$ and $`Y_i`$ such that
1. Each $`H_i`$ is non-zero and contained in $`𝔥`$.
2. The span of $`\{H_i,X_i,Y_i\}`$ is a subalgebra of $`𝔤`$ isomorphic (in the obvious way) to $`\mathrm{𝗌𝗅}(2;)`$.
3. The set $`\{H_1,\mathrm{}H_l\}`$ is a basis for $`𝔥`$.
Note that (in most cases) the set of all $`H_i`$’s, $`X_i`$’s, and $`Y_i`$’s ($`i=1,2,\mathrm{}l`$) do not span $`𝔤`$. In the case $`𝔤=\mathrm{𝗌𝗅}(3;)`$, $`l=2`$, and the span of $`H_1,X_1,Y_1,H_2,X_2,Y_2`$ represents only six of the eight dimensions of $`\mathrm{𝗌𝗅}(3;)`$. Nevertheless the subalgebras $`\{H_i,X_i,Y_i\}`$ play an important role.
We are now ready to state the main theorem.
###### Theorem 6.33.
Let $`𝔤`$ be a complex semisimple Lie algebra, $`𝔥`$ a Cartan subalgebra, and $`\{\alpha _1,\mathrm{}\alpha _l\}`$ a simple system of roots. Let $`\{H_1,\mathrm{}H_l\}`$ be as in Theorem 6.32. Then
1. In each irreducible representation $`\pi `$ of $`𝔤`$, the $`\pi (H)`$’s are simultaneously diagonalizable.
2. Each irreducible representation of $`𝔤`$ has a unique highest weight.
3. Two irreducible representations of $`𝔤`$ with the same highest weight are equivalent.
4. If $`\mu _0`$ is the highest weight of an irreducible representation of $`𝔤`$, then for $`i=1,2,\mathrm{}l`$, $`\mu _0(H_i)`$ is a non-negative integer.
5. Conversely, if $`\mu _0𝔥^{}`$ is such that $`\mu _0(H_i)`$ is a non-negative integer for all $`i=1,2,\mathrm{}l`$, then there is an irreducible representation of $`𝔤`$ with highest weight $`\mu _0`$.
The weights $`\mu _0`$ as in 4) and 5) are called dominant integral weights.
### 6.8. Exercises
1. Show that for any pair of $`n\times n`$ matrices $`X`$ and $`Y`$,
$$[X^{tr},Y^{tr}]=[X,Y]^{tr}\text{.}$$
Using this fact and the fact that $`X_i^{tr}=Y_i`$ for $`i=1,2,3`$, explain the symmetry between $`X`$’s and $`Y`$’s in the commutation relations for $`\mathrm{𝗌𝗅}(3;)`$. For example, show that the relation $`[Y_1,Y_2]=Y_3`$ can be obtained from the relation $`[X_1,X_2]=X_3`$ by taking transposes. Show that the relation $`[H_1,Y_2]=Y_2`$ follows from the relation $`[H_1,X_2]=X_2`$.
2. Recall the definition of the dual $`\pi ^{}`$ of a representation $`\pi `$ from Exercise 14 of Chapter 5. Consider this for the case of representations of $`\mathrm{𝗌𝗅}(3;)`$.
a) Show that the weights of $`\pi ^{}`$ are the negatives of the weights of $`\pi `$.
b) Show that if $`\pi `$ is the irreducible representation of $`\mathrm{𝗌𝗅}(3;)`$ with highest weight $`(m_1,m_2)`$ then $`\pi ^{}`$ is the irreducible representation with highest weight $`(m_2,m_1)`$.
Hint: If you identify $`V`$ and $`V^{}`$ by choosing a basis for $`V`$, then $`A^{tr}`$ is just the usual matrix transpose.
3. Let $`𝔥`$ denote the subspace of $`\mathrm{𝗌𝗅}(3;)`$ spanned by $`H_1`$ and $`H_2`$. Let $`G`$ denote the group of all matrices $`A\mathrm{𝖲𝖴}(3)`$ such that $`\mathrm{Ad}A`$ preserves $`𝔥`$. Now let $`G_0`$ denote the group of all matrices $`A\mathrm{𝖲𝖴}(3)`$ such that $`\mathrm{Ad}A`$ is the identity on $`𝔥`$, i.e., such that $`\mathrm{Ad}A(H_1)=H_1`$ and $`\mathrm{Ad}A(H_2)=H_2`$. Show that $`G_0`$ is a normal subgroup of $`G`$. Compute $`G`$ and $`G_0`$. Show that $`G/G_0`$ is isomorphic to the Weyl group $`W`$.
4. a) Verify Theorems 6.31 and 6.32 explicitly for the case $`𝔤=\mathrm{𝗌𝗅}(n;)`$.
b) Consider the task of trying to prove Theorem 6.33 for the case of $`\mathrm{𝗌𝗅}(n;)`$. Now that you have done (a), what part of the proof goes through the same way as for $`\mathrm{𝗌𝗅}(3;)`$? At what points in the proof of the corresponding theorem for $`\mathrm{𝗌𝗅}(3;)`$ did we use special properties of $`\mathrm{𝗌𝗅}(3;)`$?
Hint: Most of it is the same, but there is one critical point which we do something which does not generalize to $`\mathrm{𝗌𝗅}(n;)`$.
## Chapter 7 Cumulative exercises
1. Let $`G`$ be a connected matrix Lie group, and let $`\mathrm{Ad}:G\mathrm{𝖦𝖫}(𝔤)`$ be the adjoint representation of $`G`$. Show that
$$\mathrm{ker}(\mathrm{Ad})=Z(G)$$
where $`Z(G)`$ denotes the center of $`G`$. If $`G=𝖮(2)`$, compute $`\mathrm{ker}(\mathrm{Ad})`$ and $`Z(G)`$ and show that they are not equal.
Hint: You should use the fact that if $`G`$ is connected, then every $`AG`$ can be written in the form $`A=e^{X_1}e^{X_2}\mathrm{}e^{X_n}`$, with $`X_i𝔤`$.
2. Let $`G`$ be a finite, commutative group. Show that the number of equivalence classes of irreducible complex representations of $`G`$ is equal to the number of elements in $`G`$.
Hint: Use the fact that every finite, commutative group is a product of cyclic groups.
3. a) Show that if $`R𝖮(2)`$, and $`detR=1`$, then $`R`$ has two real, orthogonal eigenvectors with eigenvalues $`1`$ and $`1`$.
b) Let $`R`$ be in $`𝖮(n)`$. Show that there exists a subspace $`W`$ of $`^n`$ which is invariant under both $`R`$ and $`R^1`$, and such that $`dimW=1`$ or $`2`$. Show that $`W^{}`$ (the orthogonal complement of $`W`$) is also invariant under $`R`$ and $`R^1`$. Show that the restrictions of $`R`$ and $`R^1`$ to $`W`$ and to $`W^{}`$ are orthogonal. (That is, show that these restrictions preserve inner products.)
c) Let $`R`$ be in $`𝖮(n)`$. Show that $`^n`$ can be written as the orthogonal direct sum of subspaces $`W_i`$ such that
1. 1) Each $`W_i`$ is invariant under $`R`$ and $`R^1`$,
2. 2) Each $`W_i`$ has dimension $`1`$ or $`2`$, and
3. 3) If $`dimW_i=2`$, then the restriction of $`R`$ to $`W_i`$ has determinant one.
d) Show that the exponential mapping for $`\mathrm{𝖲𝖮}(n)`$ is onto. Make sure you use the fact that the elements of $`\mathrm{𝖲𝖮}(n)`$ have determinant one.
Note: This provides an alternative proof that the group $`\mathrm{𝖲𝖮}(n)`$ is connected.
4. Determine, up to equivalence, all of the finite-dimensional, irreducible (complex-linear) representations of the Lie algebra $`\mathrm{𝗌𝗅}(2;)\mathrm{𝗌𝗅}(2;)`$. Can your answer be expressed in terms of a sort of “highest weight”?
Hint: Imitate the proof of the classification of the irreducible representations of $`\mathrm{𝗌𝗅}(2;)`$.
5. Consider the irreducible representation $`(\pi ,V)`$ of $`\mathrm{𝗌𝗅}(3;)`$ with highest weight $`(0,2)`$. Following the procedure in Chapter 6, Section 5, determine
1) The dimension of $`V`$.
2) All of the weights of $`\pi `$.
3) The multiplicity of each of the weights. (That is, the dimension of the corresponding weight spaces.)
## Chapter 8 Bibliography
1. Theodor Bröcker and Tammo tom Dieck, Representations of Compact Lie Groups. Springer-Verlag, 1985.
A good reference for basic facts on compact groups and their representations, including characters and orthogonality relations. Analyzes representations from a more analytic and less algebraic viewpoint than other authors.
2. William Fulton and Joe Harris, Representation theory. A First Course. Graduate Texts in Mathematics, 129. Readings in Mathematics, Springer-Verlag, 1991. Has lots of examples. Written from an algebraic point of view.
3. Sigurdur Helgason, Differential Geometry, Lie Groups, and Symmetric Spaces. Academic Press, 1978.
A good reference for a lot of things. Includes structure theory of semisimple groups.
4. James E. Humphreys, Introduction to Lie Algebras and Representation Theory. Springer-Verlag, 1972.
A standard reference for the Lie algebra side of things (no Lie groups).
5. N. Jacobson, Lie Algebras. Interscience Tracts No. 10, John Wiley and Sons, 1962.
Another good reference for Lie algebras.
6. Anthony W. Knapp, Lie groups: beyond an introduction. Birkhauser, 1996. Good complement to Helgason on such matters as structure theory of Lie groups. As title suggests, not the place to start, but a good reference.
7. W. Miller, Symmetry Groups and Their Applications. Academic Press.
Oriented toward applications to physics. Includes theory of finite groups.
8. Jean-Pierre Serre, Complex Semisimple Lie Algebras. Springer-Verlag, 1987.
A very concise summary of structure theory and representation theory of semisimple Lie algebras.
9. Jean-Pierre Serre, Linear Representations of Finite Groups. Springer-Verlag.
An introduction to both complex and modular representations of finite groups.
10. Barry Simon, Representations of finite and compact Lie groups, American Mathematical Society, 1996. Covers much of the same material as Bröcker and tom Dieck, but from a more analytical perspective.
11. Frank W. Warner, Foundations of Differentiable Manifolds and Lie Groups. Springer-Verlag, 1983.
Key word in the title is foundations. Gives a modern treatment of differentiable manifolds, and then proves some important, non-trivial theorems about Lie groups, including the relationship between subgroups and subalgebras, and the relationship between representations of the Lie algebra and of the Lie group.
12. V.S. Varadarajan, Lie Groups, Lie Algebras, and Their Representations. Springer-Verlag, 1974.
A comprehensive treatment of both Lie groups and Lie algebras.
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# Contact between laboratory instruments and equations of quantum mechanics
## 1 Introduction
To build a quantum computer is to arrange for laboratory instruments to produce results in accord with models expressed in equations of quantum mechanics. These quantum-mechanical models are of two types:
1. instrument-independent models focused on relating the multiplication of unitary operators to the solving of problems of interest, and
2. models that tie the multiplication of unitary operators (as well as steps of state preparation and measurement) to the use of laboratory instruments.
Instruments that could implement a quantum computer would be valuable if their results could substitute for a more costly classical calculation defined by a model of type 1. If a scientist has such instruments, however, using them requires a model of type 2 which tells the scientist how to use the instruments; it is these quantum-mechanical models of how instruments function that are the subject of this report.
Models of type 2 are subject to revision in the course of using instruments, as we found in arranging for a nuclear-magnetic-resonance (NMR) spectrometer to implement a model of an NMR quantum computer. In addition to its electromagnet that holds a liquid sample, the spectrometer has a Classical, digital Process-control Computer (CPC), a variable radio-frequency (r.f.) transmitter controlled by the CPC, and an r.f. receiver that reports back to the CPC which then computes spectra and displays them. Our colleagues synthesized the active constituent of the liquid for a 5-spin quantum computer, aiming to have the spectrometer behave roughly as described by a certain model $`\alpha `$ (of type 2), characterized by a Hamiltonian for 5 linearly-coupled spin-active nuclei in the presence of the variable r.f. field. But while monitoring the chemical synthesis they found spectra that better fit a model $`\beta `$ (also of type 2) characterized by a more complex hamiltonian with one spin-spin coupling weaker than desired and with extra couplings beyond those wanted. So model $`\beta `$ rather than the original model $`\alpha `$ was used to design commands by which to operate the 5-spin quantum computer, showing that implementers have to negotiate with the instruments to refine their models.
Another lessen from the 5-spin endeavor is the secondary role of fields and particles in the design of instruments. We need models of how the spectrometer works to tell us what commands the CPC must issue to the transmitter in order to solve the Deutsch-Jozsa problem, and the historically available constructs and examples from which to construct these models are fields, particles and their couplings; these however are needed only as pieces from which to compose models to describe the instruments. A change from model $`\alpha `$ to model $`\beta `$ changes the couplings, and more severe changes change the particles and fields.
Is the need for alternative models serious? Whether one is working with a mathematical model or working with instruments, it is easy to assume that what goes on in the model, with its fields and particles, also goes on in the instruments, at least once the “right” model is found. But recently it was proved that multiple, inequivalent models to describe a set-up of instruments contend for acceptance not just in the early stages of an experiment but throughout. By assuring many models and hence many configurations of particles and fields to describe a single set-up of instruments, this proof punctures the idea that particles and fields (or, more generally, states and operators) reside in instruments, somehow uniquely situated in them, if only one could see them. This puncture relieves a widespread confusion between models and instruments by demonstrating a certain independence between the two. Given that independence, one has the question: of what does contact between instruments and models consist?
For purposes of an analysis that recognizes their independence, we focus on the contact between quantum-mechanical equations and laboratory instruments that takes place in computer files of a Classical, digital Process-control Computer (CPC). Within a session in which both equation writing and the use of instruments are mediated by the CPC, a scientist uses the CPC to:
* compute (classically!) with various sets of equations of quantum mechanics that define what the instruments are supposed to do, and
* send commands to instruments intended to implement those equations and record results produced by the instruments.
Quantum mechanical models appropriate to the CPC control of instruments are formulated Section 2. Section 3 addresses the a question of finding models from which to determine commands for any particular set-up of instruments, leading to an analysis of sample sizes of experimental tests of models as a function of precision. Consistent with prior estimates, in quantum searching the required sample size rises exponentially with the number of bits.
Section 4 shows lower bounds on the precision of timing needed for a CPC to manage quantum-computing instruments. Section 5 introduces the concept of a lattice of models to show that testing in physics and engineering cannot be universal like the test of a derivation in mathematics; instead, a scientist working with instruments, trying for progressively higher precision in their accord with models, must re-evaluate properties by which to restrict a set of models, requiring repeated resort to guesswork.
## 2 Quantum mechanical models for CPC-controlled instruments
For one example of contact between instruments and equations, suppose a scientist has instruments which would implement a quantum computer if they were sent correct commands; the scientist then faces the question of what commands the CPC should transmit and when it should transmit to them in order e.g. to implement a quantum gate expressed as a unitary matrix $`U_j`$.
To produce a core set of models from which a scientist can choose to describe the working of particular instruments, we suppose that during a CPC-mediated session some instruments are controlled by CPC-transmitted commands from a set $`B`$ of possible commands, where $`B`$ and $``$ is the set of all finite binary strings. Our models express the probability of an outcome of instruments in response to a command $`bB`$ sent to the instruments by the CPC, as follows. Let $``$ be a separable Hilbert space. Let $`𝒱_B`$, $`𝒰_B`$, and $`_B`$ be the sets of all functions $`|v`$, $`U`$, and $`M`$, respectively, with
$`|v:B`$ $``$ $`,`$
$`U:B`$ $``$ $`\{\text{unitary operators on }\},\text{and}`$
$`M:B`$ $``$ $`\{\text{hermitian operators on }\}.`$
The core models exhibit discrete spectra for all $`M`$:
###### Property 1
$$(bB)M(b)=\underset{j}{}m_j(b)M_j(b),$$
(1)
where $`m_j:B`$ (with $``$ denoting the real numbers) is the $`j`$-th eigenvalue of $`M`$, and $`M_j`$ is the projection onto the $`j`$-th eigenspace (so $`M_jM_k=\delta _{j,k}M_j`$).
Let $`\mathrm{Pr}(j|b)`$ denote the probability of obtaining the $`j`$-th outcome, given transmission by the CPC of a command $`b`$. Although not commonly seen in texts, this probability of an outcome given a command is the hinge pin for focusing on quantum mechanical language used to describe what a scientist can find by using instruments. Quantum mechanics constrains all these models to satisfy
###### Property 2
$$\mathrm{Pr}(j|b)=v(b)|U^{}(b)M_j(b)U(b)|v(b),$$
(2)
where the $``$ denotes the hermitian adjoint.
Any choice of command set $`B`$ and of functions from the sets $`𝒱_B`$, $`𝒰_B`$, and $`_B`$ produces a quantum-mechanical model $`(|v,U,M)_B`$. Two models $`(|v,U,M)_B`$ and $`(|v^{},U^{},M^{})_B`$ generate the same probabilities $`\mathrm{Pr}(j|b)`$ if they are unitarily equivalent, meaning there exists a $`Q:B\{\text{unitary operators on }\}`$ such that $`(bB)|v^{}(b)=Q(b)|v(b)`$, $`U^{}(b)=Q(b)U(b)Q^{}(b)`$ and $`M^{}(b)=Q(b)M(b)Q^{}(b)`$. For this reason, any model $`(|v,U,M)_B`$ can be reduced to $`(|v^{},\mathrm{𝟏},M)_B`$, where $`|v^{}=U|v`$ and $`M^{}=M`$ or, alternatively to $`(|v,\mathrm{𝟏},M^{})_B`$ where $`M^{}=U^{}MU`$.
As was recently proved, in order that measured data can select a single best fitting model from a set of models (up to unitary equivalence), additional restrictions are necessary to narrow down the set of models to a much smaller set than that defined by properties 1 and 2, because many inequivalent models can be found to fit exactly and conceivable record of outcomes. Furthermore, it was proved that these restrictions cannot be derived from quantum mechanics nor from the measured data, so that imposing them takes guesswork on the part of the scientist. This core set of models is a subset of models available using more general formulations. Because guesswork is necessary even to resolve choices among models of the core set, it follows that guesswork is necessary also to resolve the broader choices of models from supersets that include more models, e.g. models involving positive-operator-valued measures or superoperators or both.
It is conventional in endeavors aiming at quantum computing to assume three additional properties to narrow the set of models; this is community-endorsed guesswork:
###### Property 3
The command $`b`$ is the concatenation of separate commands for the three types of operations, so that
$$b=b_vb_Ub_M,$$
(3)
where here the $``$ denotes concatenation of commands.
According to these models, one can vary any one of the three while holding the other two fixed. This specializes Eq. (2) to the more restrictive form:
$$\mathrm{Pr}(j|b)=v(b_v)|U^{}(b_U)M_j(b_M)U(b_U)|v(b_v).$$
(4)
An additional constraining guess characterizes models widely used in the analysis of quantum computers, a guess prompted by the desire to generate a unitary transformation as a product of other unitary transformations that serve as “elementary quantum gates.” For example, a scientist may want to generate the unitary transformation $`U(b_{U,1})U(b_{U,2})`$ by causing the CPC to transmit some $`b_U`$. For quantum computing to have an advantage over classical computing, the determination of this $`b_U`$ in terms of $`b_{U,1}`$ and $`b_{U,2}`$ must be of polynomial complexity . It is usually assumed that $`b_U`$ is the simplest possible function of $`b_{U,1}`$ and $`b_{U,2}`$, as follows.
Let $`B_UB`$ be a set of instrument-controlling commands, thought of as strings that can be concatenated. Suppose the function $`U`$ has the form $`U(b_1b_2)=U(b_2)U(b_1)`$ for all $`b_1b_2B_U`$ (note reversal of order). Then we say the function $`U`$ respects concatenation.
###### Property 4
Quantum computation employs a subset of models in which $`U`$ respects concatenation.
Finally, the theories widely used in quantum computing assume something about timing:
###### Property 5
the unitary transformation commanded by any command $`b_U`$s takes a state $`|v`$ at one time into a a state $`|v^{}`$ at a time later by some amount $`T(b_U)`$.
###### Remark 1
To appreciate contact between instruments and equations, it is essential to see Properties 1 through 5 not as properties of laboratory instruments, but as properties that a scientist can choose to demand of models. Whether the instruments act that way is another question. There are reasons, relaxation and other forms of decoherence among them, to expect limits to the precision with which instruments can behave in accord with properties 3 through 5. All five properties are used often enough to be conventions, in the sense that a convention is a guess endorsed by a community.
### 2.1 Statistically significant differences between models
In practice, a scientist has little interest in a model chosen so that its probabilities exactly fit measured relative frequencies. Rather, the scientist wants a simpler model with some appealing structure that comes reasonably close to fitting. Quantum mechanics encourages this predilection, because on account of statistical variation in the sample mean, functions that perfectly fit outcomes on hand at one time are not apt to fit perfectly outcomes acquired subsequently. We show here that accepting statistics no way takes away from the proof that measurements and equations by themselves cannot link models to instruments.
One needs a criterion for the statistical significance of a difference between two quantum-mechanical models (or between a model and measured relative frequencies). Here we limit our attention to models $`\alpha `$ and $`\beta `$ which have a set $`B`$ of commands in common and for which the spectra of $`M_\alpha `$ and $`M_\beta `$ are the same. For a single command $`b`$, the question is whether the difference between the probability distributions $`\mathrm{Pr}_\alpha (|b)`$ and $`\mathrm{Pr}_\beta (|b)`$ is bigger than typical fluctuations expected in $`N(b)`$ trials. An answer is that two distributions are indistinguishable statistically in $`N(b)`$ trials unless
$$N(b)^{1/2}d(\mathrm{Pr}_\alpha (|b),\mathrm{Pr}_\beta (|b))>1,$$
(5)
where $`d`$ is the statistical distance defined in Eq. (10) of a paper by Wooters. Furthermore, Wooters’s Eq. (12) shows for two models $`\alpha `$ and $`\beta `$ that differ only in the function $`|v`$,
$$d(\mathrm{Pr}_\alpha (|b),\mathrm{Pr}_\beta (|b))\mathrm{cos}^1|v_\alpha (b)|v_\beta (b)|.$$
(6)
To help judge the significance of the difference between two models with respect to a set B of commands common to them, a scientist who chooses some weighting of different commands can define a weighted average of $`d(\mathrm{Pr}_\alpha (|b),\mathrm{Pr}_\beta (|b))`$ over all $`bB`$. The same holds if model $`\beta `$ is replaced by relative frequencies of outcomes interpreted from measured results.
It has been proved that the set of models statistically indistinguishable from a given model can be much larger than would be the case if the “$``$” of (6) were an equality.
###### Proposition 2.1
For any set of outcomes, two models $`\alpha `$ and $`\beta `$ of the form $`(|v,\mathrm{𝟏},M)_B`$ can perfectly fit the relative frequencies of the outcomes (Proposition 2.1) and yet be mutually orthogonal in the sense that $`v_\alpha |v_\beta =0`$
Wooters extended the definition of statistical difference to unit vectors representing quantum states. While for any two unit vectors, there exist measurement operators that maximize the statistical distance between them, for any such operator there exist other vectors, mutually orthogonal, that have zero statistical distance relative to this operator. For this reason, among others, statistics still leaves the scientist needing something beyond calculation and measurement to determine a model, for the set of models closer than $`ϵ`$ in weighted statistical distance to certain measured results certainly includes all the models that exactly fit the data and, without special restrictions dependent on guesses, this set includes models that are mutually orthogonal. Models close to given measured data are not necessarily close to each other in the predictions they make.
## 3 Large sample size needed to decide which model better fits measured data
By recognizing that models are distinct from instruments, it is easy to ask a question of the form: suppose model $`\alpha `$ assumed to describe instruments, is wrong and another model, call it $`\beta `$, describes the instruments better with respect to outcomes in some situation. Even if instruments behave in accord with a model $`\alpha `$, it can be a lot of work to show this, and without showing it, one does not know it. While there are many models having properties 1 and 2 that fit any given outcomes of instruments exactly, making the the models indistinguishable with respect to fit, there are many other models that do not fit. Given a model $`\alpha `$ that fits and a model $`\gamma `$ that does not, it was shown that models $`\alpha `$ and $`\gamma `$ are statistically distinguishable with respect to their fit to data measured for a command $`b`$ only if the statistical distance between $`\alpha `$ and $`\gamma `$, which we call $`ϵ`$, is big enough in relation to sample size $`N(b)`$; distinguishability requires
$$N(b)ϵ^2.$$
(7)
This raises the question of the size of the statistical distance $`ϵ`$. Very small values of $`ϵ`$ are required to implement quantum searching, implying large sample size, as follows. The task of searching is often expressed as the task of finding the value (assumed to exist and to be unique) for which a binary function $`f`$ of $`n`$ bits is 1. The function $`f`$ is expressed by an oracle, represented by a unitary transformation $`U_f`$, diagonal in the computational basis of dimension $`N=2^n`$, with 1 for each value except for $``$1 in the place for the argument for which $`f`$ takes the value 1. For the example of the function $`f_0`$ such that $`f(0)=1`$ and $`f(x)=0`$ for $`x0`$, $`U_f=U_0=\mathrm{𝟏}2|00|`$, The algorithm for quantum search can be viewed as implementing a model that calls for three steps, (the second of which is repeated many times):
1. prepare a state
$$|w\stackrel{\mathrm{def}}{=}N^{1/2}\underset{j=0}{\overset{N1}{}}|j$$
(8)
2. on the order $`N^{1/2}`$ times, apply $`U_f`$ followed by $`U_w=\mathrm{𝟏}2|ww|`$; and
3. make a nondegenerate measurement diagonal in the computational basis.
This works because $`U_w`$ reflects about the hyperplane perpendicular to $`|w`$, $`U_f`$ reflects about the hyperplane perpendicular to the basis vector corresponding to the special value (which is $`|0`$ for the case $`f=f_0`$), and the product of the two reflections acts as a rotation moving the starting vector $`|w`$ toward the special vector for $`f`$.
Now we show explicitly how a small error in a unitary transformation can cause a failure. Suppose that the instruments behave not in accord with the model $`\alpha `$ described but in accord with a model $`\beta `$ that differs from $`\alpha `$ in that in place of $`U_w`$, the transformation is $`U_{\stackrel{~}{w}}`$ which lacks the component $`|0`$:
$$|\stackrel{~}{w}\stackrel{\mathrm{def}}{=}(N1)^{1/2}\underset{j=1}{\overset{N1}{}}|j.$$
(9)
It is easy to check that the angle $`\theta `$ between $`|w`$ and $`|\stackrel{~}{w}`$ defined by $`\theta =\mathrm{cos}^1|w|\stackrel{~}{w}|`$ is given by
$$\theta =cos^1[(11/N)^{1/2}]N^{1/2}=2^{n/2}.$$
(10)
Correspondingly, the error between the desired unitary transformation $`U_w`$ and the transformation $`U_{\stackrel{~}{w}}`$ that describes what the instruments do is readily calculated to be
$$ϵ=U_wU_{\stackrel{~}{w}}=2|\mathrm{sin}\theta |2N^{1/2}=2^{1n/2},$$
(11)
which is exponentially small in the number $`n`$ of bits. Though exponentially small, the error $`ϵ`$ completely destroys the quantum computation for the case of $`f_0`$, because (by construction) $`0|\stackrel{~}{w}=0`$, with the result that
$$(U_{\stackrel{~}{w}}U_0)^2=\mathrm{𝟏},$$
(12)
so the repeated applications of $`U_{\stackrel{~}{w}}U_0`$ (in step 2, above) accomplish nothing, and no magnification of the desired component $`|0`$ is achieved, so with an error in $`U_w`$ as in Eq. (11), the outcome has negligible probability of corresponding to the right answer.
### 3.1 Lower bound on sample size required to verify performance
Exploring contact between instruments and equations requires not just the recognition of multiple models for given instruments, but also modeling the commands sent to instruments. To verify that a command $`b_U`$ generates statistics corresponding to a desired unitary matrix $`U_1`$ requires trying the command out in conjunction with related commands $`b_{v,j}`$ to prepare state vectors ($`j`$ ranges over a sufficiently large set). The number of these vectors must be larger than the dimension of the vector subspace relevant to problems in which $`U_1`$ is used. This means for $`n`$-bit quantum computing the number of vectors is greater than $`N=2^n`$. For each vector, it follows from Eqs. (7) and (11) that the sample size required to show that a command $`b_U`$ produces statistics consistent with any model to within the precision required for quantum searching is large indeed, namely $`N(b_U)>ϵ^2>2^{n2}`$, whence it follows that the number of trials (the sample size) needed to verify experimentally that command $`b_U`$ accords within $`ϵ`$ with a model that says it generates $`U_1`$ is greater than $`2^nN(b_U)`$, which in turn is greater than $`2^{2n2}`$. With less than this amount of testing, the likelihood that quantum computing instruments will perform to the required precision is essentially zero.
## 4 High precision required for timing of CPC commands
While within some margin, the physics of a classical computer is insulated from its manipulating of symbols, that physics is not insulated from instruments commanded by a CPC; indeed the precision of timing of the CPC matters critically to the successful functioning of instruments that it commands. Per property 5 of section 2, each unitary transformation takes a state at an earlier time to a state at a later time. Thus a unitary transformation happens not all at once, but over a time duration that depends on how the instruments implement the transformation. A written command $`b_U`$ acts as a musical score. Like sight reading at a piano, executing a program containing the command $`b_U`$ requires converting the character string $`b_U`$—the score—into precisely timed actions—like the striking of keys at the piano. In this analogy, the piano keys, so to speak, include the output buffers that control the amplitude, phase, frequency, and polarization of lasers of an ion-trap quantum computer or of radio-frequency transmitters for a nuclear-magnetic-resonance (NMR) quantum computer.
For this reason, executing a command $`b_U`$ requires parsing it into pieces (signals) and sending each signal at its proper time, the specification of which is contained in the string $`b_U`$. Either the CPC that executes a program in which $`b_U`$ is written parses the command into signals and transmits each signal at its appointed time, or the instruments receiving the command $`b_U`$, unparsed, contain programmable counters operating in conjunction with a clock that do this timed parsing. Such programmable counters themselves constitute a special-purpose CPC. So either the scientist’s CPC must execute commands by issuing an appropriately timed sequence of signals, or some other CPC attached to the instruments must do this. Either way, the capacity to execute programmed motion in step with a clock is a requirement for a CPC, distinct from and in addition to requirements to act as a Turing machine.
Implicit in models used in quantum computing is an additional property, to do with the effect of the mistiming of signals transmitted by a CPC executing a command $`b_U`$ for a unitary transformation. The effect of mistiming is to generate some (unwanted) command $`\stackrel{~}{b}_U`$ in place of $`b_U`$. To think about the effect of mistiming, one can start with the simple case of a mistimed NOT-gate for one bit. Two such gate operations concatenated result in the identity, which in the Bloch picture is a rotation by $`2\pi `$. If it takes $`T`$(NOT) seconds to perform the NOT-gate, there is an angular rotation rate (in radians/s) of $`\omega =\pi /T(\text{NOT})`$. To avoid making an error greater than $`ϵ`$ in the angle of the state vector, it is then necessary to have the error $`\mathrm{\Delta }T`$ allowable for the CPC in the timing of signals satisfy:
$$\mathrm{\Delta }T/T(\text{NOT})ϵ/\omega T(\text{NOT})=ϵ/\pi .$$
(13)
For the search algorithm, Eq. (11) implies
$$\mathrm{\Delta }T/T(\text{NOT})<2^{n/2}.$$
(14)
The best precision of timing available from hydrogen masers is something like 1 part in 10<sup>15</sup>, which fails to be adequate if $`n>30\mathrm{ln}10/\mathrm{ln}2<100`$; i.e. the precision of the best hydrogen masers is less than the precision of timing needed if $`n>99`$.
## 5 Modeling calls for guesswork
In this section we review propositions proved elsewhere concerning guesswork required to link quantum-mechanical models to instruments and contribute a new one; in addition, we describe how the set of models forms a lattice that a scientist navigates in searching for a model that fits data to within some prescribed weighted statistical distance; finally we describe guesswork needed to replace a model that fits the behavior of instruments at one level of precision with a model that fits more precisely.
### 5.1 Trying for a model that fits the instruments
Consider a scientist searching for a model $`\alpha `$ within a weighted statistical distance $`ϵ`$ of relative frequencies of outcomes obtained by use of the instruments. The scientist needs to choose a model from some large set $`S`$ of models. To illustrate what is involved, suppose $`S`$ is the set of models of section 2, constrained only to exhibit properties 1 and 2. As proved earlier, we have the following three propositions:
###### Proposition 5.1
Given any recorded counts of measured outcomes associated with any set $`B`$ of commands, the set of models satisfying properties 1 and 2 contains many unitarily inequivalent models $`(|v,U,M)_B`$, each of which fits perfectly the relative frequencies of outcomes.
###### Proposition 5.2
For any set of outcomes, two models $`\alpha `$ and $`\beta `$ of the form $`(|v,\mathrm{𝟏},M)_B`$ can perfectly fit the relative frequencies of the outcomes (Proposition 5.1) and yet be mutually orthogonal in the sense that $`v_\alpha |v_\beta =0`$
###### Proposition 5.3
For measured data to uniquely decide to within unitary equivalence which quantum-mechanical model of a set of models best fits experimental results interpreted as outcomes by a criterion of least statistical distance (or any other plausible criterion), the set of models must first be sufficiently narrowed, and this narrowing is underivable from the results and the basic properties 1 and 2 of quantum mechanics.
To these we now append:
###### Proposition 5.4
Any model $`(|v,U,M)_B`$ working in the Hilbert space $``$ can be mimicked exactly for all commands of $`B`$ by two other mutually perpendicular models working with any command set $`\stackrel{~}{B}B`$ in the Hilbert space $``$.
Proof: To generate one of the mimicking models, replace $`|v`$ by $`|\stackrel{~}{v}|w`$ while for the other use $`|v|\stackrel{~}{v}|w_{}`$, for any $`|\stackrel{~}{v}|_B=|v`$ and any $`|w,|w_{}:\stackrel{~}{B}`$ such that $`w|w_{}=0`$; for both mimicking models, replace $`U`$ by $`U\mathrm{𝟏}`$ and $`M`$ by $`M\mathrm{𝟏}`$. An easy check confirms that these models have the claimed properties of producing the same probabilities for commands of $`B`$ as does the given model, and that the two mimicking models are orthogonal to each other. $`\mathrm{}`$
All the subsets of $`S`$ constitute a lattice under set intersection and union; we call this a lattice of models. Any record of measured results interpreted as outcomes, together with a weighted statistical distance $`ϵ`$, defines a subset of $`S`$ consisting of those models that have the property being within $`ϵ`$ in distance to the relative frequencies of the outcomes. The propositions show that $`S`$ is “too big” a set, in that it always contains models that are drastically different from each other, indeed orthogonal. That means the subset contains models that, by any plausible measure, are unappealing.
Each property applicable to $`S`$ defines a subset of $`S`$; e.g. property 3 defines a subset, and property 4 a subset of that subset. Indeed, the lattice of models can be viewed as isomorphic to a lattice of properties. Our scientist hopes to guess properties to produce a subset within which the models that are close to measured results are close to each other in the predictions they make, and have properties that the scientist likes. Actually, the scientist wants to guess properties that narrow $`S`$ down to one model that appeals to the scientist and best fits the measured results.
This can require a long journey. Starting with $`S`$, the scientist guesses a list of properties, each of which defines a subset. The list of properties defines a subset $`S^{}S`$ that shares all these properties, which is just the intersection of all the subsets defined by the properties separately. Having defined this smaller set $`S^{}`$, the scientist explores the weighted distance between the models of this subset and recorded measured data, and also explores the closeness of the models that fit the data, if any, to each other. If carried out, this exploration results in one of three cases:
1. there are widely disparate models in $`S^{}`$ that all fit the data; i.e. $`S^{}`$ is too big, so the scientist needs to guess additional properties to shrink to fewer models;
2. no models of $`S^{}`$ fit well, leading to the need to drop a property from the list to get a new subset of $`S`$, which includes $`S^{}`$;
3. some models of $`S^{}`$ fit the data well and these are close to each other in their implications, leading to the use of one of these models, say model $`\alpha `$.
### 5.2 Ratcheting up precision takes fresh guesswork
Suppose the following:
1. A scientist cycles through these cases, arriving eventually at case 3 with a model $`\alpha `$, with $`ϵ`$ as the weighted statistical distance between the model and relative frequencies of outcomes obtained from the testing of instruments.
2. Using the model to define commands for the CPC to transmit to instruments results in successful quantum searching for functions of $`n1`$ bits, but the commands do not work for $`n`$ bits; per Eq. (11), requires a reduction of error by $`2^{1/2}`$.
3. The scientist needs to produce commands that more precisely generate behavior of the instruments described by quantum gates specified as unitary matrices $`U_j`$ for $`j=1,\mathrm{},G`$, and needs to achieve weighted statistical distance $`ϵ^{}=2^{1/2}ϵ`$, or less.
#### 5.2.1 Can searching supplant the need for a better model?
If a scientist has a command $`b_j`$ that generates behavior of instruments consistent to within $`ϵ`$ of a desired unitary transformation $`U_j`$, can the scientist likely find a command that is precise to within $`ϵ^{}=2^{1/2}ϵ`$ by searching, thus skipping the need to get a more precise model? That depends on the how many commands have to tried, and the number of repetitions of each command necessary to obtain a sample size big enough to discriminate between models to within a statistical distance of $`ϵ^{}`$.
To determine this number of commands that have to be tried, we introduce a notion of an $`ϵ`$-grid on a metric space, here the Lie group $`SU(N)`$ with the metric induced by the spectral norm. By an $`ϵ`$-grid on a metric space $`𝒮`$, we mean a finite subset $`𝒢_ϵ𝒮`$ such that for any $`x𝒮`$, $`y𝒢_ϵ`$ such that $`|xy|<ϵ`$. To be sure of finding a command $`\stackrel{~}{b}_j`$ that produces behavior within $`ϵ^{}`$ of $`U_j`$, it is necessary to explore all the points of an $`ϵ`$’-grid of the $`ϵ`$-ball in $`SU(N)`$ around $`U(b_j)`$. There are at least as many such points as the ratio of the volume of an $`ϵ`$ ball to the volume of an $`ϵ`$’-ball in $`SU(N)`$. That ratio is $`(ϵ/ϵ^{})^d`$, where $`d`$ is the dimension of $`SU(N)`$, so $`d=N^21=2^{2n}1`$. Thus the number of points that have to be explored to guarantee finding one that improves the precision by enough to expand quantum search by one bit is, for just one gate command, $`2^{d/2}`$, with $`d4^n`$. Starting with precision just adequate for a 4-bit search, to reach precision just adequate for a 5-bit search requires exploring about $`10^{154}`$ commands. Without counting the large number of trials needed to get a sample size big enough, blind searching is hopeless as a way to find better commands. Finding a better model to generate commands, and thus avoiding a search over all these commands, appears indispensable.
#### 5.2.2 Trying for a better model
A model $`\alpha `$ found in the beginning of a quantum-computing experiment by trying out small problems for a small sample sizes and a small set of commands $`B_{\mathrm{test}}`$ will likely be inadequate to generate the precision required to deal with larger problems, such as factorizing larger integers or searching a larger list. With blind searching ruled out, a better model, call it $`\beta `$, is necessary from which to determine more precisely the commands and their timing.
How does looking for and deciding on a better model work? Can this be done by a universally applicable program, run by the CPC that controls the quantum computing instruments? Hope for this automation lies mostly the idea that searching for a better model is adequately defined by the objective of a better fit to measured results and by a desire for simplicity in the model; the hope is perhaps spurred also by the allure of the universal test used in mathematical logic, isolated from laboratory instruments, for testing a claimed derivation, a test that can be thought of as a universally applicable program for a Turing machine.
This hope is dashed by Wittgenstein’s contribution to the foundations of both psychology and set theory, when he made unequivocally clear that words are what people say in various scenes, so that what the speaking of a word means depends on the scene in which the word is spoken. This applies to words transmitted by a CPC as commands to computer-controlled instruments. The evaluation of the words cannot be independent of the evaluation of the activity and actions of the instruments generated by those words. (Without this connection any concept of testing falls in confusion, because the same word, whether uttered by a person or transmitted as a command to an instrument by a CPC, produces actions that depend on context beyond the power of verbal or mathematical analysis, as follows from the previously cited proof of the need for guesswork to link models to instruments.) This means that, even under our bold assumption that suitable commands exist, the finding of a model with certain properties that fits measured data at one level of precision by no means assures the existence of a model having those properties and fitting the same measured data at greater precision, let alone data obtained from larger samples and a larger set of commands. Thus when larger problems are tackled, with their demand for increased precision, the scientist can expect to have to go back through the lattice of models, encountering all three cases, to find new, possibly more complex properties by which to narrow the set of models. This implies a need for a scientist operating a quantum computer to switch back and forth between modes of quantum computing and modes of refining the model and the commands it calls for.
## 6 Concluding remarks
It has been striking and surprising to find in the course of this work that contact between instruments and equations is a barely discovered discipline of physics, awaiting further investigation. The example of quantum searching shows up some of the issues of contact especially vividly. The quantum computing problem of factorizing and the Deutsch-Jozsa problem require a number of gate transformations polynomial rather than exponential in the number of bits, and so the challenges to their implementation are less extreme. Nonetheless implementing any device described quantum mechanically at progressive levels of precision calls for navigating a lattice of models, coping with ambiguity among inequivalent models, dealing with sample sizes made necessary by quantum indeterminacy, and managing via machinery that is specified classically the timing of signals at the precision demanded by quantum mechanics. Here we have provided some tools for a promising discipline. Since they are of general application, not limited to quantum computing, these tools are widely applicable to efforts that join theory and experiment in quantum physics.
## Acknowledgments
We thank Amr Fahmy for contributing to our thoughts and discussions at many points during this work. In addition to him, we are indebted to Steffen Glaser, and Raimund Marx, and Wolfgang Bermel for introducing us to the subtleties of laboratory work aimed at nuclear-magnetic-resonance quantum computers. We learned the view of quantum searching used here from David Mumford.
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# Peak effect in CeRu2: history dependence and supercooling
## I introduction
Recent theoretical and experimental studies based on high T<sub>C</sub> superconductors (HTSC) suggest the existence of at least two distinctly resolved solid phases of vortex matter which are distinguished from the high-temperature high-field vortex liquid. These two vortex solid phases are referred to as low-field quasi-ordered solid or Bragg-glass and high-field disordered solid or vortex-glass (see Ref.8). The important question now is, what is the the order of the thermodynamic phase transitions (if any) between the various vortex phases? The Bragg-glass has long range order and it is expected to melt to vortex liquid at high temperature through a first order transition. Experimentally, the indication of a first order transition usually comes via a hysteretic behaviour of various properties, not necessarily thermodynamic ones. In HTSC samples also initial suggestions of a first order melting transition came via distinct hysteresis observed in transport property measurements . The confirmatory tests of first order transition ofcourse involve the detection of discontinuous change in thermodynamic observables and the estimation of latent heat, and this has subsequently been achieved for vortex melting in HTSC materials through magnetization and calorimetric measurements . There also exists a less rigorous class of experimental tests which involves the study of phase co-existence and supercooling across a first order transtion. This kind of experiment has also come out to be pretty informative for the first order melting transition of the Bragg-glass in Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> (BSCCO) (Ref.15).
With the establishment of the first order nature of the Bragg-glass to vortex liquid transition line, the focus in the recent years has shifted to the Bragg-glass to vortex-glass transition. In various HTSC materials peak-effect(PE) or fish-tail is used to track this field induced transition from Bragg-glass to vortex-glass, and this transition is obsereved to be a sharp transtion . However, the exact nature of this transition –whether it is a continuous or a first order transition– is not established yet. Very recent magneto-optics studies on single crystal samples of BSCCO claim the presence of phase-coexistence and supercooling across the Bragg-glass-vortex-glass phase transition. These in turn suggest the possibility of a first order phase transition. Newer theoretical developments have also taken place in the field of vortex-matter physics during last few years to understand the various field/disorder induced phenomena. This activity in vortex matter physics can even have deep correlations with a more general area namely disorder/pressure induced melting/amorphization in real solids .
Very recently interesting history dependence of PE has been reported for naturally untwinned and detwinned single crystal samples of YBa<sub>2</sub>Cu<sub>3</sub>O<sub>y</sub>(YBCO) with 6.908$``$y$``$6.999 . We have earlier observed very similar history dependence of PE in the low temperature (T$`{}_{C}{}^{}`$6.1K) C15-Laves phase superconductor CeRu<sub>2</sub> . We believe that the $`\mathrm{"}\mathrm{𝑚𝑖𝑛𝑜𝑟}`$ $`\mathrm{ℎ𝑦𝑠𝑡𝑒𝑟𝑒𝑠𝑖𝑠}`$ $`\mathrm{𝑙𝑜𝑜𝑝}\mathrm{"}`$ technique used in our early study of history dependence of PE in CeRu<sub>2</sub> and also in the similar recent studies on YBCO belongs to that class of experiments which can investigate the phase coexistence and supercooling across a first order phase transition.
Interesting Fermi surface topology and enhanced pauli paramagnetism of CeRu<sub>2</sub> have given rise to various interesting possibilities, starting from an exotic non s-wave superconducting ground state to a field induced change in the microscopic superconducting order parameter associated with the onset of a Fulde-Ferrel-Larkin-Ovchinnikov (FFLO) state . Onset of FFLO state can cause softening of the vortex lattice and in turn enhanced pinning and PE . It is quite clear that if a field induced first order transtion from a BCS state to FFLO state actually takes place in any system, a PE with definite history effects (typically associated with a first order transition) will be observed, and this has been motivating us, possibly as a red herring, in our study of vortex matter in CeRu<sub>2</sub>. The question is can CeRu<sub>2</sub> actually sustain a FFLO state ? Our own macroscopic magnetization and transport-property measurements cannot provide any proof for (or against) the existance of FFLO state, and at this moment the global debate on this issue remain unsettled.
In the light of various recent developments described above regarding vortex matter, we believe that results obtained in CeRu<sub>2</sub> should be re-examined. In this paper we shall
1. present newer results on the history dependence of PE in a single crystal sample of CeRu<sub>2</sub>. Recently we have extended the classical theory for supercooling across first order phase transitions to the case when both field and temperature are control variables and have shown how the observable region of metastability depends on the path followed in the space of these two variables . We show that our present experimental results are in consonance with these theoretical predictions, and hence reinforce the idea of a first order transition from one kind of vortex solid to another.
2. compare our results on CeRu<sub>2</sub> with the recently reported history dependence of PE in single crystal samples of YBCO and point out the similarity and difference between these two disparate class of materials. The relevance of the results on the phase coexistence and supercooling in BSCCO will also be discussed.
3. discuss a possible origin of a first order transition between two kinds of vortex solid in CeRu<sub>2</sub> within the realm of the recent theoretical developments .
Preliminary results of the present work have been presented in the recently held LT-22 conference in Helsinki .
## II experimental
In contrast to the polycrystalline samples of CeRu<sub>2</sub> used in our earlier studies showing history effects associated with PE, in the present study we use a single crystal sample of CeRu<sub>2</sub> (T$`{}_{C}{}^{}`$6.1K). The details of the preparation and characterization of this sample can be found in Ref.41.
Magnetization measurements were performed using a commercial SQUID magnetometer (Quantum Design MPMS5). We have used a 2 cm scan length in the ’fixed-range’ mode to minimize the sample movement in the inhomogeneous field of the superconducting magnet. In the ’auto-range’ mode the sample goes through multiple movements while the system software searches for the most sensitive gain useful for the signal level detected. We carried out a separate preliminary run using the auto-range mode to identify the appropriate gain for the given experimental conditions and then performed a final run in the ’fixed-range’ mode. In the case of 2cm scan-length, the field inhomogeneity in an applied field of 20 kOe is $``$2 Oe. We have concluded earlier that in an isothermal field scan, as long as the field for full penetration at a particular field value is substantially larger than the field inhomogeneity during the sample measurement, the error in the results in the particular type of measurements reported here will be negligible. Inspite of all these cross-checks, in the light of general doubts concerning the measurement procedure using commercial SQUID magnetometers, it has become important to reproduce the observed history effects using other techniques which minimize the sample movement. In a recent work we have shown the existence of the history dependence of PE in the isothermal field variation of magnetization in polycrystalline samples of CeRu<sub>2</sub> using an axial-VSM (Oxford Instruments). In the present work we have used a transverse-VSM (Oxford Instruments) to get supporting results. In this transverse VSM the sample is placed in the centre of the pick-up coil and the superconducting magnet assembly, where the magnetic field inhomogeneity is $``$0.01% over 1 cm diameter spherical volume (DSV).
In contrast to the axial-VSM, the direction of sample vibration in the transverse-VSM is perpendicular to the applied field. No significant change is observed in the results by varying the sample vibration amplitude between 0.5 and 1.5mm; this rules out any distinct role of magnetic field inhomogeneity. In any case the field inhomogeneity encountered here is obviously much smaller than that encountered even in the 2cm-scan of SQUID magnetometer.
In magnetization hysteresis measurements, we draw an isothermal magnetization (M) versus field (H) curve by cycling H between $`\pm H_{C2}`$ at various temperatures (T) below T<sub>C</sub>. These T of interest are reached by cooling through T<sub>C</sub> in absence of any applied field H i.e. in zero-field- cooled (ZFC) mode. From various H-points on this isothermal M-H curve (or envelope curve), a $`\mathrm{𝑚𝑖𝑛𝑜𝑟}`$ $`\mathrm{ℎ𝑦𝑠𝑡𝑒𝑟𝑒𝑠𝑖𝑠}`$ $`\mathrm{𝑙𝑜𝑜𝑝}`$ (MHL) can be be drawn either by decreasing H from the ascending field branch of the M-H curve (or lower envelope curve) or by increasing H from the descending field branch (or upper envelope curve). We designate these MHLs as (MHL)<sub>ZFC</sub>. One can also initiate MHLs from various H-points obtained by field cooling (FC) through T<sub>C</sub>. In such a case the M-value at the starting H-point normally lies inbetween the upper and lower envelope curve (see Ref.27 and references cited therein). An MHL can be initiated from such FC H-points either by increasing or by decreasing H. We designate these MHLs as (MHL)<sub>FC</sub>. Within the realm of a single-component Bean’s crtical state model, all these MHLs (both ZFC and FC) are expected to saturate by reaching the enevelope curve . Such a behaviour has actually been observed experimentally in various type-II superconductors . The same is also observed in all H-regimes, except a finite field regime encompassing at least a part of the PE regime, in various samples of CeRu<sub>2</sub> including the present single crystal sample. In this latter regime the MHLs do not behave in accordance with the critical state models, and we have used this anomalous behaviour of MHLs to study the interesting history dependence of PE in CeRu<sub>2</sub> . Very similar technique involving the (MHL)<sub>ZFC</sub> has very recently been used to study the history dependence of PE in single crystals of YBCO .
Apart from tracking the history dependence of PE the study of MHLs can provide with at least two more useful information in our present kind of study. First, it provides a way of estimating the field for full penetration at a particular H of interest . The field inhomogeneity $`\delta `$H of the magnet causes the sample to effectively follow an MHL during measurements with a SQUID magnetometer . Since $`\delta `$H rises with scan length, MHLs allow one to estimate the error $`\delta `$M in measurements made with various scan lengths. One can then choose a scan length such that $`\delta `$M is much smaller than the magnetization hysteresis, or, equivalently $`\delta `$H is much smaller than the field for full penetration. Second, a fair amount of information regarding the influence of surface barrier can be obtained from the nature of the approach (linear or non-linear) of these MHLs to the envelope curve . These two information are useful even where the M-H curves are in accordance with the critical state models . Thus the technique of MHLs has universal applicability in determining (i) the effect of field inhomogeneity of the magnet on a measurement (ii) the importance of surface barrier in the magnetization study.
While presenting newer results, the present study on a single crystal sample of CeRu<sub>2</sub> using both SQUID magnetometer and VSM, provides also a cross-check on our earlier results obtained on polycrystalline samples mainly using a SQUID magnetometer. This will put the experimental situation regarding the PE in CeRu<sub>2</sub> on a more firm ground.
## III Results and discussion
Fig. 1 shows M-H plots of the CeRu<sub>2</sub> single crystal at T=4.5K, obtained using both SQUID magnetometer and VSM. These M-H curves are obtained by isothermally cycling H between $`\pm `$25kOe. Since H$`{}_{C2}{}^{}(4.5K)`$21.5kOe is less than 25 kOe, these provide the envelope hysteresis curve within which all the MHLs should be contained. This envelope M-H curve shows two distinct irreversible regimes separated by an almost reversible regime (see Fig. 1). While this intermediate regime appears quite reversible in the SQUID measurement (see inset of Fig. 1(a)), perceptible irreversibility is observed in the VSM measurement (see inset of Fig. 1(b)). This field-induced enhanced magnetization-irreversibility in the high field regime is the so called peak-effect (PE) and this is the subject of main interest in the present work. We note that in Fig.1 the onset field H$`{}_{a}{}^{}{}_{}{}^{}`$ of the PE in the ascending field cycle is distinctly different from the field H$`{}_{d}{}^{}{}_{}{}^{}`$ where the PE ceases to exist in the descending field cycle. Note that H$`{}_{a}{}^{}{}_{}{}^{}`$ and H$`{}_{d}{}^{}{}_{}{}^{}`$ obtained from the SQUID and VSM measurements are a bit different. We attribute this quantitative difference to (i)the difference in the magnetic field inhomogeneity encountered during SQUID and VSM measurements, (ii)possible minor difference in temperature in two different machines. Also the magnitude of measured M is smaller in the VSM measurement which we attribute to different sample orientation and the associated demagnetization factor. However, it should be noted that the actual measurement time involved in a VSM measurement, where the field is swept with a constant rate (100 Oe/s in the present case) is faster than in a SQUID measurement where the field is stabilized with a pause time (10 sec in the present case) before each measurement. In systems with finite magnetization relaxation, faster measurement with VSM would yield larger magnetization value. This is actually observed in the field regime just below the PE where magnetization hysteresis obtained with VSM measurements is perceptibly higher \[see Fig.1(a) and Fig.1(b)\].
We shall now present results obtained in the form of MHLs measured at closely spaced field intervals, after preparing the vortex lattice within the following experimental protocols :
1. Zero field cool (ZFC) the sample to the temperature of measurement, switch on a field less than -H<sub>C2</sub>, and then increase the field isothermally to reach various points on the lower envelope curve. (MHL)<sub>ZFC</sub>’s are drawn by reducing the field isothermally.
2. After the above step, increase the field to a value greater than H<sub>C2</sub> and then reduce the field to reach the various points on the upper envelope curve, while maintaining the isothermal condition. (MHL)<sub>ZFC</sub>’s are drawn by increasing the field isothermally.
3. Field cool (FC) the sample in various (positive) fields from a temperature substantially above T<sub>C</sub>. After stabilizing the temperature of interest the (MHL)<sub>FC</sub>’s can be drawn both by increasing and decreasing the field isothermally.
The MHLs at the onset of the PE regime obtained within the above experimental protocols do not conform with the critical-state models;they do not show the expected merger with the envelope curve (see Fig.2). While the (MHL)<sub>ZFC</sub>’s initiated from the lower envelope curve at H=18.25, 18.75 and 19 kOe saturate without touching the upper envelope curve (see Fig. 2(a)), the (MHL)<sub>ZFC</sub>’s initiated from the upper envelope curve at H=18.4 kOe overshoots the lower envelope curve before reaching saturation (see Fig. 2(b)). The (MHL)<sub>FC</sub>’s obtained following the FC path also overshoot the envelope curve (see Fig. 2(c)). For the sake of clarity and conciseness we show only few representative MHLs.
We have earlier reported anomalous behaviour of (MHL)<sub>ZFC</sub>’s obtained within the protocol no.1 in polycrystalline samples of CeRu<sub>2</sub>. Tenya et al have reported anomalous (MHL)<sub>ZFC</sub>s under protocol no.2. Ravikumar et al (Ref.50) and ourselves (Ref.29(b)) reported anomalous (MHL)<sub>FC</sub>’s under protocol no.3 in single crystal and polycrystalline samples of CeRu<sub>2</sub> respectively. Actually the anomalous character of the (MHL)<sub>ZFC</sub> initiated from the upper envelope curve was visible in a relatively less prominent manner in our earlier studies of polycrystalline samples as well (see Fig.3 of Ref.29 (a)). However, from our standard cross-checks we were not sure whether the observed result was beyond error bar, hence did not emphasise much on that. (On the other hand we were aware of the metastable character of those (MHL)<sub>ZFC</sub>’s drawn from the upper envelope curve, since they readily shattered on field cycling (see Fig. 3 of Ref. 29(a).) In the present work for the first time all the three different kinds of (MHL)s are shown on the same sample of CeRu<sub>2</sub>.
We shall now reproduce all these anomalous aspects of MHLs in the vicinity of PE using a transverse-VSM. Fig.3 shows the anomalous nature of the two (MHL)<sub>ZFC</sub>’s drawn from the lower envelope curve at H=19.25 kOe and 20 kOe, and two (MHL)<sub>FC</sub>’s drawn by reducing the field at H=17.25 and 19 kOe. In this transverse-VSM it is relatively difficult to stabilize the temperature for T $`<`$10K and actual temperature can vary by $`\pm `$0.025K between different runs and some time even during a single complete run. This leads to a slight scatter in the data, but from multiple runs of the same experiment we ensure that the observed anomalous features are certainly well beyond the error bars. For the two kinds of MHLs described above we need to obtain the envelope M-H curve in a separate experimental cycle. However, the anomalous nature of the (MHL)<sub>ZFC</sub>’s drawn from the upper envelope curve can be highlighted from a single experimental cycle. For this purpose we increase the field isothermally to field values well above H<sub>C2</sub> (thus drawing the lower envelope curve) and then reach the various field points of interest on the upper envelope curve by isothermal reduction of the field. From such field points we draw (MHL)<sub>ZFC</sub>’s by reversing the field sweep direction. In Fig.4 we show the MHLs initiated from the upper envelope curve at H=18.75 and 19.25 kOe and they distinctly overshoot the lower envelope curve before saturation. These representative MHLS provide support to the results obtained earlier with the SQUID magnetometer. Various envelope curves and MHLs presented in Fig.3 and 4 are obtained with field sweep rate 100 Oe/sec. We have also checked the qualitative aspects of our results by varying the field sweep rate between 20 and 200 Oe/sec.
Interesting history effects associated with the PE have been reported recently in detwinned and naturally untwinned single crystals of YBCO . These history effects in YBCO are exactly in the same form of anomalous isothermal (MHL)<sub>ZFC</sub>’s drawn from the upper and lower envelope M-H curve of CeRu<sub>2</sub> as shown in Fig. 2(a), 2(b), 3(a) and 4. In these reports, however, there is no mention of any measurement involving MHLs in a field cooled vortex state.
We have earlier associated these anomalous features in polycrystalline samples of CeRu<sub>2</sub> with a field induced first order transition to a superconducting mixed state with enhanced pinning properties . While the multivaluedness in the saturation of the (MHL)<sub>ZFC</sub> in the ascending field cycle was attributed to the nucleation and growth of the high field phase, the overshooting of the envelope M-H curve by the (MHL)<sub>ZFC</sub>’s in the descending field cycle and by the (MHL)<sub>FC</sub>, was thought to be a result of supercooling of the high field phase across the proposed first order transition . The present results on a good quality single crystal of CeRu<sub>2</sub> reinforce this picture.
We shall now provide newer evidence to support the conjecture of a first order phase transition in CeRu<sub>2</sub>. Extending the classical theory for supercooling across first order phase transitions to the case when both field and temperature are control variables, we have shown theoretically that the observable region of metastability depends on the path followed in this space of two variables, with variation of field providing a source of fluctuations in the supercooled state . We have predicted that a disordered phase can be supercooled upto the limit of metastability T(H) only if T is lowered in constant H. If the T<sub>C</sub>(H) line is crossed by lowering H at constant T, then supercooling will terminate at T<sub>0</sub>(H) which lies above the T(H) line . If T<sub>C</sub> falls with rising field, then (T<sub>0</sub>(H)-T(H)) rises with rising field . As narrated below both these predictions are experimentally found to be true in CeRu<sub>2</sub>. We have found that the anomalous features in (MHL)<sub>FC</sub> in CeRu<sub>2</sub> continue to exist in a (H,T) regime which is well below the PE regime. In this regime the (MHL)<sub>ZFC</sub>’s show normal behaviour as expected within the critical state models. This clearly suggests that the FC vortex state can be supercooled more than the isothermal ZFC state. Collating the H values where the various MHLs first show the anomalous behaviour at various T in our SQUID magnetometer measurements, in Fig.5 we present a (H,T) phase diagram. The distinct identity of the H$`{}_{a}{}^{}{}_{}{}^{}`$(T) line (which indicates the onset of the PE regime in the isothermal ascending field cycle) and the H$`{}_{d}{}^{}{}_{}{}^{}`$(T) line (at which the PE regime terminates in the isothermal descending field cycle) was infact earlier taken as an indication of a first order transition . H$`{}_{d}{}^{}{}_{}{}^{}`$(T) and H$`{}_{FC}{}^{}{}_{}{}^{}`$(T) lines are akin to the T<sub>0</sub>(H) and T$`{}_{}{}^{}(H)`$ lines respectively in our theoretical study . Experimentally H$`{}_{FC}{}^{}{}_{}{}^{}`$ at a particular temperature is defined as the H value down to which the anomalous behaviour in the (MHL)<sub>FC</sub>’s is observed. Similarly below H$`{}_{d}{}^{}{}_{}{}^{}`$(T) (defined earlier) the (MHL)<sub>ZFC</sub>’s drawn from the upper envelope curve show normal behaviour namely, they merge with the lower envelope curve without any overshooting. As shown in Fig.5 the H$`{}_{FC}{}^{}{}_{}{}^{}`$(T) line is lying distinctly below the H$`{}_{d}{}^{}{}_{}{}^{}`$(T) line and (H$`{}_{d}{}^{}{}_{}{}^{}`$(T) - H$`{}_{FC}{}^{}{}_{}{}^{}`$(T)) increases with the increase in H. Although anomalous behaviour of (MHL)<sub>ZFC</sub>’s and (MHL)<sub>FC</sub>’s in CeRu<sub>2</sub> has been highlighted by various groups during last three years, the path dependence of this anomalous behaviour in (H,T) space is definitely new and has not been reported so far. These new results, which are in accordance with our theoretical prediction , will provide further support for the existence of a first order transition in the vortex matter phase diagram of CeRu<sub>2</sub>.
The observed history effects in YBCO are interpreted in terms of a transition from a low field elastic vortex lattice to a high field plastic vortex lattice . We have previously discussed this mechanism in the context of PE in polycrystalline samples of CeRu<sub>2</sub>, and argued that a first order transition probably has an edge over this mechanism in explaining the history dependent phenomena associated with PE (see Ref.29(b)). Since then our view point is reinforced by the transport study of the PE on a single crystal sample of CeRu<sub>2</sub> . We have shown that the FC state at the vicinity of the PE regime is metastable in nature and it can be shattered easily with a small field cycling (see Ref. 29(b) and 35). This observation can be relatively easily explained in terms of the metastable nature of a supercooled state and its sensitivity to any environmental fluctuation . In a very recent study of magnetization in untwinned single crystal of YBCO using micro Hall probe, history effects in FC measurements have now been reported . Relatively complicated (H,T) phase diagram of YBCO needs a contrived path to ensure proper FC measurements and this probably may be the reason why FC measurements were not reported in earlier studies . In contrast to the suggestion that the observed history effects are properties of the high field plastic vortex lattice , it is asserted that both the low field and the high field vortex lattice are robust in nature, and the history effects and metastability are associated with the transition regime from the low field to high field phase. This recent observation makes the possibility of a first order transition in YBCO much stronger, and hence the similarity with CeRu<sub>2</sub>.
Contrary to our earlier suggestion based on bulk magnetization study , a very recent magneto optical study in good single crystals of BSCCO has claimed that the solid-solid transition in the vortex matter is indeed a first order transtion accompanied by supercooling . If this claim turns out to be true, it is possible that the supercooled (H,T) regime in BSCCO is quite narrow and/or fragile. So in our magnetization measurement we have either missed the supercooled regime and/or the fluctuation induced during the measurement procedure might have shattered the supercooled phase.
At this juncture we must point out some important difference in the history effects associated with PE in CeRu<sub>2</sub> and YBCO. While the history effects associated with PE were observed only in naturally untwinned and detwinned single crystals of YBCO and vanish quite readily with the change in oxygen stoichiometry , the same effects are quite robust in CeRu<sub>2</sub> and are observed with all the characteristic features in good quality single crystal, polycrystal, off-stoichiometric polycrystal and Nd-doped polycrystal samples of CeRu<sub>2</sub> . These results suggest that, in contrast with CeRu<sub>2</sub> the characteristic features associated with the vortex solid-solid transition in YBCO are quite sensitive to the defect. Also the negative dynamic creep at the onset of PE regime of YBCO (as reported recently (Ref.25)), is not observed with identical measurements on CeRu<sub>2</sub> . On the other hand, while the PE ceases to exist above a temperature T$`{}_{}{}^{}`$0.92 T<sub>C</sub> in all kinds of samples of CeRu<sub>2</sub>, there is no report (to our knowledge) so far of a vortex-solid to vortex-liquid melting transition in CeRu<sub>2</sub>.
We shall now discuss some recent developments in the field of vortex matter physics since they appear relevant to our experimental results in CeRu<sub>2</sub>. At low tempeartures and at high fields disorder dominates in vortex matter and topological defects proliferate, resulting in highly disordered solid . We shall now explore the possibility of a disorder induced first order transition in the vortex matter which can lead to a softening of the vortex lattice. PE will be used as an observable effect of such lattice softening. (This is in contrast with the field induced FFLO state where the lattice softening is due to a micorscopic change in the superconducting order parameter .) The point defects in the underlying crystalline lattice can cause transverse wandering of the vortex lines and this frozen-in wanderings can destroy the long-range order of the vortex lattice (Ref.4). This is analogous to the action of the thermal noise and can generate topological defects in the vortex lattice. The role of topological defects in the superconducting mixed state of type-II superconductors has been a subject of interest over the years and has come under closer scrutiny recently in the context of vortex-solid . Frey, Nelson and Fisher suggested that in the low temperature-high field regime of vortex matter, topological defects in the form of vacancies and interstitials can start proliferating leading to an intermediate supersolid state. The vortex supersolid is characterized by the coexistence of crystalline order and a finite equilibrium density of vacancy and interstitial defects . The vortex supersolid then can transform continuously into a vortex liquid state . Vortex liquid state is characterized by unbound dislocation loops as well as finite density of vacancy and interstitials . The exact nature of the transition from defect free vortex-solid to supersolid transition is not quite clear and possibility exist for both a continuous and a first order transition . Caruzzo and Yu have also considered the possibility of a first order transition to a supersoftened solid induced by interstitial and vacancy line defects in vortex lattice. Although Caruzzo and Yu mainly discuss the cases of phase transitions as a function of temperature, their theoretical approach actually is based on earlier works on disorder induced softening and first order transition in real solids . A supersoftened solid can be quite relevant in our present discussion on the field/disorder induced transitions in vortex soilds of various superconductors.
It will not be totally out of place to mention here that defect induced melting and solid state amorphization is a distinct possibility in real solids and has remained a subject of continued interest . A first order transition in such cases can be characterized by a discontinuous increase in point defects.
## IV Conclusion
Based on our results concerning history dependence of PE, we suggest the existence of a first order phase transition from one kind of vortex solid to another in the vortex matter phase diagram of CeRu<sub>2</sub>. The high-field phase can be supercooled by reducing either of the two control variables viz H and T. The extent of supercooling observed depends on the path followed in this space of two variables. In the light of very recent observation of history dependence of PE in detwinned sample of YBCO both in the ZFC and FC measurements , it will be now of interest to check experimentally for a similar (T,H) path dependence and possibility of a first order solid-solid phase transition in the vortex phase diagram of YBCO.
## V acknowledgement
We acknowledge Mohammed, S. Hebert, G. Perkins, L. F. Cohen and A. D. Caplin for various help in the experiments involving the transverse-VSM and many useful discussion. We also acknowledge Y. Radzyner, D. Giller, A. Shaulov and Y. Yeshurun for useful discussion. We thank Dr. A. D. Huxley for providing us with the single crystal sample of CeRu<sub>2</sub> used in the present study.
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# Contents
## 1 Introduction
Two-dimensional quantum field theories formulated on a surface with boundaries are important arena in various fields of study. Among others they provide the starting point for open string theories, in particular for the study of D-branes. D-branes have proven to be indispensable elements in string theories. The interplay between the target space properties of D-branes, as sources for RR fields , and how they couple to string worldsheet has been very important. In most of the recent studies of D-branes from the worldsheet point of view, boundary conformal field theories, namely conformal field theories with conformally invariant boundary conditions, have been the focus of attention. However, we believe that one can learn more, especially on the off-shell properties of string theories, by studying quantum field theories on a worldsheet with boundaries where the bulk theories and/or the boundary conditions are not necessarily conformally invariant. We also expect that such a study may open a way to reveal the geometric principle of string theories.
One of the main aims of this work is to initiate the study of D-brane associated with quantum field theories with and without conformal invariance, employing supersymmetry as the basic constraint. In particular we study D-branes in supersymmetric sigma models on Kähler manifolds and their mirror description in terms of D-branes in Landau-Ginzburg theories. We will consider both the conformal case (where target space has vanishing $`c_1`$ and is a CY manifold) as well as asymptotically free theories (where $`c_10`$) which have mass gaps in many cases. We will mainly consider D-branes corresponding to holomorphic cycles on the Kähler manifold which are mirror, as we will discuss, to Lagrangian submanifolds on the Landau-Ginzburg side.
Along the way we find some interesting similarities and differences between various aspects of D-branes for the massive sigma models and the conformal one. In particular we see how Brane creation also occurs for massive theories as we change the parameters in the sigma model. We also define a notion of intersection between two Lagrangian D-branes in the massive theory which is a refined version of the classical intersection of the cycles in the Calabi-Yau realization of it (in particular the inner product in the massive case is neither symmetric nor anti-symmetric). Also we apply the machinery of D-branes that we develop for Landau-Ginzburg theories to the LG realization of $`𝒩=2`$ minimal models. In this way we find a purely geometric realization of Verlinde algebra for bosonic $`SU(2)`$ WZW model at level $`k`$, as well as the modular transformation matrix. Also we are able to shed light on an old observation of Kontsevich connecting “helices of exceptional bundles” on Fano varieties with soliton numbers of certain Landau-Ginzburg theories, as a consequence of mirror symmetry acting on D-branes.
The organization of this paper is as follow. In section 2 we review aspects of LG solitons in $`𝒩=2`$ theories . In section 3 we discuss D-branes for supersymmetric sigma models and LG theories. In this section we will consider both holomorphic and Lagrangian D-branes. For the most of this paper we will mainly concentrate on holomorphic D-branes (“B-type”) in the context of supersymmetric sigma model and Lagrangian D-branes (“A-type”) in the context of LG models. We define and study boundary states corresponding to such D-branes following the study in string theory and conformal field theory . In section 4 we discuss the phenomenon of D-brane creation of massive LG theories, and show how these results give a reinterpretation of the connection between R-charges of chiral fields at the ultra-violet fixed point and the soliton numbers of its massive deformation discovered in . In section 5 we apply these results to the study of $`𝒩=2`$ minimal models and show how aspects of conformal theory, including certain properties of Cardy states, and its relations with Verlinde algebra, as well as its overlap with Ishibashi states in terms of modular transformation matrix, can be derived in a purely geometric way. In section 6 we derive, using the results of the mirror of certain D-branes on Fano varieties in terms of D-branes in the mirror LG models. In section 7 we apply the study of D-branes to the LG mirrors of Fano varieties and uncover beautiful mirror interpretation for helices of exceptional bundles on Fano varieties in terms of D-branes of the LG mirror. In section 8 we discuss connecting the LG mirror for the case of non-compact geometries in Calabi-Yau (such as $`\text{I}\text{P}^2`$ inside a CY threefold) discussed in to a local non-compact geometric mirror as was used in . Moreover we show how this is related in the case of local non-compact threefolds to the probe description in F-theory and its BPS states.
While completing this work, a paper , which has some overlap with our discussions in sections 3 and 5, appeared.
## 2 BPS Solitons in $`𝒩`$=2 Landau Ginzburg Theories
In this section, we review some basic facts on Landau-Ginzburg models, especially on the spectrum of BPS solitons and the relation to Picard-Lefshetz theory of vanishing cycles. The action for a Landau Ginzburg model of $`n`$ chiral superfields $`\mathrm{\Phi }_i`$ ($`i=1,\mathrm{},n`$) with superpotential $`W(\mathrm{\Phi })`$ is given by
$$S=\mathrm{d}^{\mathrm{\hspace{0.17em}2}}x\left[\mathrm{d}^4\theta K(\mathrm{\Phi }_i,\overline{\mathrm{\Phi }}_i)+\frac{1}{2}\left(\mathrm{d}^2\theta W(\mathrm{\Phi }_i)+\mathrm{d}^2\overline{\theta }\overline{W}(\overline{\mathrm{\Phi }}_i)\right)\right].$$
(2.1)
Here $`K(\mathrm{\Phi }_i,\overline{\mathrm{\Phi }}_i)`$ is the Kähler potential which defines the Kähler metric $`g_{i\overline{j}}=_i_{\overline{j}}K(\mathrm{\Phi }_i,\overline{\mathrm{\Phi }}_i)`$. If the superpotential $`W(\mathrm{\Phi })`$ is a quasi-homogeneous function with an isolated critical point (which means $`dW=0`$ can only occur at $`\mathrm{\Phi }_i=0`$) then the above action for a particular choice of $`K(\mathrm{\Phi },\overline{\mathrm{\Phi }})`$ is believed to define a superconformal theory . For a general superpotential the vacua are labeled by critical points of $`W`$, i.e., where
$$\varphi ^i(x)=\varphi _{}^i,_iW|_{\varphi _{}^i}=0.$$
(2.2)
The theory is purely massive if all the critical points are isolated and non-degenerate, which means that near the critical points $`W`$ is quadratic in fields. We assume this and label the non-degenerate critical points as $`\{\varphi _a|a=1,\mathrm{},N\}`$. In such a case the number of vacua of the theory is equal to the dimension of the local ring of $`W(\mathrm{\Phi })`$, $`=\frac{𝐂𝖨[\mathrm{\Phi }]}{_{\varphi ^i}W}`$. When we have more than one vacuum we can have solitonic states in which the boundary conditions of the fields at the left spatial infinity $`x^1=\mathrm{}`$ is at one vacuum and is different from the one at right infinity $`x^1=+\mathrm{}`$ which is in another vacuum. The geometry of solitons and their degeneracies have been extensively studied in which we will now review.
Consider a massive Landau Ginzburg theory with superpotential $`W(\mathrm{\Phi }_i)`$. Solitons are static solutions, $`\varphi ^i(x^1)`$, of the equations of motion interpolating between different vacua i.e., $`\varphi ^i(\mathrm{})=\varphi _a^i`$ and $`\varphi ^i(+\mathrm{})=\varphi _b^i`$, $`ab`$. The energy of a static field configuration interpolating between two vacua is given by
$`E_{ab}`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑x^1\left\{g_{i\overline{j}}{\displaystyle \frac{d\varphi ^i}{dx^1}}{\displaystyle \frac{d\overline{\varphi }^{\overline{i}}}{dx^1}}+\frac{1}{4}g^{i\overline{j}}_iW_{\overline{j}}\overline{W}\right\}`$
$`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}dx^1|{\displaystyle \frac{d\varphi ^i}{dx^1}}\frac{\alpha }{2}g^{i\overline{j}}_{\overline{j}}\overline{W}|^2+\text{Re}((\overline{\alpha }(W(b)W(a))).`$
Where $`g_{i\overline{j}}=_i_{\overline{j}}K`$ is the Kähler metric and $`\alpha `$ is an arbitrary phase. By choosing an appropriate $`\alpha `$ we can maximize the second term. Since $`\alpha `$ is a phase it is clear that the second term is maximum when phase of $`W(b)W(a)`$ is equal to $`\alpha `$. This implies a lower bound on the energy of the configuration,
$$E_{ab}|W(b)W(a)|.$$
(2.4)
In fact the central charge in the supersymmetry algebra in this sector is $`(W(b)W(a))`$. BPS solitons saturate this bound and therefore satisfy the equation,
$$\frac{d\varphi ^i}{dx^1}=\frac{\alpha }{2}g^{i\overline{j}}_{\overline{j}}\overline{W},\alpha =\frac{W(b)W(a)}{|W(b)W(a)|}.$$
(2.5)
An important consequence of the above equation of motion of a BPS soliton is that along the trajectory of the soliton the superpotential satisfies the equation,
$$_{x^1}W=\frac{\alpha }{2}g^{i\overline{j}}_iW_{\overline{j}}\overline{W}.$$
(2.6)
Now since the metric $`g^{i\overline{j}}`$ is positive definite, we know $`g^{i\overline{j}}_iW_{\overline{j}}\overline{W}`$ is real, and therefore the image of the BPS soliton in the W-plane is a straight line connecting the corresponding critical values $`W(a)`$ and $`W(b)`$.
The number of solitons between two vacua is equal to the number of solutions of eq. (2.5) satisfying the appropriate boundary conditions. The general way to count the number of solitons has been determined in and we will review it in the next subsection. Here we note that for the case of a single chiral superfield the number of solitons between two vacua can also be determined using eq. (2.6). Since the image of the soliton trajectory is a straight line in the $`W`$-plane therefore by looking at the pre-image of the straight line connecting the corresponding critical values in the $`W`$-plane we can determine the number of solitons between the two vacua. But since the map to the $`W`$-plane is many to one, not every pre-image of a straight line in the $`W`$-plane is a soliton. It is possible for the trajectory to start at a critical point follow a path whose image is a straight line in the $`W`$-plane and end on a point which is not a critical point but whose image in the $`W`$-plane is a critical value. The BPS solitons are those pre-images of the straight line in the $`W`$-plane which start and end on the critical points.
### 2.1 Vanishing cycles
It was shown in that the soliton numbers also have a topological description in terms of intersection numbers of vanishing cycles. The basic idea is to solve the soliton equation (2.5) along all possible directions emanating from one of the critical points. In other words the idea is to study the “wave front” of all possible solutions to the $`(\text{2.5})`$.
With no loss of generality we may assume $`\alpha =1`$. Near a critical point $`\varphi _a^i`$ we can choose coordinates $`u_a^i`$ such that,
$$W(\varphi )=W(\varphi _a)+\underset{i=1}{\overset{n}{}}(u_a^i)^2.$$
(2.7)
In this case it is easy to see that the solutions to (2.5) will have an image in the $`W`$-plane which is on a positive real line starting from $`W(\varphi _a)`$. Consider a point $`w`$ on this line. Then the space of solutions to (2.5) emanating from $`u_a^i=0`$ over this $`w`$ is a real $`(n1)`$ dimensional sphere defined by
$$\underset{i=1}{\overset{n}{}}(\text{Re}(u_a^i))^2=ww_a,\text{Im}(u_a^i)=0.$$
(2.8)
where $`w_a=W(\varphi _a)`$. Note that as we take $`ww_a`$ the sphere vanishes. This is the reason for calling these spheres “vanishing cycles”. As we move away, the wavefront will no longer be as simple as near the critical point, but nevertheless over each point $`w`$ on the positive real line emanating from $`w_a=W(\varphi _a)`$ the pre-image is a real $`(n1)`$ dimensional homology cycle $`\mathrm{\Delta }_a`$ in the $`n1`$ complex manifold defined by $`W^1(w)`$. Similarly as we move from $`w_b`$ toward $`w_a`$ there is a cycle $`\mathrm{\Delta }_b`$ evolving according to the soliton equation eq. (2.5) (this would correspond to $`\alpha =1`$). For a fixed value of $`w`$ we can compare $`\mathrm{\Delta }_a`$ and $`\mathrm{\Delta }_b`$. Solitons originating from $`\varphi _a`$ and traveling all the way to $`\varphi _b`$ correspond to the points in the intersection $`\mathrm{\Delta }_a\mathrm{\Delta }_b`$. This number, counted with appropriate signs is the intersection number of the cycles, $`\mathrm{\Delta }_a\mathrm{\Delta }_b`$. The intersection number counts the number of solitons weighed with $`(1)^F`$ for the lowest component of each soliton multiplet . This is independent of deformation of the D-terms. In particular this measures the net number of solitons that cannot disappear by deformations in the D-terms. We will denote this number by $`A_{ab}`$ and sometimes loosely refer to it as the number of solitons between $`a`$ and $`b`$. We thus have
$$A_{ab}=\mathrm{\Delta }_a\mathrm{\Delta }_b.$$
(2.9)
Note that to calculate the intersection numbers we have to consider the two cycles $`\mathrm{\Delta }_a`$ and $`\mathrm{\Delta }_b`$ in the same manifold $`W^1(w)`$. Since the intersection number is topological, a continuous deformation does not change them and hence we can actually calculate them using some deformed path in the $`W`$-plane (rather than the straight line) as long as the path we are choosing is homotopic to the straight line. How we transport the cycle along the path will not change the intersection numbers as that is topological and nothing is discontinuous, as long as the paths have the same homotopy class in the $`W`$-plane with the critical values deleted. One way, but not the only way, to transport vanishing cycles along arbitrary paths, is to use the soliton equation $`(\text{2.5})`$ but instead of having a fixed $`\alpha `$, as would be the case for a straight line, choose $`\alpha `$ to be $`e^{i\theta }`$ where $`\theta `$ denotes the varying slope of the path.
Let us fix a point $`w`$ in the $`W`$-plane. For each critical point $`a`$ of $`W`$, we choose an arbitrary path in the $`W`$-plane emanating from $`W(a)`$ and ending on $`w`$, but not passing through other critical values. This yields $`N`$ cycles $`\mathrm{\Delta }_a`$ over $`W^1(w)`$ and it is known that these form a complete basis for the middle-dimensional homology cycles of $`W^1(w)`$. Moreover, if we choose different paths the vanishing cycle we get is a linear combination of the above and the relation between them is known through the Picard Lefshetz theory as we will now review.
### 2.2 Picard-Lefshetz monodromy
As we have discussed the basis for the vanishing cycle over each point $`w`$ in the $`W`$-plane depends on the choice of paths connecting it to the critical point. Picard-Lefshetz monodromy relates how the basis changes if we change paths connecting $`w`$ to the critical values. This is quite important for the study of solitons and leads to a jump in the soliton numbers. To explain the physical motivation for the question, consider three critical values $`W(a),W(b)`$ and $`W(c)`$ depicted in Fig. 2(a), with no other critical values nearby. Suppose we wish to compute the number of solitons between them. According to our discussion above we need to connect the critical values by straight lines in the $`W`$-plane and ask about the intersection numbers of the corresponding cycles. As discussed above this is the same, because of invariance of intersection numbers under deformation, as the intersection numbers of the vanishing cycles over the point $`w`$ connecting to the three critical values as shown in Fig. 2(a). Thus the soliton number is $`A_{ij}=\mathrm{\Delta }_i\mathrm{\Delta }_j`$. However suppose now that we change the superpotential $`W`$ so that the critical values change according to what is depicted in Fig. 2(b), and that the $`W(b)`$ passes through the straight line connecting $`W(a)`$ and $`W(c)`$. In this case to find the soliton numbers between the $`a`$ vacuum and the $`c`$ vacuum we have to change the homotopy class of the path connecting $`w`$ to the critical value $`W(a)`$ as depicted by Fig. 2(b).
In particular the homology element corresponding to vanishing cycle $`a`$ changes $`\mathrm{\Delta }_a\mathrm{\Delta }_a^{}`$ and we need to find out how it changes. Picard-Lefshetz theory gives a simple formula for this change. In particular it states that
$`\mathrm{\Delta }_a^{}=\mathrm{\Delta }_a\pm (\mathrm{\Delta }_a\mathrm{\Delta }_b)\mathrm{\Delta }_b.`$ (2.10)
The sign in the above formula is determined once the orientation of the cycles are fixed and will depend on the handedness of the crossing geometry (see ). This is perhaps most familiar to string theorists in the context of moduli space of Riemann surfaces, where if we consider a point on the moduli space of Riemann surfaces where a 1-cycle shrinks to zero, as we go around this point, all the other cycles intersecting it will pick up a monodromy in the class of the vanishing cycle (the case of the torus and the $`\tau \tau +1`$ is the most familiar case, where the $`b`$ cycle undergoes a monodromy $`bb+a`$).
As a consequence of the above formula we can now find how the number of solitons between the $`a`$ and the $`c`$ vacuum change. We simply have to take the inner product $`\mathrm{\Delta }_a^{}\mathrm{\Delta }_c`$ and we find
$`A_{ac}^{}=A_{ac}\pm A_{ab}A_{bc}.`$
### 2.3 Non-compact $`n`$ Cycles
An equivalent description which will be important for later discussion involves defining soliton numbers in terms of the intersection numbers of $`n`$ real dimensional non-compact cycles which are closely related to the $`n1`$ dimensional vanishing cycles we have discussed. The idea is to consider the basis for the vanishing cycles in the limit where the point $`we^{i\theta }\mathrm{}`$. Let us consider the case where $`\theta =0`$. In this case we are taking $`w`$ to go to infinity along the positive real axis. Let us assume that the imaginary part of the critical values are all distinct. In this case a canonical choice of paths to connect the critical points to $`w`$ is along straight lines starting from the critical values $`W(a)`$ stretched along the positive real axis. We denote the corresponding non-compact $`n`$ dimensional cycles by $`\gamma _a`$. Then we have
$$W(\gamma _a)=I_a,\text{and}\gamma _a\mathrm{\Delta }_a|_{w+\mathrm{}},$$
(2.11)
where
$$I_a\{w_a+t|t[0,\mathrm{})\}.$$
(2.12)
Two such cycles are shown in Figure 3.
Let $`B`$ be the region of $`𝐂^n`$ where $`\mathrm{Re}W`$ is larger than a fixed value which is chosen sufficiently large. The non-compact cycles $`\gamma _a`$ can be viewed as elements of the homology group $`H_n(𝐂^n,B)`$ corresponding to $`n`$-cycles with boundary in $`B`$, and again it can be shown that they provide a complete basis for such cycles.
For a pair of distinct critical points, $`a`$ and $`b`$, the non-compact cycles $`\gamma _a`$ and $`\gamma _b`$ do not intersect with each other, since their images in the $`W`$-plane are parallel to each other (and are separate from each other in the present situation). In this situation we consider deforming the second cycle $`\gamma _b`$ so that its image in the $`W`$-plane is rotated with an infinitesimally small positive angle $`ϵ`$ against the real axis. We denote this deformed cycle by $`\gamma _b^{}`$. We define the “intersection number” of $`\gamma _a`$ and $`\gamma _b`$ as the geometric intersection number of $`\gamma _a`$ and $`\gamma _b^{}`$. Depending on whether $`\mathrm{Im}W(a)`$ is smaller or larger than $`\mathrm{Im}W(b)`$, the images of $`\gamma _a`$ and $`\gamma _b^{}`$ in the $`W`$-plane do not or do intersect with each other.
In the former case the “intersection number” is of course zero. In the latter case as shown in Fig. 4, the intersection number $`\gamma _a\gamma _b^{}`$ is counted by going to the point on the $`W`$-plane where their images intersect and asking what is the intersection of the corresponding vanishing cycles $`\mathrm{\Delta }_a\mathrm{\Delta }_b`$. Thus the intersection of these $`n`$-dimensional cycles has the information about the soliton numbers. In particular if there are no extra critical values between the $`I_a`$ and $`I_b`$ we will have
$`\gamma _a\gamma _b^{}=A_{ab},ab.`$ (2.13)
If there are extra critical values between $`I_a`$ and $`I_b`$ then these intersection numbers are related to the soliton numbers by the Picard-Lefshetz action as discussed before.
In the next section we will see that the cycles $`\gamma _a`$ defined through parallel transport by the soliton equation (2.5) can be viewed as D-branes for LG models that preserve half of the supersymmetries on the worldsheet. There we will also see that the “intersection number” of $`\gamma _a`$ and $`\gamma _b`$ as defined above can be interpreted as the supersymmetric index for the worldsheet theory of open strings stretched between these cycles.
### 2.4 Examples
In this section we are going to discuss some examples for the soliton numbers in the case of LG models. We will concentrate on LG models representing $`𝒩=2`$ minimal models as well as the LG models mirror to $`\text{I}\text{P}^N`$ sigma models.
#### 2.4.1 Deformed $`𝒩=2`$ Minimal models
$`𝒩=2`$ minimal models are realized as the infra-red fixed point of LG models . The soliton numbers for the deformed version of these theories has been studied in detail for the massive deformations of the A-series minimal models , which we will now review.
The $`k`$-th minimal model is described by an LG theory with one chiral superfield $`X`$ with superpotential
$$W(X)=\frac{1}{k+2}X^{k+2}.$$
(2.14)
If we add generic relevant operators to the superpotential we can deform this theory to a purely massive theory. In this case we will get $`k+1`$ vacua and we can ask how many solitons we get between each pair. For example if we consider the integrable deformation,
$$W(X)=\frac{1}{k+2}X^{k+2}X,$$
(2.15)
then there are $`k+1`$ vacua which are solutions of $`dW=0`$ given by $`X=e^{\frac{2\pi in}{k+1}},n=0,\mathrm{},k`$. In this case one can count the preimage of the straight lines in the W-plane and ask which ones connect critical points and in this way compute the number of solitons. It turns out that in this case there is exactly one soliton connecting each pair of critical points. If we deform $`W`$ the number of solitons will in general change as reviewed above. In this case one can show (by taking proper care of the relevant signs in the soliton number jump) that there is always at most one soliton between vacua. The precise number can be determined starting from the above symmetric configuration (see ). The analog of the non-compact 1-cycles $`\gamma _i`$ in this case will be discussed in more detail in section 4 after we discuss their relevance as D-branes in section 3. They are cycles in the $`X`$-plane which asymptote an $`(k+2)`$-th root of unity as $`X\mathrm{}`$. That there are $`k+1`$ inequivalent such homology classes for $`H_1(𝐂,ReW=\mathrm{})`$ is related to the fact that there are $`k+1`$ such classes defined by $`\gamma `$’s up to linear combinations.
#### 2.4.2 $`\text{I}\text{P}^{N1}`$
We next consider the $`\text{I}\text{P}^{N1}`$ sigma model. The soliton matrix of the non-linear sigma model with target space $`\text{I}\text{P}^{N1}`$ can be computed directly by studying the $`tt^{}`$ equations and their relations to soliton numbers . This has been done in . The $`tt^{}`$ equations are, however, very difficult to solve for more non-trivial spaces such as toric del Pezzos. The mirror LG theory obtained in provides a simple way of calculating the soliton matrix. We start with the case $`N=2`$ where we can present explicit solutions to the soliton equation.
The Landau Ginzburg theory which is mirror to the non-linear sigma model with $`\text{I}\text{P}^1`$ target space is the $`𝒩=2`$ sine-Gordon model with the superpotential,
$$W(x)=x+\frac{\lambda }{x}.$$
(2.16)
Here $`x=\mathrm{e}^y`$ is the single valued coordinate of the cylinder $`𝐂𝖨^\times `$ and $`\mathrm{log}\lambda `$ corresponds to the Kähler parameter of $`\text{I}\text{P}^1`$. The critical points are $`x_{}^\pm =\pm \sqrt{\lambda }`$ with critical values $`w_{}^\pm =\pm 2\sqrt{\lambda }`$. As mentioned in the previous section the BPS solitons are trajectories, $`x(t)`$, starting and ending on the critical points such that their image in the W-plane is a straight line,
$$x(t)+\frac{\lambda }{x(t)}=2\sqrt{\lambda }(2t1),t[0,1].$$
(2.17)
This is a quadratic equation with two solutions given by,
$$x(t)_\pm =\sqrt{\lambda }(2t1)\pm 2i\sqrt{\lambda }\sqrt{tt^2}=\sqrt{\lambda }e^{\pm i\text{tan}^1\frac{2\sqrt{tt^2}}{2t1}}.$$
(2.18)
Since $`x_+(t)=x_{}(t)^{}`$ and $`|x_+(t)|=|\sqrt{\lambda }|`$, there are two solitons between the two vacua such that their trajectories in the x-plane lie on two half-circles as shown in Fig. 5(a). Since $`x`$ is a $`𝐂𝖨^\times `$ coordinate we can consider the x-plane as a cylinder. Soliton trajectories on the cylinder are shown in Fig. 5(b). This description is useful in determining the intersection numbers of middle dimensional cycles. As described in the previous section the number of solitons between two critical points is given by the intersection number of middle dimensional cycles starting from the critical points. In our case there are two such cycles which are the preimages of two semi-infinite lines in the $`W`$-plane starting at the critical values as shown in Fig. 6(a). The preimage of these cycles on the cylinder is shown in Fig. 6(b).
The cycles in the x-space intersect only if the lines in the $`W`$-plane intersect each other and the intersection number in this case is two.
We now turn to the study of solitons of the $`\text{I}\text{P}^{N1}`$ sigma model. The LG theory mirror to the non-linear sigma model with $`\text{I}\text{P}^{N1}`$ target space has superpotential
$$W(X)=\underset{k=1}{\overset{N1}{}}X_k+\frac{\lambda }{X_1\mathrm{}X_{N1}}.$$
(2.19)
This superpotential has $`N`$ critical points given by
$$X_i^{(a)}=\mathrm{e}^{\frac{2\pi ia}{N}}i=1,\mathrm{},N1;a=0,\mathrm{},N1,$$
(2.20)
with the critical values
$$w_aW(\stackrel{}{X}^{(a)})=N\mathrm{e}^{\frac{2\pi ia}{N}}.$$
(2.21)
Here unlike the previous case of $`\text{I}\text{P}^1`$, to be able to solve for the preimage of a straight line, we will make an assumption about the soliton solution (For the case of $`\text{I}\text{P}^2`$ and its blowups we will also find another way to count the soliton numbers as will be discussed in section 8). Even though we will not justify this ansatz, the results we find are consistent with what is known based on $`tt^{}`$ equations. We assume that the soliton trajectory is determined by a function $`f(t)`$ such that
$$X_1=X_2=\mathrm{}=X_k=f(t)^{Nk},X_{k+1}=X_{k+2}=\mathrm{}=X_N=f(t)^k.$$
(2.22)
This parameterization of the solution satisfies the constraint $`_{i=1}^NX_i=1`$ by construction. With this ansatz the straight line equation in the $`W`$-plane becomes (for $`\lambda `$=1)
$$P(f):=kf^{Nk}+(Nk)f^k=N(1t+te^{\frac{2\pi ik}{N}}),$$
(2.23)
where the right hand side is the straight line $`w(t)`$ starting from $`w(0)=N`$ and ending on $`w(1)=N\mathrm{e}^{\frac{2\pi ik}{N}}`$. Here we have chosen the parameter $`t`$ running in the range $`[0,1]`$ that is linear in the $`W`$-plane. We are interested in the solutions which start at $`t=0`$ from $`X_i^{(0)}`$ and end at $`t=1`$ on $`X_i^{(k)}`$. This implies that $`f(0)^{Nk}=f(0)^k=1`$ and $`f(1)^{Nk}=f(1)^k=e^{\frac{2\pi ik}{N}}`$. Thus the number of solitons which satisfy eq. (2.22) is given by the number of solutions to eq. (2.23) such that $`f(0)=1`$ and $`f(1)=e^{\frac{2\pi i}{N}}`$. We will show that there is only a single solution which satisfies these conditions.
Since $`P^{}(1)=0`$ and $`P^{\prime \prime }(1)0`$, where prime denotes a differentiation with respect to $`f`$, only two trajectories start from $`f=1`$. Thus it follows that the number of solutions is less than or equal to two. From eq. (2.23) it is clear that $`f`$ can be real only at $`t=0`$. Thus a trajectory cannot cross the real axis for $`t>0`$. For $`t`$ very close to zero one of the trajectories move into the upper half plane. Since the trajectory in the upper half plane cannot cross the real axis it cannot end on $`e^{\frac{2\pi ik}{N}}`$ <sup>1</sup><sup>1</sup>1Very close to $`f=1`$ the two solutions are given by $`f_\pm =1\pm \sqrt{\frac{2t(e^{\frac{2\pi ik}{N}}1)}{k(Nk)}}`$. Thus there can be at most one solution.
To show that there actually exists a solution we will construct a solution whose image in the $`W`$-plane is homotopic to the straight line $`w(t)`$. Consider the function $`f_{}(t)=e^{\frac{2\pi i}{N}t}`$ where $`t[0,1]`$. Since
$$|P(f_{}(t))|=|ke^{2\pi it}+(Nk)||ke^{2\pi it}|+(Nk)=N,$$
(2.24)
the image of $`f_{}(t)`$ in the W-plane always lies inside the circle of radius $`N`$ and only intersects the circle for $`t=0`$ and $`t=1`$ at $`w=w_0`$ and $`w=w_k`$ respectively. Thus the image is homotopic to the straight line $`w(t)`$ and therefore there exists a solution $`f_0(t)`$ homotopic to $`f_{}(t)`$ with the required properties.
Since permuting the $`N`$ coordinates among themselves does not change the superpotential, it follows that we can choose any $`k`$ coordinates to be equal to $`f^{Nk}`$ and the remaining $`(Nk)`$ coordinates equal to $`f^k`$. Thus we see that there are $`\left(\genfrac{}{}{0pt}{}{N}{k}\right)`$ solitons between the critical points $`X_i^{(0)}`$ and $`X_i^{(k)}`$ consistent with the ansatz of eq. (2.22). In fact this is the same number anticipated by the study of tt\* equations . Note that if the $`\text{I}\text{P}^{N1}`$ has a round metric having $`SU(N)`$ symmetry, then the solitons should form representations of this group. In fact the permutations of $`X_i`$ can be viewed as the Weyl group of the $`SU(N)`$, as is clear from the derivation of the mirror in this case . It thus follows, given how the permutation acts on the solutions we have found, that in this case the solitons connecting vacua $`k`$ units apart correspond to $`k`$ fold anti-symmetric tensor product of the fundamental representation of $`SU(N)`$, a result which was derived from the large $`N`$ analysis of this theory .
## 3 D-Branes in $`𝒩=2`$ Supersymmetric Field Theories
In this section, we study the $`N=(2,2)`$ supersymmetric field theory formulated on a 1+1 dimensional worldsheet with boundaries. We mainly consider supersymmetric sigma models and Landau-Ginzburg models. We find boundary conditions that preserve half of the (worldsheet) supersymmetry. (See for earlier works.) We define and compute the supersymmetric index of a theory on an interval. We also analyze the $`𝒩=2`$ boundary entropy defined as the pairing of the boundary states and the supersymmetric ground states.
### 3.1 The Supersymmetric Boundary Conditions
Let us consider a supersymmetric sigma model on a Kähler manifold $`X`$ of dimension $`n`$ with a superpotential $`W`$. We denote the Kähler metric with respect to local complex coordinates $`z^i`$ as $`g_{i\overline{ȷ}}`$. We formulate the theory on the strip $`\mathrm{\Sigma }=𝐑\times I`$ where $`I`$ is an interval $`0x^1\pi `$ and $`𝐑`$ is parametrized by the time coordinate $`x^0`$. Here, without any loss in generality we have fixed the size of the interval to a fixed length. Changing the length of the interval is equivalent to changing the parameters in the action according to the RG flow.
The action of the system is given by
$`S`$ $`=`$ $`{\displaystyle \underset{\mathrm{\Sigma }}{}}\mathrm{d}^2x\{g_{i\overline{ȷ}}^\mu \varphi ^i_\mu \overline{\varphi }^{\overline{ȷ}}+{\displaystyle \frac{i}{2}}g_{i\overline{ȷ}}\overline{\psi }_{}^{\overline{ȷ}}(\stackrel{}{D}\stackrel{}{}_0+\stackrel{}{D}\stackrel{}{}_1)\psi _{}^i+{\displaystyle \frac{i}{2}}g_{i\overline{ȷ}}\overline{\psi }_+^{\overline{ȷ}}(\stackrel{}{D}\stackrel{}{}_0\stackrel{}{D}\stackrel{}{}_1)\psi _+^i`$ (3.1)
$`{\displaystyle \frac{1}{4}}g^{\overline{ȷ}i}_{\overline{ȷ}}\overline{W}_iW{\displaystyle \frac{1}{2}}(D_i_jW)\psi _+^i\psi _{}^j{\displaystyle \frac{1}{2}}(D_{\overline{ı}}_{\overline{ȷ}}\overline{W})\overline{\psi }_{}^{\overline{ı}}\overline{\psi }_+^{\overline{ȷ}}`$
$`+R_{i\overline{k}j\overline{l}}\psi _+^i\psi _{}^j\overline{\psi }_{}^{\overline{k}}\overline{\psi }_+^{\overline{l}}\},`$
where $`\overline{\psi }^{\overline{ȷ}}\stackrel{}{D}\stackrel{}{}_\mu \psi ^i=\overline{\psi }^{\overline{ȷ}}(D_\mu \psi )^i(D_\mu \overline{\psi })^{\overline{ȷ}}\psi ^i`$. See for other notations. The above action is the same as the component expression of (2.1) up to a boundary term. We require the equations of motion for the fields $`\varphi ^i,\psi _\pm ^i`$ to be local. This yields the following conditions on the boundary $`\mathrm{\Sigma }`$
$`g_{IJ}\delta \varphi ^I_1\varphi ^J=0,`$ (3.2)
$`g_{IJ}(\psi _{}^I\delta \psi _{}^J\psi _+^I\delta \psi _+^J)=0,`$ (3.3)
where $`\varphi ^I,\psi _\pm ^I`$ and $`g_{IJ}`$ are the components of the fields and the metric with respect to the real coordinates of the target space.
Under the supersymmetry transformation
$`\delta \varphi ^i=ϵ_+\psi _{}^iϵ_{}\psi _+^i,`$ (3.4)
$`\delta \psi _+^i=i\overline{ϵ}_{}(_0+_1)\varphi ^i+ϵ_+F^i,`$ (3.5)
$`\delta \psi _{}^i=i\overline{ϵ}_+(_0_1)\varphi ^i+ϵ_{}F^i,`$ (3.6)
where
$$F^i=\frac{1}{2}g^{i\overline{ȷ}}_{\overline{ȷ}}\overline{W}+\mathrm{\Gamma }_{jk}^i\psi _+^j\psi _{}^k,$$
(3.7)
the action varies as
$`\delta S`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mathrm{\Sigma }}{}}\mathrm{d}x^0\{ϵ_+(g_{i\overline{ȷ}}(_0+_1)\overline{\varphi }^{\overline{ȷ}}\psi _{}^i+{\displaystyle \frac{i}{2}}\overline{\psi }_+^{\overline{ı}}_{\overline{ı}}\overline{W})+ϵ_{}(g_{i\overline{ȷ}}(_0_1)\overline{\varphi }^{\overline{ȷ}}\psi _+^i{\displaystyle \frac{i}{2}}\overline{\psi }_{}^{\overline{ı}}_{\overline{ı}}\overline{W})`$
$`+\overline{ϵ}_+(g_{i\overline{ȷ}}\overline{\psi }_{}^{\overline{ȷ}}(_0+_1)\varphi ^i+{\displaystyle \frac{i}{2}}\psi _+^i_iW)+\overline{ϵ}_{}(g_{i\overline{ȷ}}\overline{\psi }_+^{\overline{ȷ}}(_0_1)\varphi ^i{\displaystyle \frac{i}{2}}\psi _{}^i_iW)\}.`$
If the boundary were absent, the action would be invariant under the full $`(2,2)`$ supersymmetry and the following four supercurrents would be conserved.
$`G_\pm ^0=g_{i\overline{ȷ}}(_0\pm _1)\overline{\varphi }^{\overline{ȷ}}\psi _\pm ^i{\displaystyle \frac{i}{2}}\overline{\psi }_{}^{\overline{ı}}_{\overline{ı}}\overline{W},G_\pm ^1=g_{i\overline{ȷ}}(_0\pm _1)\overline{\varphi }^{\overline{ȷ}}\psi _\pm ^i{\displaystyle \frac{i}{2}}\overline{\psi }_{}^{\overline{ı}}_{\overline{ı}}\overline{W},`$
$`\overline{G}_\pm ^0=g_{i\overline{ȷ}}\overline{\psi }_\pm ^{\overline{ȷ}}(_0\pm _1)\varphi ^i\pm {\displaystyle \frac{i}{2}}\psi _{}^i_iW,\overline{G}_\pm ^1=g_{i\overline{ȷ}}\overline{\psi }_\pm ^{\overline{ȷ}}(_0\pm _1)\varphi ^i+{\displaystyle \frac{i}{2}}\psi _{}^i_iW.`$
In what follows, we determine the boundary conditions on the fields $`\varphi ^i,\psi _\pm ^i`$ that preserve half of the supersymmetry. We also wish to maintain the translation symmetry that maps the worldsheet boundary to itself, which is the time translation in the present set-up. There are essentially two possibilities for the unbroken supercharges ;
(A) $`Q=\overline{Q}_++\mathrm{e}^{i\alpha }Q_{}`$ and $`Q^{}=Q_++\mathrm{e}^{i\alpha }\overline{Q}_{}`$,
(B) $`Q=\overline{Q}_++\mathrm{e}^{i\beta }\overline{Q}_{}`$ and $`Q^{}=Q_++\mathrm{e}^{i\beta }Q_{}`$.
Here $`\mathrm{e}^{i\alpha }`$ and $`\mathrm{e}^{i\beta }`$ are arbitrary phases. In both cases, the supercharges satisfy $`\{Q,Q^{}\}=2H`$, up to a possible central term. The variation parameters for these supersymmetries are $`ϵ_{}=\mathrm{e}^{i\alpha }\overline{ϵ}_+`$ for (A) while $`ϵ_{}=\mathrm{e}^{i\beta }ϵ_+`$ for (B). Conservation of the charges $`Q`$ and $`Q^{}`$ requires that the spatial component of the corresponding currents vanish at the boundary $`\mathrm{\Sigma }`$: $`\overline{G}_+^1+\mathrm{e}^{i\alpha }G_{}^1=G_+^1+\mathrm{e}^{i\alpha }\overline{G}_{}^1=0`$ for (A), and $`\overline{G}_+^1+\mathrm{e}^{i\beta }\overline{G}_{}^1=G_+^1+\mathrm{e}^{i\beta }G_{}^1=0`$ for (B).
The conditions we are interested in are the ones associated with D-branes wrapped on a submanifold $`\gamma `$ of $`X`$. Namely, we require the worldsheet boundary to be mapped to $`\gamma `$. In such a case, the derivative along the boundary $`_0\varphi ^I`$ as well as an allowed variation $`\delta \varphi ^I`$ at the boundary must be tangent to $`\gamma `$. The locality condition (3.2) then tells that $`_1\varphi ^I`$ must be normal to $`\gamma `$. The other condition (3.3) is satisfied if $`\psi _{}^I`$ and $`\psi _+^I`$ are related by an orthonormal transformation $`\psi _{}^I=M_J^I\psi _+^J`$, $`g_{IJ}M_K^IM_L^J=g_{KL}`$. In fact, supersymmetry requires this and determines the matrix $`M_J^I`$, as we now show. For simplicity, we set the phases $`\mathrm{e}^{i\alpha }`$ and $`\mathrm{e}^{i\beta }`$ to be equal to $`1`$; the general case can be easily recovered by $`U(1)_V`$ and $`U(1)_A`$ rotations. Then, both (A) and (B) contains an $`N=1`$ subalgebra generated by the variations with parameter $`ϵ_+=iϵ`$ and $`ϵ_{}=iϵ`$ where $`ϵ`$ is real. Expressed in the real coordinates, the action varies as
$$\delta S=\frac{iϵ}{2}\underset{\mathrm{\Sigma }}{}dx^0\left\{g_{IJ}_0\varphi ^I(\psi _{}^J\psi _+^J)g_{IJ}_1\varphi ^I(\psi _{}^J+\psi _+^J)\frac{i}{2}(\psi _{}^I+\psi _+^I)_I(W\overline{W})\right\}.$$
(3.9)
Since $`_0\varphi ^I`$ and $`_1\varphi ^i`$ are tangent and normal to $`\gamma `$, the invariance of the action requires $`iϵ(\psi _{}^I\psi _+^I)`$ and $`iϵ(\psi _{}^I+\psi _+^I)`$ to be normal and tangent to $`\gamma `$ respectively. This means that
$$\psi _{}^I=\{\begin{array}{cc}\psi _+^I\hfill & I:\text{tangent}\hfill \\ \psi _+^I\hfill & I:\text{normal}\hfill \end{array}$$
(3.10)
for a choice of coordinates that separates the tangent and normal directions. Furthermore, invariance of $`S`$ requires $`W\overline{W}`$ to be a constant along $`\gamma `$.
As we will see, A-type supersymmetry requires $`\gamma `$ to be a middle dimensional Lagrangian submanifold whose image in the $`W`$-plane is a straight line, while B-type supersymmetry requires $`\gamma `$ to be a holomorphic submanifold on which $`W`$ is a constant.
### A-Type Supersymmetry
We first consider A-type supersymmetry with the trivial phase $`\mathrm{e}^{i\alpha }=1`$, which is generated by the variations with parameters $`ϵ_{}=\overline{ϵ}_+`$ and $`\overline{ϵ}_{}=ϵ_+`$. The bosonic fields $`\varphi ^i`$ transform as
$`\delta \varphi ^i`$ $`=`$ $`ϵ_+\psi _{}^i\overline{ϵ}_+\psi _+^i`$ (3.11)
$`=`$ $`ϵ_1(\psi _{}^i\psi _+^i)+iϵ_2(\psi _{}^i+\psi _+^i),`$
where $`ϵ_1`$ and $`ϵ_2`$ are the real and the imaginary parts of $`ϵ_+`$; $`ϵ_+=ϵ_1+iϵ_2`$. This shows that, for a real parameter $`ϵ`$, $`ϵ(\psi _{}^i\psi _+^i)`$ and $`iϵ(\psi _{}^i+\psi _+^i)`$ are the holomorphic components of tangent vectors of $`\gamma `$. On the other hand, $`N=1`$ supersymmetry requires $`iϵ(\psi _{}^i\psi _+^i)`$ and $`iϵ(\psi _{}^i+\psi _+^i)`$ are the holomorphic components of normal and tangent vectors of $`\gamma `$ respectively. Thus, multiplication by $`i=\sqrt{1}`$ on the holomorphic components sends tangent vectors to normal vectors and vice versa. Namely, the cycle $`\gamma `$ must be a middle dimensional Lagrangian submanifold of $`X`$ (where $`X`$ is considered as a symplectic manifold defined by the Kähler form). The supersymmetry transformation of the tangent vector $`ϵ(\psi _{}^i\psi _+^i)`$ is
$$\delta [ϵ(\psi _{}^i\psi _+^i)]=2iϵϵ_1(_0\varphi ^i)+2iϵϵ_2(i_1\varphi ^i+F^i).$$
(3.12)
This must again be tangent to $`\gamma `$ when $`\delta \varphi ^i=0`$ (i.e. $`\psi _\pm ^i=0`$ for which $`F^i=\frac{1}{2}g^{i\overline{ȷ}}_{\overline{ȷ}}\overline{W}`$). We note that $`iϵϵ_1`$ and $`iϵϵ_2`$ are real parameters, $`(iϵϵ_i)^{}=iϵ_iϵ=iϵϵ_i`$. Since $`_0\varphi ^i`$, $`i_1\varphi ^i`$ are both tangent to $`\gamma `$, the vector $`v^i=g^{i\overline{ȷ}}_{\overline{ȷ}}\overline{W}`$ must be tangent to $`\gamma `$. This is consistent with the requirement of $`N=1`$ supersymmetry that $`W\overline{W}`$ must be a constant along $`\gamma `$, since the vector $`v^I`$ annihilates $`W\overline{W}`$
$$v^I_I(W\overline{W})=|W|^2|W|^2=0.$$
(3.13)
It is easy to see that under these conditions the action $`S`$ is invariant under the full A-type supersymmetry with $`\mathrm{e}^{i\alpha }=1`$ and also that no other condition is required.
It is easy to recover the phase $`\mathrm{e}^{i\alpha }`$ using the $`U(1)_V`$ symmetry which rotates the fermions as $`\psi _\pm ^i\mathrm{e}^{i\alpha /2}\psi _\pm ^i`$ and $`\overline{\psi }_\pm ^{\overline{ȷ}}\mathrm{e}^{i\alpha /2}\overline{\psi }_\pm ^{\overline{ȷ}}`$ and the superpotential as $`W\mathrm{e}^{i\alpha }W`$. The boundary conditions on the fermions (3.10) are rotated accordingly. The condition on the cycle is also rotated: the cycle $`\gamma `$ is a middle dimensional Lagrangian submanifold of $`X`$ with respect to the Kähler form, whose image in the $`W`$-plane is a straight line in the $`\mathrm{e}^{i\alpha }`$-direction.
The axial $`U(1)`$ R-symmetry is not broken by the boundary condition. Indeed, $`ϵ(\psi _{}^i\psi _+^i)`$ and $`iϵ(\psi _{}^i+\psi _+^i)`$, which are holomorphic component of tangent vectors to $`\gamma `$ (for $`\mathrm{e}^{i\alpha }=1`$), are rotated within themselves by the axial rotation $`\psi _{}^i\mathrm{e}^{\pm i\theta _A}\psi _{}^i`$. If not anomalous (i.e. if $`c_1(X)=0`$), the axial R-charge is conserved and the spatial component of the current must vanish at the boundary, $`J_A^1=0`$ at $`\mathrm{\Sigma }`$. We note that the conserved supercharges $`Q`$ and $`Q^{}`$ have axial charge $`1`$ and $`1`$ respectively.
The basic example of a cycle satisfying this condition is the wave-front trajectory emanating from a critical point of the superpotential. We consider here the case $`\mathrm{e}^{i\alpha }=1`$ for simplicity. Let $`p_{}X`$ be a non-degenerate critical point of $`W`$ and let us consider a wave-front trajectory $`\gamma _p_{}`$ emanating from $`p_{}`$ in the positive real direction. As discussed in section 2 this corresponds to the totality of all potential soliton solutions starting at $`p_{}`$ whose image in the $`W`$-plane is stretched along the positive real axis. We recall that the one parameter family of maps $`f_t`$ generated by the vector field $`v^i=g^{i\overline{ȷ}}_{\overline{ȷ}}\overline{W}`$ acts on $`\gamma _p_{}`$ and any point on it is mapped by $`f_t`$ to $`p_{}`$ in the limit $`t\mathrm{}`$. By definition, $`\gamma _p_{}`$ is of middle dimension and the image in the $`W`$-plane is a straight-line in the real direction. To see that $`\gamma _p_{}`$ is a Lagrangian submanifold of $`X`$ with respect to the Kähler form $`\omega =ig_{i\overline{ȷ}}\mathrm{d}z^i\mathrm{d}\overline{z}^{\overline{ȷ}}`$, it is crucial to note that
$$i_v\omega =i(g_{i\overline{ȷ}}v^i\mathrm{d}\overline{z}^{\overline{ȷ}}g_{i\overline{ȷ}}\mathrm{d}z^i\overline{v}^{\overline{ȷ}})=i\mathrm{d}(\overline{W}W),$$
(3.14)
and hence
$$_v\omega =\mathrm{d}i_v\omega +i_v\mathrm{d}\omega =0.$$
(3.15)
Thus, $`\omega `$ is invariant under the diffeomorphisms $`f_t`$. Let $`V_1`$ and $`V_2`$ be tangent vectors of $`\gamma _p_{}`$ at any point. Since the Kähler form is $`f_t`$-invariant, $`\omega (f_tV_1,f_tV_2)=(f_t^{}\omega )(V_1,V_2)`$ is independent of $`t`$. However, in the limit $`t\mathrm{}`$, the vectors $`f_tV_i`$ become the zero vector at $`p_{}`$. Thus, we have shown $`\omega (V_1,V_2)=0`$. Namely, $`\gamma _p_{}`$ is Lagrangian.
### B-Type Supersymmetry
We next consider B-type supersymmetry with the phase $`\mathrm{e}^{i\beta }=1`$ which is generated by $`ϵ_{}=ϵ_+`$ and $`\overline{ϵ}_{}=\overline{ϵ}_+`$. The bosonic fields $`\varphi ^i`$ transform as
$$\delta \varphi ^i=ϵ_+(\psi _{}^i+\psi _+^i).$$
(3.16)
Since $`ϵ_+`$ is a complex parameter, this shows that the tangent space to $`\gamma `$ is invariant under the multiplication by $`i=\sqrt{1}`$ on the holomorphic components. Namely, the cycle $`\gamma `$ must be a complex submanifold of $`X`$. The supersymmetry transformation of the tangent vector $`\psi _{}^i+\psi _+^i`$ is $`\delta (\psi _{}^i+\psi _+^i)=2i\overline{ϵ}_+_0\varphi ^i`$ which is indeed tangent to $`\gamma `$. On the other hand, the normal vector $`\psi _{}^i\psi _+^i`$ transforms as
$$\delta (\psi _{}^i\psi _+^i)=2i\overline{ϵ}_+_1\varphi ^i+ϵ_+g^{i\overline{ȷ}}_{\overline{ȷ}}\overline{W},$$
(3.17)
at $`\psi _\pm ^i=0`$ for which $`\delta \varphi ^i=0`$. This must again be normal to $`\gamma `$. Since $`_1\varphi ^i`$ is normal to $`\gamma `$, this requires that $`n^i=g^{i\overline{ȷ}}_{\overline{ȷ}}\overline{W}`$ is also a normal vector to $`\gamma `$. Namely, for a tangent vector $`v^i`$ we have
$$0=g_{i\overline{ȷ}}v^i\overline{n}^{\overline{ȷ}}=v^i_iW.$$
(3.18)
Thus, not only the imaginary part $`W\overline{W}`$ but $`W`$ itself must be a constant on $`\gamma `$. It is easy to see that under these conditions the action is invariant under the full B-type supersymmetry with $`\mathrm{e}^{i\beta }=1`$ and also that no other condition is required.
It is again easy to recover the phase $`\mathrm{e}^{i\beta }`$ using the $`U(1)_A`$ symmetry which rotates the fermions as $`\psi _\pm ^i\mathrm{e}^{\pm i\beta /2}\psi _\pm ^i`$ and $`\overline{\psi }_\pm ^{\overline{ȷ}}\mathrm{e}^{i\beta /2}\overline{\psi }_\pm ^{\overline{ȷ}}`$. The boundary conditions on the fermions (3.10) are rotated accordingly, but the condition on the cycle remains the same: the cycle $`\gamma `$ is a complex submanifold of $`X`$ on which $`W`$ is a constant.
The vector $`U(1)`$ R-symmetry is not broken by the boundary condition. Indeed, the tangent vector $`ϵ(\psi _{}^i+\psi _+^i)`$ to $`\gamma `$ (for $`\mathrm{e}^{i\beta }=1`$) is rotated by phase under the vector rotation $`\psi _{}^i\mathrm{e}^{i\theta _V}\psi _{}^i`$ and hence remains tangent. If not broken by the superpotential, the vector R-charge is conserved and the spatial component of the current must vanish at the boundary, $`J_V^1=0`$ at $`\mathrm{\Sigma }`$. We note that the conserved supercharges $`Q`$ and $`Q^{}`$ has the vector charge $`1`$ and $`1`$ respectively.
#### 3.1.1 Inclusion of the B-field
We can deform the theory by adding the following term to the action (3.1)
$$\frac{1}{2}\underset{\mathrm{\Sigma }}{}B_{IJ}d\varphi ^I\mathrm{d}\varphi ^J,$$
(3.19)
where $`B=\frac{1}{2}B_{IJ}\mathrm{d}x^I\mathrm{d}x^J`$ is a closed two-form on the manifold $`X`$. This term alters the condition (3.2) of locality for the bosonic equations of motion as
$$\delta \varphi ^I(g_{IJ}_1\varphi ^J+B_{IJ}_0\varphi ^J)=0,$$
(3.20)
but the condition (3.3) for the fermionic equations of motion remains the same.
We look for the boundary conditions associated with the D-branes wrapped on a cycle $`\gamma `$ in $`X`$ which preserves A-type or B-type supersymmetry. By definition and by the requirement (3.20), the bosonic fields must obey the boundary conditions
$$\begin{array}{cc}g_{IJ}_1\varphi ^J+B_{IJ}_0\varphi ^J=0,\hfill & I:\text{tangent},\hfill \\ _0\varphi ^I=0,\hfill & I:\text{normal},\hfill \end{array}$$
(3.21)
where we have chosen the coordinates that separate the tangent and the normal directions. For invariance under the $`N=1`$ supersymmetry generated by $`ϵ_+=iϵ`$ and $`ϵ_{}=iϵ`$ with $`ϵ`$ being real (which is contained in both (A) and (B) supersymmetries with the trivial phases), the following boundary conditions on the fermions are required:
$$\begin{array}{cc}g_{IJ}(\psi _{}^J\psi _+^J)B_{IJ}(\psi _{}^J+\psi _+^J)=0,\hfill & I:\text{tangent},\hfill \\ \psi _{}^I+\psi _+^I=0,\hfill & I:\text{normal}.\hfill \end{array}$$
(3.22)
This also guarantees the condition (3.3). We also obtain the condition that the imaginary part of $`W`$ is a constant along $`\gamma `$.
Proceeding as in the case without $`B`$-field, we obtain the following conditions on the cycle $`\gamma `$ for the A- and B-type supersymmetry to be preserved. We only state the conditions for the cases with the trivial phase $`\mathrm{e}^{i\alpha }=\mathrm{e}^{i\beta }=1`$ since the generalization is clear.
A-type Supersymmetry
$`\gamma `$ is a middle dimensional Lagrangian submanifold of $`X`$ on which (not only the Kähler form but also) the $`B`$-field is annihilated, $`B|_\gamma =0`$. The image in the $`W`$-plane must be a straight line in the real direction.
B-type Supersymmetry
$`\gamma `$ is a complex submanifold of $`X`$. $`B`$-field evaluated on the holomorphic tangent vectors to $`\gamma `$ is zero, $`(B|_\gamma )^{(2,0)}=0`$. Also, $`W`$ must be a constant on $`\gamma `$.
#### 3.1.2 Coupling to the Gauge Fields on the Branes
We can couple the worldsheet boundaries to the gauge fields on the branes. In the case of the strip $`\mathrm{\Sigma }=𝐑\times I`$, this corresponds to adding to the action (3.1) the terms
$$\underset{\mathrm{\Sigma }}{}A_Md\varphi ^M=\underset{x^1=\pi }{}dx^0_0\varphi ^{M_b}A_{M_b}^{(b)}\underset{x^1=0}{}dx^0_0\varphi ^{M_a}A_{M_a}^{(a)},$$
(3.23)
where $`A^{(a)}`$ and $`A^{(b)}`$ are the $`U(1)`$ gauge fields on the branes $`\gamma _a`$ and $`\gamma _b`$ on which the left and the right boundaries end. (We use $`M,N,\mathrm{}`$ for coordinate indices on the branes.) If the left and the right boundaries are coupled to the same gauge field $`A`$ that extends to the whole target space $`X`$, the boundary terms (3.23) can be written as
$$\frac{1}{2}\underset{\mathrm{\Sigma }}{}F_{IJ}d\varphi ^I\mathrm{d}\varphi ^J$$
(3.24)
where $`F=\frac{1}{2}F_{IJ}\mathrm{d}x^I\mathrm{d}x^J`$ is the curvature of the gauge field, $`F=\mathrm{d}A`$. Thus, in this case, we can treat the gauge field coupling in the same way as the coupling to the $`B`$-field. In particular, we have the local equations of motion and $`N=1`$ supersymmetry by imposing the boundary conditions (3.21) and (3.22) with $`BF`$. When the cycle $`\gamma `$ is a middle dimensional Lagrangian submanifold of $`X`$ whose image in the $`W`$-plane is a straight line, the theory is invariant under A-type supersymmetry if the gauge field is flat on the cycle, $`F|_\gamma =0`$. When the cycle $`\gamma `$ is a complex submanifold of $`X`$ on which $`W`$ is a constant, the theory is invariant under B-type supersymmetry if the gauge field has a $`(1,1)`$ curvature on $`\gamma `$, $`(F|_\gamma )^{(2,0)}=0`$, namely, if $`A|_\gamma `$ determines a holomorphic line bundle on $`\gamma `$.
The conclusion obtained above remains valid even if the left and the right boundary components are coupled to different gauge fields that are defined only on the branes. Thus, one can deform the A-type supersymmetric theory by flat gauge fields on $`\gamma `$ while B-type supersymmetric theory can be deformed by holomorphic line bundles on $`\gamma `$.
### An Alternative Formulation for B-type D-branes
For A-type D-branes with a flat gauge field, the boundary condition given by (3.21) and (3.22) (with $`B_{IJ}F_{IJ}`$) is the same as the standard one
$$\begin{array}{cc}_1\varphi ^I=0,\psi _{}^I\psi _+^I=0,\hfill & I:\text{tangent},\hfill \\ _0\varphi ^I=0,\psi _{}^I+\psi _+^I=0\hfill & I:\text{normal}.\hfill \end{array}$$
(3.25)
However, it is in general different from (3.25) for B-type D-branes where the gauge field is not necessarily flat. There is actually an alternative formulation for B-type D-branes where we still impose the standard boundary condition (3.25). It is easy to see that (with $`\psi ^I:=(\psi _+^I+\psi _{}^I)/2`$)
$$\underset{\mathrm{\Sigma }}{}dx^0\left\{_0\varphi ^MA_M+iF_{MN}\psi ^M\psi ^N\right\}$$
(3.26)
is invariant by itself under the B-type supersymmetry if the gauge field is holomorphic, $`F_{mn}=F_{\overline{m}\overline{n}}=0`$. Thus, instead of (3.23), one can add the boundary term (3.26) without breaking the B-type supersymmetry of the bulk action (3.1) which holds under (3.25). We note however that the equations of motion for the fields $`\varphi ^I`$, $`\psi _\pm ^I`$ are modified by boundary terms. This formulation was used in to study the fluctuation of the target space gauge fields in string theory.
### Non-Abelian Gauge Fields
One can generalize the above analysis to non-abelian $`U(k)`$ gauge group . In this case, the path-integral weight $`\mathrm{exp}(iS)`$ is accompanied by the matrix factors
$$P_\mathrm{\Sigma }\mathrm{exp}\left(i_\mathrm{\Sigma }dx^0\left\{_0\varphi ^MA_M+iF_{MN}\psi ^M\psi ^N\right\}\right)$$
(3.27)
where $`P_\mathrm{\Sigma }`$ is (the product of) the path-ordering along the boundary $`\mathrm{\Sigma }`$. Under the standard boundary condition (3.25), the weight (3.27) is invariant under A-type supersymmetry if $`A`$ is flat, $`F_{MN}=0`$, while it is invariant under B-type supersymmetry if $`A`$ is holomorphic, $`F_{mn}=F_{\overline{m}\overline{n}}=0`$.
### 3.2 Supersymmetric Ground States
As in any supersymmetric field theory, in the theory on the segment $`I=[0,\pi ]`$ with the boundary condition that preserves A or B-type supersymmetry, one can define the supersymmetric index $`\mathrm{Tr}(1)^F`$ which is invariant under deformations of the theory. We denote this index as
$$I(a,b)=\mathrm{Tr}(1)^F,$$
(3.28)
where $`a`$ and $`b`$ are the boundary conditions at the left and the right boundaries.<sup>1</sup><sup>1</sup>1The index (3.28) for supersymmetric D-branes in Calabi-Yau manifolds was studied in . We shall compute this index in the two basic examples; A-type D-branes in Landau-Ginzburg models and B-type D-branes in non-linear sigma models. Actually, in the LG models (not only the index but also) the complete spectrum of supersymmetric ground states can be determined. This can also be done for non-linear sigma models under a certain condition on the cohomology of the gauge bundles. For simplicity, we set the phases $`\mathrm{e}^{i\alpha }=1`$ and $`\mathrm{e}^{i\beta }=1`$.
#### 3.2.1 Landau-Ginzburg Models
Let us consider a LG model with superpotential $`W`$. We assume that the bosonic potential $`U=|W|^2`$ diverges at infinity in the configuration space $`X`$. We also assume that there is no non-trivial $`B`$ field and we will not consider coupling to gauge field on the branes. Let $`a`$ and $`b`$ be two non-degenerate critical points of $`W`$. We consider the wave-front trajectories $`\gamma _a`$ and $`\gamma _b`$ emanating from $`a`$ and $`b`$ in the positive real direction in the $`W`$-plane. We assume for now that the half-lines $`W(\gamma _a)`$ and $`W(\gamma _b)`$ are separated in the imaginary direction, and there is no other critical values of $`W`$ between them. We consider the theory on $`[0,\pi ]`$ where the left boundary $`x^1=0`$ is mapped to $`\gamma _a`$ and the right boundary $`x^1=\pi `$ is mapped to $`\gamma _b`$. For the boundary condition described earlier, the theory is invariant under A-type supersymmetry generated by the supercharges $`Q=\overline{Q}_++Q_{}`$ and $`Q^{}=Q_++\overline{Q}_{}`$, which are expressed as
$`Q`$ $`=`$ $`\sqrt{2}{\displaystyle _0^\pi }dx^1\left\{\overline{\psi }_+^{\overline{ȷ}}\left(g_{i\overline{ȷ}}(_0+_1)\varphi ^i+{\displaystyle \frac{i}{2}}_{\overline{ȷ}}\overline{W}\right)+\psi _{}^i\left(g_{i\overline{ȷ}}(_0_1)\overline{\varphi }^{\overline{ȷ}}+{\displaystyle \frac{i}{2}}_iW\right)\right\},`$
$`Q^{}`$ $`=`$ $`\sqrt{2}{\displaystyle _0^\pi }dx^1\left\{\overline{\psi }_{}^{\overline{ȷ}}\left(g_{i\overline{ȷ}}(_0_1)\varphi ^i{\displaystyle \frac{i}{2}}_{\overline{ȷ}}\overline{W}\right)+\psi _+^i\left(g_{i\overline{ȷ}}(_0+_1)\overline{\varphi }^{\overline{ȷ}}{\displaystyle \frac{i}{2}}_iW\right)\right\}.`$
The supercharges $`Q`$ and $`Q^{}`$ are nilpotent and satisfy the anti-commutation relation
$$\{Q,Q^{}\}=4(H+\mathrm{\Delta }\mathrm{Im}W),$$
where $`\mathrm{\Delta }\mathrm{Im}W=\mathrm{Im}W(b)\mathrm{Im}W(a)`$ is the separation of the two half-lines in the imaginary direction. We shift the definition of the Hamiltonian as $`\stackrel{~}{H}=H+\mathrm{\Delta }\mathrm{Im}W`$ so that the supersymmetry algebra takes the standard form $`\{Q,Q^{}\}=4\stackrel{~}{H}`$. Since $`\mathrm{\Delta }\mathrm{Im}W`$ is a constant, this is done simply by the shift of the action
$$\stackrel{~}{S}=S\underset{x^1=\pi }{}dx^0\mathrm{Im}W+\underset{x^1=0}{}dx^0\mathrm{Im}W.$$
(3.30)
The index can be defined by $`I(a,b)=\mathrm{Tr}(1)^F\mathrm{e}^{\beta \stackrel{~}{H}}`$, and only the ground states with energy $`\stackrel{~}{H}=0`$ can contribute to this. One can see from the expressions (LABEL:QpmbQpm) that $`Q=Q^{}=0`$ for a static configuration such that
$$_1\varphi ^i=\frac{i}{2}g^{i\overline{ȷ}}_{\overline{ȷ}}\overline{W}.$$
(3.31)
Namely, the supersymmetry is classically preserved for a static configuration that goes from $`\gamma _a`$ to $`\gamma _b`$, straight down in the negative imaginary direction of the $`W`$-plane. Such a configuration would indeed have $`H=\mathrm{\Delta }\mathrm{Im}W`$ or $`\stackrel{~}{H}=0`$ and satisfy the required boundary condition. We note that there is no such configuration if $`\mathrm{Im}W(a)<\mathrm{Im}W(b)`$. In such a case $`I(a,b)=0`$.
Now let us compute the index. We are considering the situation as depicted in Figure 8 where the arrowed line from $`A`$ to $`B`$ is a straight segment in the negative imaginary direction of the $`W`$-plane. We consider the wave-front at the point $`B`$ along the straight line from $`W(b)`$ and another wave-front at $`B`$ along the broken segment starting from $`W(a)`$ and bending at the point $`A`$. From the general theory of singularities, the two wave-fronts have intersection number $`\mathrm{\Delta }_a\mathrm{\Delta }_b`$, the same as the soliton number between $`a`$ and $`b`$. This means that there are $`\mathrm{\Delta }_a\mathrm{\Delta }_b`$ paths from $`\gamma _a`$ to $`\gamma _b`$ that maps to the straight segment from $`A`$ to $`B`$ in the $`W`$-plane. Since this holds for any starting point $`A`$, there are $`\mathrm{\Delta }_a\mathrm{\Delta }_b`$ families of such paths parametrized by $`w_1:=\mathrm{Re}A=\mathrm{Re}B`$. It may appear that there are infinitely many solutions to (3.31) and therefore infinite degeneracy of supersymmetric ground states. However, we note that the length of $`x^1`$ that is required to go from $`\gamma _a`$ to $`\gamma _b`$ depends on each path and does not necessarily coincide with $`\pi `$. The required length of $`x^1`$ for each path $`P`$ is given by
$$\mathrm{\Delta }x^1=\left|2\underset{P}{}\frac{\mathrm{dIm}W}{|W|^2}\right|.$$
(3.32)
Only the path with $`\mathrm{\Delta }x^1=\pi `$ defines a classical supersymmetric ground state. If the starting point $`A`$ or the end point $`B`$ is the critical value $`W(a)`$ or $`W(b)`$, the required length is infinity $`\mathrm{\Delta }x^1=+\mathrm{}`$. In the massive theory where the bosonic potential $`U=|W|^2`$ diverges at infinity, $`\mathrm{\Delta }x^1`$ approaches zero when $`w_1=\mathrm{Re}A`$ goes to infinity. Thus, for each of the $`\mathrm{\Delta }_a\mathrm{\Delta }_b`$ families, $`\mathrm{\Delta }x^1`$ is roughly a decreasing function as a function of $`w_1`$. If it is a monotonic function, the function (3.32) cut through $`\mathrm{\Delta }x^1=\pi `$ exactly once and hence the contribution to the index of that family is $`1`$. However, one may encounter a family where it cuts through $`\mathrm{\Delta }x^1=\pi `$ more than once as depicted in Figure 9.
In such a case, we make use of the fact that the index is invariant under the deformation of the theory. In particular, we can rescale the superpotential as $`W\mathrm{e}^tW`$. This changes the function (3.32) as $`\mathrm{\Delta }x^1\mathrm{e}^t\mathrm{\Delta }x^1`$. For an appropriate choice of $`\mathrm{e}^t`$ one can make $`\mathrm{e}^t\mathrm{\Delta }x^1`$ to cut through $`\pi `$ exactly once. Thus, in any case, the contribution to the Witten index is $`1`$ for each family. Thus, the total index is given by
$$I(a,b)=\{\begin{array}{cc}\mathrm{\Delta }_a\mathrm{\Delta }_b\hfill & \text{if}\mathrm{Im}W(a)>\mathrm{Im}W(b),\hfill \\ 0\hfill & \text{if}\mathrm{Im}W(a)<\mathrm{Im}W(b).\hfill \end{array}$$
(3.33)
In the case $`a=b`$ there is of course a unique classical supersymmetric configuration: $`\varphi ^I(x^1)=a^I`$ (constant along $`x^1`$). Thus, we have
$$I(a,a)=1.$$
(3.34)
It is easy to generalize the above analysis to the case where there are critical values between $`W(\gamma _a)`$ and $`W(\gamma _b)`$. Let us consider the simplest case where there is one critical value $`W(c)`$. Then, the number of paths from $`\gamma _a`$ to $`\gamma _b`$ depends on whether the image in the $`W`$-plane is on the left or right of $`W(c)`$, i.e. whether $`w_1<\mathrm{Re}W(c)`$ or $`w_1>\mathrm{Re}W(c)`$; we denote these numbers as $`(\mathrm{\Delta }_a\mathrm{\Delta }_b)_<`$ and $`(\mathrm{\Delta }_a\mathrm{\Delta }_b)_>`$ respectively. Some paths on the left smoothly continue to the right of $`W(c)`$. However, some others hit the critical point $`c`$ at $`w_1=\mathrm{Re}W(c)`$ where $`\mathrm{\Delta }x^1`$ blows up to infinity. The length $`\mathrm{\Delta }x^1`$ is bounded from below by a positive value for the paths on the left of $`W(c)`$. Thus, by rescaling the superpotential if necessary, we can have a situation where the solution with $`\mathrm{\Delta }x^1=\pi `$ exists only on the right of $`W(c)`$. Thus, the index is given by $`I(a,b)=(\mathrm{\Delta }_a\mathrm{\Delta }_b)_>`$. It is obvious how to generalize this argument to the case where there are more than one critical values in the region between $`W(\gamma _a)`$ and $`W(\gamma _b)`$. The index is still given by (3.33) where it is understood that $`\mathrm{\Delta }_a\mathrm{\Delta }_b`$ stands for the intersection number of the wavefronts corresponding to the paths in the $`W`$-plane that meet with each other on the right of the right-most critical value. Note that this is the same as the “intersection number” of $`\gamma _a`$ and $`\gamma _b`$ defined in the previous section; namely, the number $`\mathrm{\#}(\gamma _a\gamma _b^{})`$ where $`\gamma _b^{}`$ is obtained by tilting $`\gamma _b`$ with an infinitesimal positive angle against the real axis in the $`W`$-plane.
The asymmetry in $`I(a,b)`$ under the interchange of $`a`$ and $`b`$ has an interesting interpretation, as we will discuss later in this paper. The mirror version of the same asymmetry is discussed in the next subsection, for holomorphic D-branes on sigma models.
One might be interested in exactly how many supersymmetric ground states are there. If the function (3.32) for a family of paths cut through $`\mathrm{\Delta }x^1=\pi `$ exactly once, there is of course one ground state from that family. However, one may find a family where the graph of $`\mathrm{\Delta }x^1`$ looks like Figure 9. In such a case, extra pairs of states may potentially become supersymmetric ground states (though do not contribute to the index). To see whether this is possible or not, we show that, under a certain assumption, the system under consideration is nothing but the supersymmetric quantum mechanics considered in applied to the infinite-dimensional space of paths. In , the system with Hamiltonian $`H=\frac{1}{2}p^2+\frac{1}{2}(h^{}(x))^2+\frac{1}{2}h^{\prime \prime }(x)(\psi \overline{\psi }\overline{\psi }\psi )`$ is considered where $`h(x)`$ is a real valued function. This system possesses supersymmetry generated by $`Q=\sqrt{2}\overline{\psi }(p+ih^{}(x))`$ and $`Q^{}=\sqrt{2}\psi (pih^{}(x))`$ which satisfy $`\{Q,Q^{}\}=4H`$. It is shown that there is a single supersymmetric ground state as long as $`h^{}(x)`$ cuts through $`h^{}(x)=0`$ odd times (no matter how many), but there is none if it cuts through $`h^{}(x)=0`$ even times. The analysis is extended in to the supersymmetric quantum mechanics on a Riemannian manifold deformed by a Morse function $`h`$. In particular, the supersymmetric ground states are realized as cohomology classes of a cochain complex constructed from the critical points of $`h`$ with the grading determined by the Morse index. Let $`\mathrm{\Omega }_{ab}X`$ be the space of paths $`[0,\pi ]X`$ from $`\gamma _a`$ to $`\gamma _b`$, with the boundary condition that the derivatives at $`x^1=0`$ and $`\pi `$ are normal to $`\gamma _a`$ and $`\gamma _b`$ respectively. From the inspection of the supercharges (LABEL:QpmbQpm), we see that the present system is nothing but the supersymmetric quantum mechanics on $`\mathrm{\Omega }_{ab}X`$ deformed in the same way as if there is a function $`h`$ on $`\mathrm{\Omega }_{ab}X`$ such that $`\delta h/\delta \varphi ^i=ig_{i\overline{ȷ}}_1\overline{\varphi }^{\overline{ȷ}}+\frac{1}{2}_iW`$, and $`\delta h/\delta \overline{\varphi }^{\overline{ȷ}}=ig_{i\overline{ȷ}}_1\varphi ^i+\frac{1}{2}_{\overline{ȷ}}\overline{W}`$; in other words
$$\delta h=_0^\pi dx^1\left\{\omega _{IJ}\delta \varphi ^I_1\varphi ^J+\delta \varphi ^I_I\mathrm{Re}W\right\},$$
(3.35)
where $`\omega `$ is the Kähler form of $`X`$, $`\omega =ig_{i\overline{ȷ}}\mathrm{d}z^i\mathrm{d}\overline{z}^{\overline{ȷ}}`$. If we choose a base point $`\varphi _0`$ of (each connected component of) $`\mathrm{\Omega }_{ab}X`$, one can “define” a function
$$h[\varphi ]=\frac{1}{2}\underset{[0,1]\times [0,\pi ]}{}\widehat{\varphi }^{}\omega +_0^\pi dx^1\mathrm{Re}W,$$
(3.36)
where $`\widehat{\varphi }(s,x^1)`$ is a homotopy in $`\mathrm{\Omega }_{ab}X`$ connecting $`\varphi `$ and $`\varphi _0`$, namely a map $`\widehat{\varphi }:[0,1]\times [0,\pi ]X`$ such that $`\widehat{\varphi }(0,x^1)=\varphi _0(x^1)`$, $`\widehat{\varphi }(1,x^1)=\varphi (x^1)`$, and obeying the Dirichlet/Neumann boundary condition at $`x^1=0`$ and $`x^1=\pi `$. Using the fact that $`\gamma _a`$ and $`\gamma _b`$ are Lagrangian and recalling the boundary condition that $`_1\varphi ^I|_\mathrm{\Sigma }`$ is normal to the brane, it is easy to see that $`h[\varphi ]`$ is invariant under a small variation of the homotopy and that it satisfies (3.35) for the variation of $`\varphi `$. For a large change of homotopy, using the fact that the cycles $`\gamma _a`$ and $`\gamma _b`$ are simply connected, we can show that $`h[\varphi ]`$ changes by $`\frac{1}{2}_C\omega `$ where $`C`$ is a two-cycle in $`X`$. We discard this subtlety by focusing our attention only to those cases where $`_C\omega =0`$ for all (compact) two-cycle $`C`$. We also assume $`c_1(X)=0`$. These assumptions hold in the class of models we consider later in this paper. Then, (3.36) is independent of the choice of homotopy and becomes a well-defined function on $`\mathrm{\Omega }_{ab}X`$. Under the assumption $`c_1(X)=0`$, $`U(1)`$ axial R-charge is conserved and the supercharge $`Q`$ has charge $`1`$. Thus, we can define the cochain complex as in graded by the axial R-charge.<sup>2</sup><sup>2</sup>2There is a subtlety for defining the grading (Morse index) from the fact that the Hessian of $`h[\varphi ]`$ has unbounded spectrum. However, one can regularize it, up to an additive constant, by the index of the corresponding Dirac-type operator which is well-defined under the assumption that $`c_1(X)=0`$ (this is related to the conservation of the axial R-charge). We set the ground of the grading so that it is $`0`$ for the critical path which is unique in the family (unlike Figure 9). The coboundary operator is defined by counting the number of instantons connecting different critical points of $`h[\varphi ]`$. We note that an instanton connecting critical paths $`\varphi _1`$ and $`\varphi _2`$ is a configuration $`\varphi (\tau ,x^1)`$ such that $`\varphi (\mathrm{},x^1)=\varphi _1(x^1)`$, $`\varphi (+\mathrm{},x^1)=\varphi _2(x^1)`$ and satisfying the following equation
$$\frac{\varphi ^i}{\tau }=i\frac{\varphi ^i}{x^1}+\frac{1}{2}g^{i\overline{ȷ}}_{\overline{ȷ}}\overline{W}.$$
(3.37)
We denote the corresponding cohomology groups as
$$\mathrm{HF}_W^p(\gamma _a,\gamma _b),$$
(3.38)
where the grading $`p`$ is given by the axial R-charge of the ground states. This is a Landau-Ginzburg generalization of the Floer homology group . Under the rescaling of the superpotential $`W\mathrm{e}^tW`$, the supercharge simply changes as $`Q\mathrm{e}^{\mathrm{\Delta }h}Q\mathrm{e}^{\mathrm{\Delta }h}`$ where $`\mathrm{\Delta }h=(\mathrm{e}^t1)_0^\pi dx^1\mathrm{Re}W`$. Since multiplication by $`\mathrm{e}^{\mathrm{\Delta }h}`$ is well-defined, the cohomology is invariant under this rescaling.<sup>3</sup><sup>3</sup>3Alternatively, as in one may construct a cochain homotopy equivalence of the reduced complexes introduced above. Then, we see that for $`\mathrm{Im}W(a)>\mathrm{Im}W(b)`$
$$\mathrm{HF}_W^p(\gamma _a,\gamma _b)=\{\begin{array}{cc}𝐑^{|\mathrm{\Delta }_a\mathrm{\Delta }_b|}\hfill & p=0,\hfill \\ 0\hfill & p0.\hfill \end{array}$$
(3.39)
In the situation as in Figure 9, we see that pairs of classical supersymmetric ground states are lifted by an instanton effect. Those states have very small (but positive) energies.
It would be interesting to generalize the above consideration to the case where there are two-cycles with $`_C\omega 0`$ and also to the case where $`c_1(X)0`$. In some of the latter cases, we expect that the cohomology (3.38) is not graded by integers, but by some cyclic group.
#### 3.2.2 Sigma Models
The other example we consider is the supersymmetric sigma model on $`X`$ with trivial superpotential $`W=0`$ where the D-branes are wrapped totally on $`X`$. We couple the left and the right boundaries to $`U(1)`$ gauge fields $`A^{(a)}`$ and $`A^{(b)}`$ respectively that define holomorphic line bundles $`E_a`$ and $`E_b`$ on $`X`$. We use the formulation where the boundary term is given by (3.26) and the boundary condition is the standard one (3.25) (where there is no normal direction in the present case). The theory is invariant under B-type supersymmetry generated by $`Q=\overline{Q}_++\overline{Q}_{}`$ and $`Q^{}=Q_++Q_{}`$. Since the boundary term (3.26) includes the time derivatives of the fields, the Noether charges are modified. Thus, the supercharge $`Q`$ is expressed as
$`Q=\sqrt{2}({\displaystyle _0^\pi }\mathrm{d}x^1\{g_{i\overline{ȷ}}(\overline{\psi }_+^{\overline{ȷ}}+\overline{\psi }_{}^{\overline{ȷ}})_0\varphi ^ig_{i\overline{ȷ}}(\overline{\psi }_+^{\overline{ȷ}}\overline{\psi }_{}^{\overline{ȷ}})_1\varphi ^i\}`$ (3.40)
$`+(\overline{\psi }_+^{\overline{ȷ}}+\overline{\psi }_{}^{\overline{ȷ}})A_{\overline{ȷ}}^{(b)}|_{x^1=\pi }(\overline{\psi }_+^{\overline{ȷ}}+\overline{\psi }_{}^{\overline{ȷ}})A_{\overline{ȷ}}^{(a)}|_{x^1=0}).`$
For the purpose of computing the index, we can focus on the zero modes ($`x^1`$-independent modes). Then from the boundary condition, the left and the right fermionic zero modes are related as $`\psi _0^i=\psi _{+0}^i`$ and $`\overline{\psi }_0^{\overline{ı}}=\overline{\psi }_{+0}^{\overline{ı}}`$. We can identify the quantum mechanical Hilbert space as the space of sections on the bundle
$$\left(T^{(0,1)}X\right)E_a^{}E_b,$$
(3.41)
on which the fermionic zero modes act as<sup>4</sup><sup>4</sup>4 Unlike in the closed string case , we do not have the factor $`TX^{(1,0)}`$ in (3.41) nor
$$\frac{1}{\sqrt{2}}g_{i\overline{ȷ}}(\overline{\psi }_{+0}^{\overline{ȷ}}\overline{\psi }_0^{\overline{ȷ}})(/z^i),\frac{1}{\sqrt{2}}(\psi _{+0}^i\psi _0^i)i_{\mathrm{d}z^i},$$
because $`\overline{\psi }_{+0}^{\overline{ȷ}}\overline{\psi }_0^{\overline{ȷ}}=0`$ and $`\psi _{+0}^i\psi _0^i=0`$ from the boundary condition.
$`{\displaystyle \frac{1}{\sqrt{2}}}(\overline{\psi }_{+0}^{\overline{ı}}+\overline{\psi }_0^{\overline{ı}})`$ $``$ $`\mathrm{d}\overline{z}^{\overline{ı}},`$ (3.42)
$`{\displaystyle \frac{1}{\sqrt{2}}}g_{i\overline{ȷ}}(\psi _{+0}^i+\psi _0^i)`$ $``$ $`i_{/\overline{z}^{\overline{ȷ}}}.`$ (3.43)
Then the supercharge $`Q`$ corresponds to the Dolbeault operator on the bundle $`E_a^{}E_b`$:
$$Q2\overline{}_A=2\mathrm{d}\overline{z}^{\overline{ı}}\left(_{\overline{ı}}+A_{\overline{ı}}^{(b)}A_{\overline{ı}}^{(a)}\right).$$
(3.44)
Thus, the Witten index, which is defined as the index of $`Q`$ operator, is equal to the index of this Dolbeault operator. By the standard index theorem, we obtain
$$I(a,b)=\chi (E_a,E_b):=\underset{X}{}\mathrm{ch}(E_a^{}E_b)\mathrm{Td}(X),$$
(3.45)
where $`\mathrm{Td}(X)`$ is the total Todd class of the tangent bundle of $`X`$ which are given by polynomials of the Chern classes (see e.g. ). It is easy to extend this analysis to the case where the bundles $`E_a`$ and $`E_b`$ have higher ranks. The conclusion remains the same as (3.45).
In general, the index (3.45) is not symmetric nor anti-symmetric under the exchange of $`a`$ and $`b`$. This is related by mirror symmetry, as we will discuss later, with the fact noted earlier, that supersymmetric index $`I(a,b)`$ for Lagrangian D-branes in LG models is neither symmetric nor anti-symmetric. However, since odd Todd classes are divisible by the first Chern class of $`X`$ , for a Calabi-Yau manifold $`\mathrm{Td}(X)`$ is a sum of $`4k`$-forms. Under the exchange $`E_a^{}E_bE_b^{}E_a`$ the Chern character changes by sign flip in the $`(4k+2)`$-form components. Thus, for a Calabi-Yau manifold of dimension $`n`$, the index $`I(a,b)`$ is symmetric for even $`n`$ and anti-symmetric for odd $`n`$ under the exchange of $`a`$ and $`b`$.
One can actually obtain an upper bound on the number of supersymmetric ground states, using a technique of section 3 of . The cohomology of operator $`Q`$ is actually invariant under the rescaling $`_1\varphi ^i\mathrm{e}^t_1\varphi ^i`$, since it is done by conjugation by $`\mathrm{e}^{tP}`$ where $`P`$ is an operator counting the number of fermions of combination $`\overline{\psi }_+^{\overline{ȷ}}\overline{\psi }_{}^{\overline{ȷ}}`$. This means that the cohomology group is independent of the width of the strip in the $`x^1`$ direction, as long as it is finite. A zero energy state should remain a zero energy state as we make the strip thinner and thinner. In particular, a ground state must correspond to a ground state of the quantum mechanics of the zero modes. In this way we obtain the upper bound on the number of ground states. The ground states of the quantum mechanics are the cohomology classes of the Dolbeault complex $`\mathrm{\Omega }^{0,p}(X,E_a^{}E_b)`$ with (3.44) as the coboundary operator. Thus, the quantum mechanical ground states are given by the Dolbeault cohomology
$$\mathrm{H}^{0,p}(X,E_a^{}E_b).$$
(3.46)
We note that the vector R-symmetry, for which $`Q`$ has charge $`1`$, is not broken in the bulk nor by the boundary condition. Thus, the grading $`p`$ of the cohomology group (3.46) is the same as the vector R-charge. The group (3.46) gives us an upper bound on the number of supersymmetric ground states, but we do not have a lower bound in general (the argument in section 3 of does not apply here)<sup>5</sup><sup>5</sup>5This is analogous to the situation in the sigma model on a worldsheet without boundary. B-type supercharge yields Dolbeault complex with the coefficient $`TX^{(1,0)}`$ in the zero mode approximation. This indeed has the correct Witten index
$$\underset{p,q}{}(1)^{p+q}dim\mathrm{H}^{0,p}(X,^qTX^{(1,0)})=\pm \chi (X).$$
However, the cohomology group $`_{p,q}\mathrm{H}^{0,p}(X,^qTX^{(1,0)})`$ itself is in general larger than the space of ground states which we know to be $`_{p,q}\mathrm{H}^{q,p}(X)`$, unless $`X`$ is a Calabi-Yau manifold.. However, if the cohomology is non-vanishing only for even $`p`$ (or only for odd $`p`$), the cohomology group (3.46) is indeed the same as the space of supersymmetric ground states. Later in this paper, we will consider a certain set of bundles such that the cohomology (3.46) vanishes except $`p=0`$ and hence it can be identified as the space of ground states.
### 3.3 The Boundary States
Let us consider a Euclidean quantum field theory formulated on a Riemann surface $`\mathrm{\Sigma }`$ with boundary circles. We choose an orientation of each component $`S^1`$ of the boundary and we call it an incoming (resp. outgoing) component if the $`90^{}`$ rotation of the positive tangent vector of $`S^1`$ (with respect to the orientation of $`\mathrm{\Sigma }`$) is an inward (outward) normal vector at the boundary. We choose the metric on $`\mathrm{\Sigma }`$ such that it is a flat cylinder near each boundary component. Suppose $`\mathrm{\Sigma }`$ has a single outgoing boundary, $`S^1=\mathrm{\Sigma }`$. The partition function on $`\mathrm{\Sigma }`$ depend on the boundary condition $`a`$ on the fields at $`\mathrm{\Sigma }`$ and we denote it by $`Z^a(\mathrm{\Sigma })`$. On the other hand, the path-integral over the fields on $`\mathrm{\Sigma }`$ defines a state $`|\mathrm{\Sigma }`$ that belongs to the quantum Hilbert space $`_{S^1}`$ at the boundary circle. We define the boundary state $`a|`$ corresponding to the boundary condition $`a`$ by the property
$$Z^a(\mathrm{\Sigma })=a|\mathrm{\Sigma }.$$
(3.47)
If $`\mathrm{\Sigma }`$ has a single incoming boundary $`\mathrm{\Sigma }=S^1`$, we have a state $`\mathrm{\Sigma }|`$ that belongs to the dual space $`_{S^1}^{}`$. For a boundary condition $`b`$ at $`S^1`$, we define the boundary state $`|b`$ by
$$Z_b(\mathrm{\Sigma })=\mathrm{\Sigma }|b,$$
(3.48)
where $`Z_b(\mathrm{\Sigma })`$ stands for the partition function on $`\mathrm{\Sigma }`$ with the boundary condition $`b`$. In general, the boundary state $`a|`$ (resp. $`|b`$) does not belong to $`_{S^1}^{}`$ (resp. $`_{S^1}`$) but is a formal sum of elements therein. If $`\mathrm{\Sigma }`$ consists of several incoming components $`S_i^1`$ and outgoing components $`S_j^1`$, we have a map $`f_\mathrm{\Sigma }:_i_{S_i^1}_j_{S_j^1}`$. The partition function on $`\mathrm{\Sigma }`$ with the boundary conditions $`\{a_j\}`$ and $`\{b_i\}`$ can be expressed using the boundary states as
$$Z_{\{b_i\}}^{\{a_j\}}(\mathrm{\Sigma })=\left(\underset{j}{}a_j|\right)f_\mathrm{\Sigma }\left(\underset{i}{}|b_i\right).$$
(3.49)
For instance, let us consider a flat finite size cylinder $`\mathrm{\Sigma }`$ of length $`T`$ and circumference $`\beta `$. With a choice of orientation in the circle direction, we have one incoming and one outgoing boundaries. We choose the boundary conditions $`b`$ and $`a`$ there. Then, the partition function is given by $`Z_b^a(\mathrm{\Sigma })=a|\mathrm{e}^{TH(\beta )}|b`$, where $`H(\beta )`$ is the Hamiltonian of the theory on the circle of circumference $`\beta `$.
This is the interpretation of the partition function from the closed string view point. On the other hand, one can interpret it from the point view of open strings. Let $`_{ab}`$ be the space of states on the interval of length $`T`$ with $`a`$ and $`b`$ as the left and the right boundary conditions and let $`H(T)`$ be the Hamiltonian generating the evolution in the circle direction. If the theory has spin half fermions and if the spin structure is periodic (anti-periodic) along the circle direction, the partition function is the trace of $`(1)^F\mathrm{e}^{\beta H(T)}`$ ($`\mathrm{e}^{\beta H(T)}`$) over $`_{ab}`$. Thus, we have
$`\mathrm{Tr}_{_{ab}}(1)^F\mathrm{e}^{\beta H(T)}={}_{_{\mathrm{RR}}}{}^{}a|\mathrm{e}^{TH(\beta )}|b_{_{\mathrm{RR}}}^{},`$ (3.50)
$`\mathrm{Tr}_{_{ab}}\mathrm{e}^{\beta H(T)}={}_{_{\mathrm{NS}^2}}{}^{}a|\mathrm{e}^{TH(\beta )}|b_{_{\mathrm{NS}^2}}^{},`$ (3.51)
where RR (NS<sup>2</sup>) shows that the fermions on the circle are periodic (anti-periodic).
Let us consider a $`(2,2)`$ supersymmetric field theory formulated on the strip $`\mathrm{\Sigma }=𝐑\times I`$ of Minkowski signature, with the boundary conditions $`a`$ and $`b`$ that preserves A- or B-type supersymmetry. We recall that the $`x^1`$ components of the supercurrents are required to obey $`\overline{G}_+^1+G_{}^1=G_+^1+\overline{G}_{}^1=0`$ for A-type supersymmetry and $`\overline{G}_+^1+\overline{G}_{}^1=G_+^1+G_{}^1=0`$ for B-type supersymmetry (for the trivial phases $`\mathrm{e}^{i\alpha }=\mathrm{e}^{i\beta }=1`$). Now let us compactify the time direction $`𝐑`$ to $`S^1`$ and continue the theory to Euclidean signature by the Wick rotation $`x^0=ix^2`$ where we choose the orientation so that $`z=x^1+ix^2`$ is a complex coordinate. The boundary conditions of the supercharges remains the same as in the Minkowski theory. If we change the coordinates as $`(x^1^{},x^2^{})=(x^2,x^1)`$, the conditions become $`\mathrm{e}^{\frac{\pi i}{4}}\overline{G}_+^{}^2^{}+\mathrm{e}^{\frac{\pi i}{4}}G_{^{}}^2^{}=\mathrm{e}^{\frac{\pi i}{4}}G_+^{}^2^{}+\mathrm{e}^{\frac{\pi i}{4}}\overline{G}_{^{}}^2^{}=0`$ for A-type supersymmetry and $`\mathrm{e}^{\frac{\pi i}{4}}\overline{G}_+^{}^2^{}+\mathrm{e}^{\frac{\pi i}{4}}\overline{G}_{^{}}^2^{}=\mathrm{e}^{\frac{\pi i}{4}}G_+^{}^2^{}+\mathrm{e}^{\frac{\pi i}{4}}G_{^{}}^2^{}=0`$ for B-type supersymmetry, where the phases $`\mathrm{e}^{\pm \frac{\pi i}{4}}`$ come from the spin of the supercurrent. This means that the boundary states satisfy
$`\left(\overline{G}_+^{}^2^{}iG_{^{}}^2^{}\right)|b=\left(G_+^{}^2^{}i\overline{G}_{^{}}^2^{}\right)|b=J_A^2^{}|b=0,`$ (3.52)
$`a|\left(\overline{G}_+^{}^2^{}iG_{^{}}^2^{}\right)=a|\left(G_+^{}^2^{}i\overline{G}_{^{}}^2^{}\right)=a|J_A^2^{}=0`$ (3.53)
for A-type supersymmetry and
$`\left(\overline{G}_+^{}^2^{}i\overline{G}_{^{}}^2^{}\right)|b=\left(G_+^{}^2^{}iG_{^{}}^2^{}\right)|b=J_V^2^{}|b=0,`$ (3.54)
$`a|(\overline{G}_+^{}^2^{}i\overline{G}_{^{}}^2^{})=a|\left(G_+^{}^2^{}iG_{^{}}^2^{}\right)=a|J_V^2^{}=0,`$ (3.55)
for B-type supersymmetry. Here we have added the conditions for conservation of the R-charge, which applies when the R-symmetry is not broken in the bulk theory. Note that, in the quantization of the closed strings, the Hermiticity condition is imposed so that $`(G_\pm ^{}^\mu ^{})^{}=\overline{G}_\pm ^{}^\mu ^{}`$ (whereas the quantization of open strings would lead to $`(G_\pm ^\mu )^{}=\overline{G}_\pm ^\mu `$). Thus, the above conditions on the boundary states are not invariant under Hermitian conjugation. If $`|b`$ and $`a|`$ correspond to the boundary conditions preserving A- or B-type supersymmetry with the phase $`\mathrm{e}^{i\alpha }`$ or $`\mathrm{e}^{i\beta }`$, the Hermitian conjugates $`\overline{b}|`$ and $`|\overline{a}`$ correspond to the boundary conditions preserving A- or B-type supersymmetry with the phase $`\mathrm{e}^{i\alpha }`$ or $`\mathrm{e}^{i\beta }`$. If the sign flip $`(1)^{F_L}`$ of the left-moving worldsheet fermions is a symmetry of the theory, the states $`\overline{b}|(1)^{F_L}`$ and $`(1)^{F_L}|\overline{a}`$ correspond to the boundary conditions preserving the A- or B-type supersymmetry with the phase $`\mathrm{e}^{i\alpha }`$ or $`\mathrm{e}^{i\beta }`$, which is the same as the original supersymmetry.<sup>1</sup><sup>1</sup>1Our convention differs from that in the reference where $`b|`$ stands for the Hermitian conjugate of $`|b`$. In particular, we do not need an extra $`(1)^{F_L}`$ insertion in the r.h.s. of eq. (3.56) that is required in the notation of .
As above, let $`a`$ and $`b`$ be the boundary conditions that preserve the same combinations of the supercharges (A-type or B-type). We can use the boundary states to represent the supersymmetric index as
$$I(a,b)={}_{_{\mathrm{RR}}}{}^{}a|\mathrm{e}^{TH(\beta )}|b_{_{\mathrm{RR}}}^{},$$
(3.56)
where $`{}_{_{\mathrm{RR}}}{}^{}a|`$ and $`|b_{_{\mathrm{RR}}}`$ are the boundary states in the RR sector. By the basic property of the index, it is independent of the various parameters, such as $`\beta `$ and $`T`$. It is an integer and therefore must be invariant under the complex conjugation that induces the replacement $`(a,b)(\overline{b},\overline{a})`$. We note, however, that the latter preserves a different combination of the supercharges compared to the original one.
#### 3.3.1 Boundary Entropy
The boundary states are in general a sum of infinitely many eigenstates of the Hamiltonian. An important information on the boundary states can be obtained by looking at the contribution by the ground state. For instance, in boundary conformal field theory, the coefficient $`g_b=0|b`$ of the expansion is known to play a role analogous to that of the central charge $`c`$ in the bulk theory and is called boundary entropy. In supersymmetric field theory, there are several supersymmetric ground states $`|i`$ in the RR sector. Thus the $`𝒩=2`$ analog of the boundary entropy would be the pairings
$$\mathrm{\Pi }_i^a={}_{_{\mathrm{RR}}}{}^{}a|i.$$
(3.57)
These overlaps were studied in especially on the relation to the period integrals, which will be further elaborated here. If the axial R-symmetry is unbroken in the bulk theory, we see from (3.53) that an A-type boundary state $`a|`$ has zero axial charge. Thus for the pairing (3.57) to be non-vanishing, the ground state $`|i`$ must also have zero axial R-charge. Likewise, if vector R-symmetry exists, the pairing (3.57) for B-type boundary state is non-vanishing only for the ground state $`|i`$ with zero vector R-charge. If the theory has a mass gap, this selection rule is vacuous since all ground states have zero R-charges. However, if there is a non-empty IR fixed point, some of the ground states can have non-vanishing R-charges and this selection rule is non-trivial. For example, in LG models all ground states have vanishing R-charges (even if it is quasihomogeneous and has a non-trivial fixed point) and the selection rule is vacuous for A-type boundary states, but in LG orbifold, there are usually ground states of nonzero axial R-charges and the selection rule is non-trivial.
If the vector (resp. axial) R-charge is conserved and integral, there is a one to one correspondence between the supersymmetric ground states and the elements of the $`ac`$ ring (resp. $`cc`$ ring) . The state $`|\varphi _i`$ corresponding to a chiral ring element $`\varphi _i`$ is the one that appears at the boundary $`S^1`$ of the semi-infinite cigar $`\mathrm{\Sigma }`$ with the insertion of $`\varphi _i`$ at the tip, where the theory is twisted to a topological field theory in the curved region.
Thus, for those states, the pairings $`{}_{_{\mathrm{RR}}}{}^{}a|\varphi _i`$ can be identified as the path-integral on the semi-infinite cigar where the boundary condition $`a`$ is imposed at the outgoing boundary and the operator $`\varphi _i`$ is inserted at the tip (see Figure 11).
For concreteness, let us consider a theory with conserved and integral axial R-charge where B-twist is possible. The operators $`\varphi _i`$ we use to define the supersymmetric ground states are the $`cc`$ ring elements. We will be interested in A-type boundary conditions $`a`$ and the corresponding boundary states $`a|`$ obeying (3.53) which in particular yields
$$a|\left(_{S^1}\overline{G}_+i_{S^1}G_{}\right)=a|\left(_{S^1}G_+i_{S^1}\overline{G}_{}\right)=0.$$
(3.58)
Here we use the closed string coordinates as in (3.53) but omit the primes. Also, we use the current notation $`G_\pm =\mathrm{d}x^1G_\pm ^2\mathrm{d}x^2G_\pm ^1`$ which are one forms with values in the spinor bundles of $`\mathrm{\Sigma }`$. After B-twisting, the currents $`\overline{G}_\pm `$ become ordinary one forms but the current $`G_{}`$ (resp. $`G_+`$) becomes a one form with values in the bundle of holomorphic (resp. antiholomorphic) one forms of $`\mathrm{\Sigma }`$.
The pairings $`\mathrm{\Pi }_{a,i}={}_{_{\mathrm{RR}}}{}^{}a|\varphi _i`$ are invariant under the twisted F-term deformations of the theory.
$$\frac{\mathrm{\Pi }_i^a}{t_{ac}}=0,\frac{\mathrm{\Pi }_i^a}{\overline{t}_{ac}}=0.$$
(3.59)
As for the F-term deformations generated by $`cc`$ ring elements, they satisfy the following equation
$$(_i\mathrm{\Pi }^a)_j=(D_i\delta _j^k+i\beta C_{ij}^k)\mathrm{\Pi }_k^a=0,(_{\overline{ı}}\mathrm{\Pi }^a)_{\overline{ȷ}}=(D_{\overline{ı}}\delta _{\overline{ȷ}}^{\overline{k}}i\beta C_{\overline{ı}\overline{ȷ}}^{\overline{k}})\mathrm{\Pi }_{\overline{k}}^a=0,$$
(3.60)
where $`\beta `$ is the circumference of the boundary circle $`S^1`$. Here $`D_i`$ is the covariant derivative defined in and $`C_{ij}^k`$ is the structure constant of the chiral ring. These can be shown as follows by the standard gymnastics in $`tt^{}`$ equation.
We start with the twisted F-term deformation which can be written as <sup>2</sup><sup>2</sup>2Note that $`(Q_+\varphi _{ac})(x)=_xG_+\varphi _{ac}(x)`$ where the contour integral is along a small circle around the point $`x`$.
$$\frac{1}{2}_\mathrm{\Sigma }\overline{Q}_{}Q_+\varphi _{ac}\sqrt{h}\mathrm{d}^2x$$
(3.61)
plus its complex conjugate. Here $`\varphi _{ac}`$ is a twisted chiral operator of axial R-charge $`2`$ (and therefore $`\overline{Q}_{}Q_+\varphi _{ac}`$ has spin zero even in the twisted theory; the spin of $`Q_+`$ cancels that of $`\varphi _{ac}`$). Now, let us divide the semi-infinite cigar into two infinite regions $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_2`$ separated by a circle $`S_{\mathrm{mid}}^1`$ as shown in Figure 12.
We first consider the integral (3.61) in the region $`\mathrm{\Sigma }_1`$. We recall that $`(Q_+\varphi _{ac})(x)=_xG_+\varphi _{ac}(x)`$ etc, where the contour integral is along a small circle around the point $`x`$. We can deform the contour of the $`\overline{G}_{}`$ integral from the small circle around $`x`$ to two boundary circles of $`\mathrm{\Sigma }_1`$; $`S^1`$ and $`S_{\mathrm{mid}}^1`$. The one on $`S_{\mathrm{mid}}^1`$ can be considered as the supercharge acting on the state at the boundary of $`\mathrm{\Sigma }_2`$. Since $`\mathrm{\Sigma }_2`$ is infinitely long, the state that appears at the boundary is a ground state. Thus, the contour integral along $`S_{\mathrm{mid}}^1`$ vanishes. The one on $`S^1`$ turns into $`i_{S^1}G_+`$ by the boundary condition (3.58). By deforming the contour back into the interior, it becomes the sum of the integral on a small circle around $`x`$ and the one on $`S_{\mathrm{mid}}^1`$. The latter vanishes for the same reason as before. The former becomes $`iQ_+Q_+\varphi _{ac}`$ which is zero from the nilpotency of $`Q_+`$. We next consider the region $`\mathrm{\Sigma }_2`$. This has a curved region but one can still deform the contour of the $`\overline{G}_{}`$ integral since it is an ordinary one-form on $`\mathrm{\Sigma }`$. The contour can be deformed to the sum of $`S_{\mathrm{mid}}^1`$ and a small circle around the tip $`x_0`$ at which $`\varphi _j`$ is inserted. The integral on $`S_{\mathrm{mid}}^1`$ can be considered as the supercharge acting on the (dual) state that appears on the incoming boundary of $`\mathrm{\Sigma }_1`$. Since $`\mathrm{\Sigma }_1`$ is infinite, the state is a ground state and thus the supercharge vanishes. The contour integral around $`x_0`$ yields $`\overline{Q}_{}\varphi _j`$ but this vanishes since $`\varphi _j`$ is a chiral operator. To summarize, the pairing $`\mathrm{\Pi }_{i,a}`$ is independent of the twisted F-term deformations.
Next, we consider the F-term deformation generated by the chiral operator $`\varphi _i`$, which can be written as
$$\frac{1}{2}_\mathrm{\Sigma }Q_{}Q_+\varphi _i+\frac{1}{2}_\mathrm{\Sigma }\overline{Q}_+\overline{Q}_{}\overline{\varphi }_{\overline{ı}}\sqrt{h}\mathrm{d}^2xi_{S^1}dx^1(\varphi _i\overline{\varphi }_{\overline{ı}}).$$
(3.62)
Here we have included the (constant) boundary term which is required to set the ground state energy at zero, as found before in this section (see (3.30)). We separate the integral into two regions as before. We first consider the region $`\mathrm{\Sigma }_1`$. Since $`\mathrm{\Sigma }_1`$ is flat, one can treat the currents $`G_\pm `$ as ordinary one forms. The contour of the $`G_{}`$ integral can be deformed to the one on $`S^1`$ and the one on $`S_{\mathrm{mid}}^1`$. The latter vanishes by the same reason as before. By the boundary condition the former turns to $`i_{S^1}\overline{G}_+`$ which can be deformed to an integral around the insertion point of $`\varphi _i`$ (plus an integral over $`S_{\mathrm{mid}}^1`$ that vanishes). This yields the term $`\frac{i}{2}\overline{Q}_+Q_+\varphi _i`$ which is the same as $`\frac{i}{2}\{\overline{Q}_+,Q_+\}\varphi _i`$ since $`\varphi _i`$ is a chiral operator. On the other hand, the same manipulation for $`G_+`$ rather than $`G_{}`$ leads to $`\frac{i}{2}\{Q_{},\overline{Q}_{}\}\varphi _i`$. By taking the average of the two we obtain
$`{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{\Sigma }_1}}Q_{}Q_+\varphi _i`$ $`=`$ $`{\displaystyle _{\mathrm{\Sigma }_1}}\left({\displaystyle \frac{i}{4}}\{\overline{Q}_+,Q_+\}\varphi _i{\displaystyle \frac{i}{4}}\{Q_{},\overline{Q}_{}\}\varphi _i\right)\mathrm{d}^2x=i{\displaystyle _{\mathrm{\Sigma }_1}}H\varphi _i\mathrm{d}^2x`$ (3.63)
$`=`$ $`i{\displaystyle _{\mathrm{\Sigma }_1}}{\displaystyle \frac{}{x^2}}\varphi _i\mathrm{d}^2x=i{\displaystyle _{S^1}}\varphi _idx^1i{\displaystyle _{S_{\mathrm{mid}}^1}}\varphi _idx^1,`$
where we have used the supersymmetry algebra $`\{Q_\pm ,\overline{Q}_\pm \}=2(HP)`$. The first term on the right hand side cancels the boundary term in (3.62). The integrand of the last term is a constant along the circle and hence we obtain $`i\beta \varphi _i`$. One can move the operator $`\varphi _i`$ toward the tip $`x_0`$ where $`\varphi _j`$ is inserted, and this will yield the term $`i\beta C_{ij}^k\varphi _k`$. Next we consider the term $`\frac{1}{2}_{\mathrm{\Sigma }_2}Q_{}Q_+\varphi _i`$. Since we have an infinite cylinder to the left of $`\mathrm{\Sigma }_2`$, by definition, this yields the term $`A_{ij}^k\varphi _k`$ where $`A_{ij}^k\mathrm{d}t^i`$ is the connection form defining the covariant derivative $`D_i`$. Thus, we obtain
$$_i\mathrm{\Pi }_j^a=A_{ij}^k\mathrm{\Pi }_k^ai\beta C_{ij}^k\mathrm{\Pi }_k^a.$$
(3.64)
This is nothing but the first equation in (3.60). The derivation of the second equation is similar.
#### 3.3.2 Period Integral as the Boundary Entropy for LG Models
We study the pairings $`\mathrm{\Pi }_i^a`$ in a LG model on a non-compact Calabi-Yau manifold $`X`$, where the axial R-charge is conserved and integral. We consider the D-brane wrapped on a wave-front trajectory $`\gamma _a`$ emanating from a critical point in the positive real direction. The corresponding boundary condition $`a`$ preserves A-type supersymmetry. As we have shown, the pairings $`\mathrm{\Pi }_i^a=a|\varphi _i`$ are invariant under twisted F-term deformations (3.59). In particular they are invariant under the Kähler deformation and can be studied by taking the large volume limit where the contribution of constant maps dominates. We thus expect the quantum mechanical expression $`\mathrm{\Pi }_i^a=_{\gamma _a}\omega _i`$ where $`\omega _i`$ are the vacuum wave functions corresponding to the chiral fields $`\varphi _i`$. It is known that $`\omega _i`$ are middle dimensional differential forms on $`X`$ which have the right dimension to be integrated over the middle dimensional cycles $`\gamma _a`$. However, we recall that, in addition to the ordinary path integral with the boundary condition $`a`$, we have the following boundary term in the Euclidean action
$$\frac{i}{2}dx^1\left(W\overline{W}\right)$$
(3.65)
which comes from the shift (3.30). This is simply $`i\beta (W\overline{W})/2`$ since the integrand is a constant. Thus, the pairing is given by
$$\mathrm{\Pi }_i^a=\underset{\gamma _a}{}\mathrm{e}^{i\beta (W\overline{W})/2}\omega _i.$$
(3.66)
The vacuum wave forms $`\omega _i`$ are in general difficult to determine. However, a simplification is expected when we make a replacement $`(W,\overline{W})(\lambda W,\overline{\lambda }\overline{W})`$ and take the limit $`\overline{\lambda }0`$ but keeping $`\lambda `$ finite. In this limit (and in the quantum mechanical approximation), the supercharges $`Q_\pm `$, $`\overline{Q}_\pm `$ correspond to the operators
$$\begin{array}{cc}\overline{Q}_+\overline{}\frac{i}{2}\lambda W,\hfill & Q_+\overline{}^{},\hfill \\ \overline{Q}_{}^{}\frac{i}{2}\lambda (\overline{}\overline{W})^{},\hfill & Q_{}.\hfill \end{array}$$
(3.67)
The vacuum wave forms $`\omega _i`$ in the $`\overline{\lambda }0`$ limit must be annihilated by these operators. Reference studies the cohomology of the Dolbeault operator deformed as $`\overline{Q}_+`$ in (3.67). Under suitable assumption about $`X`$<sup>3</sup><sup>3</sup>3 The assumption is that $`X`$ be a Stein space, where ordinary Dolbeault cohomology $`H^{p,q}(X)`$ vanishes except $`q=0`$ where it is isomorphic to the space of holomorphic $`p`$-forms., it was shown that the cohomology of $`\overline{}\frac{i}{2}\lambda W`$ is isomorphic to the cohomology of Kozsul complex given by the operator $`W`$ acting on the holomorphic forms. Furthermore, under the assumption that $`W`$ has a finite number of critical points, the latter cohomology group is non-zero only at middle dimension and is isomorphic to the underlying group of the local ring of $`W`$ which is nothing but the chiral ring of the LG model. Here we use this fact and the arguments in to study the overlap integral (3.66) in the $`\overline{\lambda }0`$ limit. The vacuum wave form $`\omega =\omega _i`$, which in particular defines a $`\overline{Q}_+`$ cohomology class, can be written as
$$\omega =\mathrm{\Omega }+(\overline{}\frac{i}{2}\lambda W)\eta ,$$
(3.68)
where $`\mathrm{\Omega }`$ is a holomorphic $`n`$-form (where $`n`$ is the complex dimension of $`X`$) and $`\eta `$ is an $`(n1)`$-form. It is clear that a holomorphic $`n`$-form $`\mathrm{\Omega }`$ is annihilated by all operators $`Q_\pm `$, $`\overline{Q}_\pm `$ in (3.67). Now, let us evaluate the overlap integral
$`\underset{\overline{\lambda }0}{lim}\mathrm{\Pi }^a`$ $`=`$ $`{\displaystyle \underset{\gamma _a}{}}\mathrm{e}^{i\beta \lambda W/2}\omega ={\displaystyle \underset{\gamma _a}{}}\mathrm{e}^{i\beta \lambda W/2}\left(\mathrm{\Omega }+(\overline{}{\displaystyle \frac{i}{2}}\lambda W)\eta \right)`$ (3.69)
$`=`$ $`{\displaystyle \underset{\gamma _a}{}}\mathrm{e}^{i\beta \lambda W/2}\mathrm{\Omega }+{\displaystyle \underset{\gamma _a}{}}\left\{\mathrm{d}\left(\mathrm{e}^{i\beta \lambda W/2}\eta \right)\mathrm{e}^{i\beta \lambda W/2}\eta \right\}.`$
The total derivative term vanishes under the assumption that the integrand vanishes at infinity of $`\gamma _a`$. Let us focus on the term involving $`\eta `$. Here we use the fact that the vacuum wave form $`\omega `$ must be annihilated by all supercharges, not just by $`\overline{Q}_+`$. In particular it must be annihilated by $`Q_{}`$. Since $`\mathrm{\Omega }`$ is trivially annihilated by $``$, this leads to the condition $`\overline{Q}_+\eta =0`$. Since $``$ and $`\overline{Q}_+`$ anti-commute with each other, this means that $`\eta `$ is annihilated by $`\overline{Q}_+`$. In particular, it can be written as
$$\eta =\mathrm{\Omega }_1+(\overline{}\frac{i}{2}\lambda W)\eta _1,$$
(3.70)
where $`\mathrm{\Omega }_1`$ is a holomorphic $`n`$-form and $`\eta _1`$ is an $`(n1)`$-form. Inserting this expression to (3.69), we obtain $`_{\gamma _a}\mathrm{e}^{i\beta \lambda W/2}(\mathrm{\Omega }+\mathrm{\Omega }_1\eta _1)`$, where again we assumed that the total derivative term vanishes. Since $`\eta `$ has no $`(0,n)`$ component, we can choose $`\eta _1`$ to have no $`(0,n1)`$ component. Continuing this procedure, we finally obtain an expression of $`\mathrm{\Pi }^a`$ as an integral over $`\gamma _a`$ of $`\mathrm{e}^{i\beta \lambda W/2}`$ times a holomorphic $`n`$-form only. The vanishing of the total derivative terms is assured by taking $`\lambda =i`$, since the exponential factor becomes $`\mathrm{e}^{\beta W/2}`$ which quickly vanishes at infinity of $`\gamma _a`$ which extends to real positive directions in the $`W`$-plane. Thus, we obtain
$$\underset{\genfrac{}{}{0pt}{}{\lambda i}{\overline{\lambda }0}}{lim}\mathrm{\Pi }_i^a=\underset{\gamma _a}{}\mathrm{e}^{\beta W/2}\mathrm{\Omega }_i.$$
(3.71)
where $`\mathrm{\Omega }_i`$ is a holomorphic $`n`$-form.<sup>4</sup><sup>4</sup>4 In this derivation we have assumed that $`X`$ is a Stein space. In the cases of interest in this paper we will be dealing this is the case for some examples. But we will also also consider cases where $`X`$ has a non-trivial $`\pi _1`$. In such a case one can repeat the arguments above for the covering space and obtain the same results. Even though we considered the limit $`\overline{\lambda }0`$ in finding this overlap between ground states and D-brane boundary states, in the conformal limit this is unnecessary (as the conformal limit corresponds to taking $`\lambda ,\overline{\lambda }0`$). We will use this result in section 5 in the context of the LG realization of minimal models.
The result obtained here was anticipated in part in and can be viewed as an interpretation of some of the observations there. The argument presented there shows that for flat coordinates, i.e. for a special choice of chiral fields, the period integrals $`\mathrm{\Pi }_i^a`$ given above satisfy the holomorphic part of the flatness equations given in eq. (3.60). The anti-holomorphic part trivializes in the limit $`\overline{\lambda }0`$, and thus we obtain the above result in this limit.
## 4 Brane Creation
As we have already discussed the D-branes preserving the A-type supersymmetry in a LG theory are Lagrangian submanifolds, and their image in the $`W`$-plane correspond to straight lines. The slope of the straight lines depend on which phase combination of A-type supercharges one preserves. In particular for $`Q_A^\alpha =\overline{Q}_++\mathrm{exp}(i\alpha )Q_{}`$ the image in the $`W`$-plane forms an angle $`\alpha `$ relative to the real axis. Moreover D-branes which lead to boundary states with finite overlap with Ramond ground states correspond to D-branes whose image in the $`W`$-plane correspond to straight lines emanating from a critical point.
We consider a LG model of $`n`$ variables which are coordinates of $`𝐂^n`$. Let us assume that the superpotential $`W`$ has $`N`$ isolated critical points $`x_1,\mathrm{},x_N`$. We denote by $`B_\alpha `$ the region in $`𝐂^n`$ on which $`\mathrm{Re}[e^{i\alpha }W]`$ is larger than a fixed large value. Let $`\gamma _i`$ be the wavefront trajectory emanating from the critical point $`x_i`$ along the straight line in the $`W`$-plane with the angle $`\alpha `$ against the real axis. As discussed before $`\gamma _i`$ are the cycles on which the D-branes can wrap without breaking the supersymmetry $`Q_A^\alpha `$. It is known that the cycles $`\gamma _i`$ form a basis of the middle-dimensional homology group $`H_n(𝐂^n,B_\alpha )`$ relative to the boundary $`B_\alpha `$. In other words $`H_n(𝐂^n,B_\alpha )`$ can be viewed as the lattice of charge for the D$`n`$ branes.
Now, let us consider a one parameter family $`\gamma _1(t)`$ of D-branes emanating from a critical point $`x_1`$. Here $`t`$ is a deformation parameter either in the couplings in the superpotential $`W`$, or the angle $`\alpha `$ in the combination of supercharges the D-brane preserves. In such a situation a special thing may happen: The image of the $`\gamma _1(t)`$ brane in the W-plane may pass through a critical value $`W(x_2)`$ at some $`t=t_0`$, so that as we go from $`t_0ϵ`$ to $`t_0+ϵ`$ the position of the critical value relative to the image of the $`D`$-brane on the W-plane, goes from one side to the other. In such a case $`\gamma _1(t_0ϵ)`$ and $`\gamma _1(t_0+ϵ)`$ will label different elements of $`H_n(𝐂^n,B_\alpha )`$, i.e. they will have different D-brane charges. In particular as discussed before,
$$[\gamma _1(t_0ϵ)]=[\gamma _1(t_0+ϵ)](\mathrm{\Delta }_1\mathrm{\Delta }_2)[\gamma _2(t_0+ϵ)].$$
(4.1)
In the context of string theory, charge conservation would imply that we have to have created $`+\mathrm{\Delta }_1\mathrm{\Delta }_2`$ of $`\gamma _2`$ branes in order to guarantee charge conservation during this process.
In the present context the same can be said if we demand continuity of the correlation functions of the 2-dimensional theory. For example if we consider the 2-dimensional theory on a cylinder with one boundary ending on the $`\gamma _1(t)`$, then the continuity of the correlation function with this boundary condition as a function of $`t`$ demands that as we change $`t`$ from $`t_0ϵ`$ to $`t_0+ϵ`$ we would have to also add to the correlation function the correlator involving boundaries on the $`\gamma _2`$ brane with multiplicity factor $`+\mathrm{\Delta }_1\mathrm{\Delta }_2`$. Also, continuity of the overlap with the Ramond ground states in the closed channel will already imply this. Note that $`\mathrm{\Delta }_1\mathrm{\Delta }_2`$ possibly being negative simply means we have an opposite orientation for the $`\gamma _2`$-brane (i.e. the rules of Grassmann integration over the fermions has picked an extra minus sign).
In string theory a similar process was discovered in where again charge conservation leads to creation of new branes.
### 4.1 Monodromy and R-Charges
Now we revisit a result obtained in which relates the number of BPS solitons in 2d theories with $`(2,2)`$ supersymmetry to the R-charges of the Ramond ground states at the conformal point. In particular we show how this result follows very naturally from the realization of D-branes in LG theories, together with the Brane creation discussed above. Our proof will be based on the case of LG theories, though the generalization to arbitrary massive theories should hold, as already shown in .
Consider an LG theory with $`N`$ isolated massive vacua. As discussed before we can associate $`N`$ natural D-branes to these vacua, one for each critical point. The image in the $`W`$-plane is a straight line emanating from the critical point and going to infinity along a line whose slope depends on the combination of A-type supercharges we are preserving. In particular for $`Q_A^\alpha =\overline{Q}_++\mathrm{exp}(i\alpha )Q_{}`$ they make an angle $`\alpha `$ relative to the real axis.
Let us start with $`\alpha =0`$ and order the $`N`$ D-branes according to the lower value for $`Im(W)`$, as depicted in Fig. 14.
Let us further assume that the critical values have a convexity compatible with the ordering of $`Im(W)`$ as shown in the figure. This can be done, by deforming the coefficients of $`W`$ if necessary. As we increase $`\alpha `$ from $`0`$ to $`\pi `$ we rotate the image of branes in the W-plane counter clockwise. As discussed in the previous section, during this process we create new branes. In particular the action of brane creation in the basis of branes emanating from the critical points $`\gamma _i`$ is rather simple: The rotation of branes by $`\pi `$ in the W-plane causes the $`\gamma _i`$ brane to cross all the other $`\gamma _j`$ branes with $`j>i`$ exactly once. Moreover during this crossover it creates $`(\mathrm{\Delta }_i\mathrm{\Delta }_j)`$ new $`\gamma _j`$ branes. This action of rotation of branes by $`\pi `$ is thus realized by an $`N\times N`$ upper triangular matrix with 1 on the diagonal and $`\mathrm{\Delta }_i\mathrm{\Delta }_j`$ for each $`i<j`$. This is denoted by
$$S=1+A,$$
(4.2)
where $`A`$ is the upper triangular matrix of inner product of $`\mathrm{\Delta }`$’s.
Now consider instead going from $`\alpha =0`$ to $`\alpha =\pi `$. In this case for each $`i>j`$ we get $`\mathrm{\Delta }_i\mathrm{\Delta }_j`$ brane creation of $`\gamma _j`$. Thus this action is realized as
$$S^t=1+A^t.$$
(4.3)
Now we consider going around from $`\alpha =0`$ to $`\alpha =2\pi `$. In this case the action on the $`\gamma _i`$ brane basis is given by
$$H=SS^t,$$
(4.4)
where we used here the fact that going from $`\alpha =\pi `$ to $`\alpha =2\pi `$ is the inverse of the action of changing $`\alpha `$ from $`0`$ to $`\pi `$.
Now consider rescaling the superpotential $`W\lambda W`$ as we send $`\lambda 0`$. In this limit we approach a conformal point. For any $`\lambda `$ the monodromy operator $`H`$ we have obtained is the same, because the rescaling of $`W`$ does not affect the relative location of D-branes or their intersection numbers.
Consider the boundary states $`|\gamma _i`$ corresponding to the $`i`$-th D-brane at $`\alpha =0`$. At the conformal point we obtain a new conserved R-charge, which is the fermion number of the right-moving fermions. In particular we have
$$Q_A^\alpha =\mathrm{exp}(i\alpha R)Q_A^0\mathrm{exp}(i\alpha R),$$
(4.5)
where $`R`$ denotes the right moving fermion number charge. Thus the $`H`$ monodromy is realized in the conformal limit as
$$H|\gamma _i=\mathrm{exp}(2\pi iR)|\gamma _i.$$
(4.6)
On the other hand we can go to a basis where the action of $`R`$ is diagonal. Note that since the $`|\gamma _i`$ have invertible overlaps with the Ramond ground states, we can choose linear combination of Ishibashi type states associated to Ramond ground states to represent $`|\gamma _i`$. We thus learn that,
$$\mathrm{Eigenvalues}[SS^t]=\mathrm{Spectrum}[\mathrm{exp}(2\pi iR)]\mathrm{on}\mathrm{Ramond}\mathrm{Ground}\mathrm{states},$$
(4.7)
which is a result of rederived in a purely D-brane language. Note that the choice we have made in the convexity of the critical values is irrelevant for the final result, in that the brane creation was derived precisely based on the continuity of physical correlation functions. The operator $`\mathrm{exp}(2\pi iR)`$ is a physical observable and the structure of brane creation guarantees that for any distribution of critical values, going around the $`W`$ plane by $`2\pi `$ will yield the same operator on the $`\gamma _i`$ brane states.
In the next section, after we discuss minimal models we show that we can make a slightly stronger statement than just equating the eigenvalues of $`SS^t`$ with the spectrum of $`\mathrm{exp}(2\pi iR)`$. Namely we can actually find the change of basis which diagonalizes $`SS^t`$ by considering the overlap of chiral fields with definite $`R`$ charges with the corresponding boundary states. That this should be possible is clear, because the chiral fields provide a basis where $`R`$ acts diagonally.
## 5 D-Branes in $`𝒩=2`$ Minimal Models
In this section, we study D-branes of $`𝒩=2`$ minimal models using their realizations as the infra-red fixed points of Landau-Ginzburg models . We will see that the D-branes of the LG models naturally gives rise to the Cardy states of the minimal models and we will be able to study their properties using purely geometric method. In particular, we will find a beautiful geometric realization of the Verlinde ring for $`SU(2)`$ level $`k`$ Wess-Zumino-Witten models as well as a simple understanding of the $`\tau \frac{1}{\tau }`$ modular transformation matrix $`S_i^j`$. We first review the construction of the D-branes in the minimal models and then see how they are realized as the D-branes in the LG models.
### 5.1 Cardy States, Ishibashi States and $`N=2`$ Minimal models
$`𝒩=2`$ minimal models are unitary $`(2,2)`$ superconformal field theories in two dimensions with central charge $`c=\frac{3k}{k+2}`$, where $`k`$ is a positive integer. They can be viewed as an $`SU(2)/U(1)`$ super-GKO construction at level $`k`$. The superconformal primary fields are labeled by three integers $`(l,m,s)`$ such that
$`l`$ $`=`$ $`0,\mathrm{},k,`$ (5.1)
$`m`$ $`=`$ $`(k+1),\mathrm{},(k+2)(\text{mod}\mathrm{\hspace{0.17em}\hspace{0.17em}2}k+4),`$
$`s`$ $`=`$ $`1,0,1,2(\text{mod}\mathrm{\hspace{0.17em}\hspace{0.17em}4}),`$
with the constraint $`l+m+s0(\text{mod}2)`$ and field identification $`(l,m,s)=(kl,m+k+2,s+2)`$. $`s=0,2`$ in the NS sector and s=$`\pm 1`$ in the Ramond sector. The two different values of $`s`$ denote the GSO parity of various states in the Ramond or NS sector. The conformal weights and the $`U(1)`$ charges of the primary fields are (mod integer),
$$h_{m,s}^l=\frac{l(l+2)m^2}{4(k+2)}+\frac{s^2}{8},q_{m,s}^l=\frac{m}{k+2}\frac{s}{2}.$$
(5.2)
The $`𝒩=2`$ chiral primary states are $`(l,l,0)`$ in the NS sector. The related Ramond states $`(l,l+1,1)`$ can be reached by spectral flow. These models can also be described by the IR fixed point of the LG model with a single chiral superfield $`X`$ with superpotential $`W=X^{k+2}`$ . The chiral primary fields $`X^l`$ correspond to the states $`(l,l,0)`$ and provide a representation of the chiral ring. Note that there are only $`k+1`$ chiral primary fields (as $`l`$ ranges from $`0`$ to $`k`$), which correspond to the $`k+1`$ ground states in the Ramond sector. However there are a total of $`(k+1)(k+2)`$ primary states in the Ramond sector (up to a choice of GSO action $`(1)^F`$).
An A-type boundary state satisfies the following boundary conditions,
$$(L_n\overline{L}_n)|B=0,(J_n\overline{J}_n)|B=0,(G_r^\pm +i\overline{G}_r^{})|B=0.$$
(5.3)
For a rational conformal field theory it was shown in that the boundary states are linear combinations of “Ishibashi states” on which the left and the right generators of the superconformal algebra are linearly related. Ishibashi state corresponding to the primary state $`(l,m,s)`$ is given by
$$|l,m,s=\underset{N}{}|l,m,s;NU\mathrm{\Omega }\overline{|l,m,s;N}.$$
(5.4)
Where $`U`$ is an anti-linear operator acting only on the right moving sector as $`U\overline{𝒪}_nU^1=(1)^{h_𝒪}\overline{𝒪}_n`$, $`\mathrm{\Omega }`$ is the mirror automorphism of the $`𝒩=2`$ algebra and the states $`|l,m,s;N`$ form an orthonormal basis of $`_{l,m,s}`$. The boundary states are particular linear combination of Ishibashi states
$$|l,m,s_{BS}=\alpha \underset{(l^{},m^{},s^{})}{}\frac{S_{l,m,s}^{l^{},m^{},s^{}}}{\sqrt{S_{0,0,0}^{l^{},m^{},s^{}}}}|l^{},m^{},s^{}.$$
(5.5)
Where the constant, $`\alpha `$, is fixed by the condition that the partition function in the open sting channel is integral linear combination of the characters. The summation in the above equation is only over allowed states modulo the field identification and $`S_j^i`$ is the matrix representation of the modular transformation $`\tau \frac{1}{\tau }`$ for the characters $`\chi _{l,m,s}(\tau )=\text{Tr}_{_{l,m,s}}q^{L_0\frac{c}{24}}`$,
$$\chi _{l,m,s}(\frac{1}{\tau })=_{(l^{},m^{},s^{})}S_{l,m,s}^{l^{},m^{},s^{}}\chi _{(l^{},m^{},s^{})}(\tau ),$$
(5.6)
and is given by,
$$S_{l,m,s}^{l^{},m^{},s^{}}=\frac{1}{\sqrt{2}(k+2)}\text{Sin}(\pi \frac{(l+1)(l^{}+1)}{k+2})e^{\frac{i\pi mm^{}}{k+2}}e^{\frac{i\pi ss^{}}{2}}.$$
(5.7)
The above identification of boundary state is motivated mainly by demanding integral expansion in the characters of $`Tr_{\alpha ,\beta }q^{L_0}`$ corresponding to open strings ending on $`\alpha `$ and $`\beta `$ D-branes. The integrality of characters in this sector follows from properties of the Verlinde algebra. We are interested in the Ramond part of the boundary state which can be obtained by restricting the sum in eq. (5.5) to Ramond states only. The properly normalized Ramond part of the boundary state is,
$$|l,m,s_{_{\mathrm{RR}}}=\sqrt{2\sqrt{2}}\underset{(l^{},m^{},s^{})_R}{}\frac{S_{l,m,s}^{l^{},m^{},s^{}}}{\sqrt{S_{0,0,0}^{l^{},m^{},s^{}}}}|l^{},m^{},s^{}.$$
(5.8)
Consider an open string in the $`(a,b)`$ sector. As we have discussed in section 3, the index $`I(a,b)=\mathrm{Tr}_{a,b}(1)^F\mathrm{e}^{\beta H}`$ corresponds in the closed string channel to an overlap in the Ramond sector boundary states $`I(a,b)=\mathrm{Tr}_{a,b}(1)^F={}_{_{\mathrm{RR}}}{}^{}a|b_{_{\mathrm{RR}}}^{}`$. Using the expression (5.8), it is straightforward to compute the index in the $`(a,b)=((l_1,m_1,s_1),(l_2,m_2,s_2))`$ sector. Since the index gets contribution from the Ramond ground states only we have,
$`I(a,b)`$ $`=`$ $`{}_{_{\mathrm{RR}}}{}^{}l_1,m_1,s_1|l_2,m_2,s_2_{_{\mathrm{RR}}}^{}`$
$`=`$ $`2\sqrt{2}{\displaystyle \underset{l=0}{\overset{k}{}}}{\displaystyle \frac{(S_{l_1,m_1,s_1}^{l,l+1,1})^{}S_{l_2,m_2,s_2}^{l,l+1,1}}{S_{0,0,0}^{l,l+1,1}}}l,l+1,1|l,l+1,1`$
$`=`$ $`\frac{2\sqrt{2}i}{\sqrt{2}(k+2)}_{l=0}^k\frac{\text{Sin}(\pi {\scriptscriptstyle \frac{(l+1)(l_1+1)}{k+2}})\text{Sin}\left(\pi {\scriptscriptstyle \frac{(l+1)(l_2+1)}{k+2}}\right)}{\text{Sin}(\pi {\scriptscriptstyle \frac{l+1}{k+2}})}e^{\frac{i\pi (l+1)(m_2m_1+1)}{k+2}}e^{\frac{i\pi (s_2s_1)}{k+2}}`$
$`=`$ $`\frac{2e^{\frac{i\pi (s_2s_1)}{2}}}{k+2}_{l=0}^k\frac{\text{Sin}(\pi {\scriptscriptstyle \frac{(l+1)(l_1+1)}{k+2}})\text{Sin}\left(\pi {\scriptscriptstyle \frac{(l+1)(l_2+1)}{k+2}}\right)\text{Sin}\left(\pi {\scriptscriptstyle \frac{(l+1)(m_2m_1+1)}{k+2}}\right)}{\text{Sin}(\pi {\scriptscriptstyle \frac{l+1}{k+2}})}.`$
Finally using the fact that $`s_2s_1`$ is an even integer in the Ramond sector we get ,
$$I((l_1,m_1,s_1),(l_2,m_2,s_2))=(1)^{\frac{s_2s_1}{2}}N_{l_1,l_2}^{m_2m_1}.$$
(5.9)
Where
$`N_{l_2,l_3}^{l_1}=\frac{2}{k+2}_{l=0}^k\frac{\text{Sin}(\pi {\scriptscriptstyle \frac{(l+1)(l_1+1)}{k+2}})\text{Sin}\left(\pi {\scriptscriptstyle \frac{(l+1)(l_2+1)}{k+2}}\right)\text{Sin}\left(\pi {\scriptscriptstyle \frac{(l+1)(l_3+1)}{k+2}}\right)}{\text{Sin}(\pi {\scriptscriptstyle \frac{l+1}{k+2}})}`$
are the $`SU(2)_k`$ fusion coefficients,
$$N_{l_1,l_2}^{m_2m_1}=\{\begin{array}{cc}1\hfill & \text{if}|l_1l_2|m_2m_1\text{min}\{l_1+l_2,2kl_1l_2\},\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}$$
(5.10)
As we have studied in section 3, the index for a pair of D-branes in the LG model can be identified as the “intersection number” of the corresponding cycles. In the LG realization of the minimal model D-branes, as we will see below, the index (5.9) can indeed be considered as such an “intersection number”.
### 5.2 D-branes in LG description
The Landau Ginzburg description of $`A_{k+1}`$ minimal model consists of a single chiral superfield $`X`$ with superpotential $`W=X^{k+2}`$ . In section 3 we saw that the D-branes in the LG description are preimages of the straight lines in W-plane starting from the critical values which correspond to vanishing cycles in the x-space fibered over the W-plane.
The superpotential $`W=X^{k+2}`$ has a single critical point $`X=0`$ of multiplicity $`k+1`$ with critical value $`w_{}=0`$. If we consider deforming the superpotential by lower powers of $`X`$ we will generically obtain $`k+1`$ isolated and non-degenerate critical points with distinct critical values $`w_i`$. We assume that $`\mathrm{Im}(w_i)`$ are separate from one another. Then as we discussed before we would get $`k+1`$ D-branes, one associated to each of the critical points. Moreover the image of these D-branes are straight lines in the W-plane going to $`+\mathrm{}`$ in the real positive direction. For large values of $`X`$ the lower order terms which deform $`W`$ are irrelevant and the D-branes approach the preimages of the positive real axis $`X^{k+2}𝐑_0𝐂`$, namely
$$X=r\mathrm{exp}\left(\frac{2\pi ni}{k+2}\right),n=0,\mathrm{},k+1,r[0,\mathrm{}).$$
(5.11)
Thus we see that the $`X`$-plane is divided up into $`k+2`$ wedge shaped regions by the $`k+2`$ lines going from the origin to infinity making an angle of $`\frac{2\pi n}{k+2}`$ with the positive real axis, we will denote such a line by $`_n`$.
Any D-brane of the deformed theory is a curve in the $`X`$-plane that will asymptote to a pair of such lines, say $`_{n_1}`$ and $`_{n_2}`$ with $`n_1n_2`$. To see this, we note that the deformed superpotential $`W`$ is approximately quadratic around any (non-degenerate) critical point $`a`$ and the preimage of the straight line emanating from $`W(a)`$ in the $`W`$-plane splits to trajectories of two points (wavefronts) starting from $`a`$. The two wavefronts approach the lines $`_{n_1}`$ and $`_{n_2}`$ as they move away from the critical point. To see that $`n_1n_2`$ it is sufficient to note that the two wavefronts can merge only at a critical point (but the $`(k+1)`$ critical values are assumed to be separate in the imaginary direction). Note that the homology class of the D-branes is completely specified by the choices of the combinations of $`k+2`$ wedges in the x-plane, and that the $`k+1`$ D-branes will be enough to provide a linear basis for the non-trivial cycles (since the sum of all wedges is homologically a trivial contractible cycle). Precisely which combination of $`k+1`$ pairs of asymptote we obtain will depend on which deformation of $`W`$ away from criticality we are considering.
In general the $`k+1`$ D-branes we obtain in this way will not intersect with each other (as their images in the W-plane do not intersect one another). Nevertheless, as discussed in section 3 the index $`I(a,b)=\mathrm{Tr}_{a,b}(1)^F`$ is not in general zero and will depend on the number of solutions to (3.31), that is, how many orthogonal gradient trajectories there are from $`a`$ to $`b`$ D-branes, with a fixed length $`x^1[0,\pi ]`$. This in turn is given by the “intersection number” of the D-branes which is defined as the geometric intersection number where the $`b`$-brane is tilted with a small positive angle in the $`W`$-plane.
We thus see that away from the conformal point there are $`k+1`$ distinct pairs of D-branes, each labeled by an ordered pair of integers ($`n_1,n_2`$) which label the asymptotes that it makes (taking into account the orientation of the D-brane). In particular $`n_1`$ and $`n_2`$ are well defined modulo $`k+2`$ and $`n_1n_2`$. We will label such a D-brane by $`\gamma _{n_1n_2}`$. However there are only $`k+1`$ such pairs for a generically deformed $`W`$. In particular we do not have both branes of the type $`\gamma _{n_1n_2}`$ and $`\gamma _{r_1r_2}`$ with $`n_1<r_1<n_2<r_2`$ for generically deformed $`W`$ as that would have required them to geometrically intersect.
Let us consider two allowed branes $`\gamma _{n_1n_2}`$ and $`\gamma _{r_1r_2}`$. We are interested in computing the Witten index in the oriented open string sector starting from the $`\gamma _{n_1n_2}`$ brane and ending on the $`\gamma _{r_1r_2}`$ brane. Let us denote this by the overlap of the corresponding boundary states, namely $`\gamma _{n_1n_2}|\gamma _{r_1r_2}`$. If none of the $`n_i`$ and $`r_i`$ are equal the branes do not intersect even when one of them is slightly tilted (as noted before in the massive theory the case $`n_1<r_1<n_2<r_2`$ is not allowed) and thus the index is zero. The more subtle case is when one of the $`n_i`$ is equal to one of the $`r_i`$. If they are both equal then we get the Witten index to be $`1`$ as discussed before. Without loss of generality we can order the branes so that $`n_1<n_2`$ and $`r_1<r_2`$ (otherwise the intersection number gets multiplied by a minus sign for each switch of order). Thus there are only four more cases to discuss: $`n_i=r_j`$, for some choice of pair of $`i,j=1,2`$. Let us also assume that $`r_1+r_2>n_1+n_2`$ (by $`r_2=n_1`$ we mean equality mod $`k+2`$, i.e. this is $`r_2=n_1+k+2`$). It turns out that in such cases
$$\gamma _{n_1n_2}|\gamma _{r_1r_2}=1\text{iff}n_i=r_i\text{for some }i$$
(5.12)
and zero otherwise. To see this, as discussed in section 2 and 3 it suffices to consider tilting the slope of the image in the $`W`$-plane of the D-brane corresponding to $`\gamma _{r_1r_2}`$ in the positive direction and seeing if they intersect or not. Tilting the slope in the $`W`$-plane in this case will also correspond to tilting the asymptotes $`r_1,r_2`$ in the positive direction and seeing if they intersect the $`n_1,n_2`$ brane. This leads to the above formula. The case of $`r_1=n_1`$ which leads to intersection number 1 and $`r_1=n_2`$ which leads to intersection number zero is depicted in the Fig. 15.
Now we come to the D-branes at the conformal point. Since the $`k+1`$ D-branes make sense arbitrary close to the conformal point, they survive in the limit of conformal point as well. But here since we have different allowed D-branes at the massive theory, depending on the choice of the deformation polynomials, we learn that all of them survive at the conformal point. Since all pairs $`(n_1,n_2)`$ are realized in terms of a D-brane for some deformation of $`W`$ (which follows from Picard-Lefshetz action discussed before) we learn that all D-branes $`\gamma _{n_1n_2}`$ exist for arbitrary unequal integers $`n_1,n_2`$ defined mod $`k+2`$, which now correspond to exact straight lines in the $`x`$ space along the half-lines given by $`n_1`$ and $`n_2`$, passing through the origin. This gives us a total of $`(k+2)(k+1)`$ D-branes, which are pairwise the same upto orientation at the conformal point. Here we are encountering an interesting effect: The number of D-branes jump as we go from the conformal point to the massive theory.
The fact that we have obtained $`(k+2)(k+1)`$ of such branes at the conformal point is very encouraging as that is exactly the same as the predicted number of Cardy states, as already discussed. Moreover, if we consider the range of parameters where $`0n_1<n_2k+1`$ we see that $`|n_1n_2|\{1,\mathrm{},k+1\}`$ and $`n_1+n_2\{0,\mathrm{}2k+2\}`$. The range of these parameters exactly corresponds to the quantum numbers $`(l,m)`$ labelling the boundary states. Note also that $`s=\pm 1`$ for the Ramond sector boundary states which we are considering. Thus we claim the following identification
$$|n_2n_1|=l+1,n_1+n_2=m,s=\mathrm{sign}(n_2n_1).$$
(5.13)
It follows from (5.13) that $`m+l+s=2n_1=0(\text{mod}\mathrm{\hspace{0.17em}2})`$ as is required. The field identification $`(l,m,s)=(kl,m+k+2,s+2)`$ also has a natural identification as shown in Fig. 16 and relates to the statement that if we change $`n_1n_2`$ and $`n_2n_1+k+2`$ we get the same brane back up to a flip in the orientation (reflected in the shift in $`s`$).
It is convenient to choose $`n_2>n_1`$ (with their differences less than $`k+2`$) in which case we have
$$n_2n_1=l+1,n_2+n_1=m,s=1,$$
(5.14)
which can be solved for $`2n_1`$ and $`2n_2`$ as
$$2n_1=ml1,2n_2=m+l+1.$$
(5.15)
We will denote the D-brane corresponding to the Ramond boundary state $`|l,m,s_{RR}`$ by $`\gamma _{l,m,s}`$. We will now provide further evidence for this identification. Along the way we find a simple geometric interpretation of Verlinde ring for $`SU(2)`$ level $`k`$, as well as certain matrix elements of modular transformations matrix $`S`$.
### 5.3 Geometric Interpretation of Verlinde Algebra
We would like to compute the Witten index at the conformal point for the open string strechted between two D-branes $`\gamma _{n_1n_2}`$ and $`\gamma _{r_1r_2}`$ and reproduce the index formula (5.9). One aspect of the formula is clear. The “intersection number” will not change if we rotate both branes by integral multiples of $`\frac{2\pi }{k+2}`$ which implies that the index will depend on $`m_2m_1`$ but not on the other combination of $`m_1`$ and $`m_2`$ (as $`m_1`$ and $`m_2`$ shift by the same amount under the rotation). Moreover the appearance of $`(1)^{\frac{s_2s_1}{2}}`$ in the intersection is also natural as that correlates with the choice of orientation on the D-branes. So without loss of generality we set $`s_1=s_2=1`$, i.e., as before we choose $`n_2>n_1`$ and $`r_2>r_1`$. Also in checking (5.9) in computing the Verlinde algebra coefficients it suffices to consider the case where $`m_2m_10`$ which is the same case as $`r_1+r_2n_1+n_2`$. With these set, it is now clear what the conditions are for obtaining overlap 1, namely we must have
$$n_1r_1<n_2r_2<n_1+k+2$$
(5.16)
and all the other cases vanish. This is simply the condition that the branes intersect as shown in Fig. 17.
Note that the condition of getting non-vanishing results in the case of equality follows from equation (5.12). Now we use (5.15) to rewrite (5.16) as
$$m_1l_11m_2l_21<m_1+l_1+1m_2+l_2+1<m_1l_1+2k+3$$
(5.17)
These four conditions can also be written as
$$|l_2l_1|m\mathrm{min}[l_1+l_2,2kl_1l_2]$$
(5.18)
where $`m=m_2m_1`$ (to show this and write all inequalities in terms of inequalities with equal signs we used the fact that $`m_2m_1`$ and $`l_2+l_1`$ are equal mod 2). This is precisely the condition for the $`SU(2)`$ level $`k`$ algebra and we have thus derived (5.9) from a purely LG point of view.
### 5.4 Period Integrals and Boundary States
As was discussed in the context of LG models and also in section 3 of this paper there is a natural pairing between the A-model boundary states and B-model chiral fields given by integrating the B-model chiral fields over the cycles representing the A-type boundary states. This kind of pairing was first noticed in and elaborated further in . For the $`A_{k+1}`$ minimal models the chiral primary fields are $`X^l`$ and therefore the inner product of the boundary state $`|l,m,s_{_{\mathrm{RR}}}`$ and the state defined by the B-model chiral field is as discussed in section 3,
$`{}_{_{\mathrm{RR}}}{}^{}l,m,s|X^l^{}`$ $`=`$ $`{\displaystyle _{\gamma _{l,m,s}}}𝑑XX^l^{}e^{W(X)},`$
where the superpotential $`W(X)=X^{k+2}`$. The image of the cycle $`\gamma _{l,m,s}`$ in the W-plane is the positive real axis. Thus from the discussion of the previous section we can see that we can parameterize the curve $`\gamma _{l,m,1}`$ in the following way,
$`\gamma _{l,m,1}:X(t)`$ $`=`$ $`(t)^{\frac{1}{k+2}}e^{\frac{i\pi (ml1)}{k+2}},t[\mathrm{},0],`$
$`=`$ $`t^{\frac{1}{k+2}}e^{\frac{i\pi (m+l+1)}{k+2}},t[0,\mathrm{}].`$
Since $`m+l+10(\text{mod}2)`$, the image in the W-plane of $`\gamma _{l,m,1}`$ is the positive real axis. We have given the curve for $`s=1`$ the curve for $`s=1`$ can be obtained from this by reversing the orientation. With the above parameterization of $`\gamma _{l,m,1}`$,
$`{\displaystyle _{\gamma _{l,m,1}}}𝑑XX^l^{}e^{W(X)}`$ $`=`$ $`\frac{e^{\frac{i\pi (ml1)(l^{}+1)}{k+2}}}{k+2}_+\mathrm{}^0t^{\frac{l^{}+1}{k+2}1}e^t𝑑t+\frac{e^{\frac{i\pi (m+l+1)(l^{}+1)}{k+2}}}{k+2}_0^+\mathrm{}t^{\frac{l^{}+1}{k+2}1}e^t𝑑t`$ (5.20)
$`=`$ $`(\frac{e^{\frac{i\pi (ml1)(l^{}+1)}{k+2}}}{k+2}+\frac{e^{\frac{i\pi (m+l+1)(l^{}+1)}{k+2}}}{k+2})_0^+\mathrm{}t^{\frac{l^{}+1}{k+2}1}e^t𝑑t`$
$`=`$ $`\frac{e^{\frac{i\pi m(l^{}+1)}{k+2}}}{k+2}(2i\text{Sin}(\pi \frac{(l+1)(l^{}+1)}{k+2})\mathrm{\Gamma }(\frac{l^{}+1}{k+2}).`$
To relate the above integral with the modular S-matrix we need to use the normalized operator $`X_{norm}^l^{}`$ ( ) in the above integral. The normalization can be determined by evaluating the matrix element $`\overline{X}^l^{}|X^l^{}`$.
To determine the matrix element $`\overline{X}^l^{}|X^l^{}`$ note that $`_{a_1,a_2}|a_1_{_{\mathrm{RR}}}_{_{\mathrm{RR}}}a_2|S^{a_1a_2}=1`$ when it is sandwitched by the ground states. Here $`|a_1_{_{\mathrm{RR}}}`$ form a basis of the Ramond boundary states and $`S^{a_1a_2}`$ is the inverse of the index matrix,
$$S^{a_1a_2}=(S^1)_{a_1a_2},S_{a_1a_2}=I(a_1,a_2)={}_{_{\mathrm{RR}}}{}^{}a_1|a_2_{_{\mathrm{RR}}}^{}.$$
(5.21)
We choose the basis such that the state $`|a_{_{\mathrm{RR}}}`$ corresponds to the D-brane $`_{a+1}_a`$, where $`a=0,\mathrm{},k`$. With this choice of the basis states we see that <sup>1</sup><sup>1</sup>1$`S_{a_1a_2}`$ is the intersection matrix not to be confused with the modular transformation matrix $`S_{l,m,s}^{l^{},m^{},s^{}}`$ for which the indices will always be written as subscript and superscript.
$`S_{a_1a_2}=\delta _{a_1,a_2}\delta _{a_1+1,a_2},S^{a_1a_2}:=(S^1)_{a_1a_2}=\{\begin{array}{cc}1\hfill & ,a_2a_1\hfill \\ 0\hfill & ,a_2<a_1.\hfill \end{array},`$ (5.24)
which follows from eq. (5.12) by taking into account the reversal of the orientation of the neighboring branes. We insert this complete set of states in the matrix element $`\overline{X}^l^{}|X^l^{}`$,
$$\overline{X}^l^{}|X^l^{}=\underset{a_1,a_2=0}{\overset{k}{}}\overline{X}^l^{}|a_1_{_{\mathrm{RR}}}S_{_{\mathrm{RR}}}^{a_1a_2}a_2|X^l^{}.$$
(5.25)
Using eq. (5.20) we see that
$`{}_{_{\mathrm{RR}}}{}^{}a_2|X^l^{}`$ $`=`$ $`{\displaystyle _{_{a_2+1}}}𝑑XX^l^{}e^{W(X)}{\displaystyle _{_{a_2}}}𝑑XX^l^{}e^{W(X)}`$ (5.26)
$`=`$ $`\frac{1}{k+2}\{e^{\frac{2\pi i(a_2+1)(l^{}+1)}{k+2}}e^{\frac{2\pi ia_2(l^{}+1)}{k+2}}\}\mathrm{\Gamma }(\frac{l^{}+1}{k+2}).`$
To calculate $`\overline{X}^l^{}|a_1_{_{\mathrm{RR}}}`$ we use the fact that, as discussed in section 3,
$$\overline{X}^l^{}|a_1_{_{\mathrm{RR}}}=_{_{\mathrm{RR}}}a_1|(1)^{F_L}|X^l^{}^{}.$$
(5.27)
Using the action
$`(1)^{F_L}|X^l^{}=e^{i\pi (\frac{\widehat{c}}{2}\frac{l^{}}{k+2})}|X^l^{}=ie^{\frac{i\pi (l^{}+1)}{k+2}}|X^l^{}`$
where $`\widehat{c}=\frac{c}{3}=\frac{k}{k+2}`$ we thus obtain
$$\overline{X}^l^{}|a_1_{_{\mathrm{RR}}}=\frac{ie^{\frac{i\pi (l^{}+1)}{k+2}}}{k+2}\{e^{\frac{2\pi i(a_1+1)(l^{}+1)}{k+2}}e^{\frac{2\pi ia_1(l^{}+1)}{k+2}}\}\mathrm{\Gamma }(\frac{l^{}+1}{k+2}).$$
(5.28)
Using eq. (5.26) and eq. (5.28) in eq. (5.25) we see that
$`\overline{X}^l^{}|X^l^{}`$ $`=`$ $`4i\frac{e^{\frac{i\pi (l^{}+1)}{k+2}}}{(k+2)^2}\text{Sin}^2(\pi \frac{l^{}+1}{k+2})\mathrm{\Gamma }(\frac{l^{}+1}{k+2})^2_{a_1,a_2=0}^kS^{a_1a_2}e^{\frac{2\pi i(a_2a_1)(l^{}+1)}{k+2}}`$ (5.29)
$`=`$ $`4i\frac{e^{\frac{i\pi (l^{}+1)}{k+2}}}{(k+2)^2}\text{Sin}^2(\pi \frac{l^{}+1}{k+2})\mathrm{\Gamma }(\frac{l^{}+1}{k+2})^2_{a_1=0}^k_{a_2=a_1}^ke^{\frac{2\pi i(a_2a_1)(l^{}+1)}{k+2}}`$
$`=`$ $`\frac{2}{k+2}\text{Sin}(\pi \frac{l^{}+1}{k+2})\mathrm{\Gamma }(\frac{l^{}+1}{k+2})^2.`$
Thus we see that
$$X_{norm}^l^{}=\sqrt{\frac{k+2}{2\text{Sin}(\frac{\pi (l^{}+1)}{k+2})}}\frac{X^l^{}}{\mathrm{\Gamma }(\frac{l^{}+1}{k+2})}.$$
(5.30)
Using the normalized operator $`X_{norm}^l^{}`$ in eq. (5.20) we get
$$_{\gamma _{l,m,s}}𝑑XX_{norm}^l^{}e^{W(X)}=\sqrt{\frac{2}{(k+2)\text{Sin}(\pi \frac{l^{}+1}{k+2})}}e^{i\pi s/2}e^{i\pi \frac{m(l^{}+1)}{k+2}}\text{Sin}(\pi \frac{(l+1)(l^{}+1)}{k+2}).$$
We can immediately recognize the r.h.s in eq. (5.4) as the coefficient of the Ishibashi state in the expansion of the boundary state i.e.,
$${}_{_{\mathrm{RR}}}{}^{}l,m,s|X^l^{}=\sqrt{2\sqrt{2}}\frac{S_{l,m,s}^{l^{},l^{}+1,1}}{\sqrt{S_{0,0,0}^{l^{},l^{}+1,1}}}$$
(5.31)
Thus we have found a beautiful realization of the modular transformation matrix in terms of classical integrals in the LG theory.<sup>2</sup><sup>2</sup>2Computation of boundary entropy in terms of kinks was carried out in in a slightly different context where the modular S-matrix for $`SU(1)_k`$ appeared in a similar way. It would be interesting to see whether and how it is related to the present discussion.
We can actually check more. Namely we know that the Ramond states corresponding to chiral fields $`|X^l^{}`$ provides a basis where the $`R`$ charge is diagonal. This implies that if we consider a basis for the D-branes, for example the one given above, $`\gamma _{n,n+1}:=_{n+1}_n`$ where $`n=0,\mathrm{},k`$ and compute the operator $`SS^t`$ where $`S`$ is intersection matrix given in eq. (5.24), then the corresponding change of basis to make it diagonal should be given by the matrix
$$M_{ab}:={}_{_{\mathrm{RR}}}{}^{}\gamma _{a,a+1}|X^b$$
(5.32)
where $`a,b\{0,\mathrm{},k\}`$. To show this we will calculate $`D:=M^1SS^tM`$ and show that it is a diagonal matrix with eigenvalues equal to the $`R`$ charge.
From eq. (5.4) and eq. (5.31) it follows that the matrix $`M`$ and its inverse is given by,
$`M_{ab}`$ $`=`$ $`i\sqrt{\frac{2}{k+2}}e^{\frac{i\pi (2a+1)(b+1)}{k+2}}\sqrt{\text{Sin}(\pi \frac{b+1}{k+2})},`$ (5.33)
$`(M^1)_{ab}`$ $`=`$ $`\sqrt{\frac{2}{k+2}}{\displaystyle \underset{c=0}{\overset{k}{}}}S^{cb}e^{\frac{2\pi ic(a+1)}{k+2}}\sqrt{\text{Sin}(\pi \frac{a+1}{k+2})}.`$
Now consider $`D_{ab}`$,
$`D_{ab}`$ $`=`$ $`{\displaystyle \underset{c,d,g=0}{\overset{k}{}}}(M^1)_{ac}S_{cd}(S^t)_{dg}M_{gb}`$ (5.34)
$`=`$ $`\frac{2i}{k+2}\sqrt{\text{Sin}(\pi \frac{a+1}{k+2})\text{Sin}(\pi \frac{b+1}{k+2})}_{g,f=0}^ke^{\frac{2\pi if(a+1)}{k+2}}(S^t)_{fg}e^{\frac{i\pi (2g+1)(b+1)}{k+2}},`$
$`=`$ $`\frac{2i}{k+2}\sqrt{\text{Sin}(\pi \frac{a+1}{k+2})\text{Sin}(\pi \frac{b+1}{k+2})}_{g=0}^k_{fg}^ke^{\frac{2\pi if(a+1)}{k+2}}e^{\frac{i\pi (2g+1)(b+1)}{k+2}},`$
$`=`$ $`\frac{2i}{k+2}\sqrt{\text{Sin}(\pi \frac{a+1}{k+2})\text{Sin}(\pi \frac{b+1}{k+2})}e^{\frac{i\pi (b+1)}{k+2}}_{g=0}^k_{fg}^ke^{\frac{2\pi if(a+1)}{k+2}}e^{\frac{2\pi ig(b+1)}{k+2}}.`$
Using the identity
$$\underset{e=0}{\overset{k}{}}\underset{fe}{\overset{k}{}}e^{\frac{2\pi if(a+1)}{k+2}}e^{\frac{2\pi ie(b+1)}{k+2}}=\delta _{a,b}\frac{(k+2)e^{\frac{i\pi (b+1)}{k+2}}}{2i\text{Sin}(\pi \frac{b+1}{k+2})}$$
(5.35)
we see that
$$D_{ab}=e^{\frac{2\pi i(b+1)}{k+2}}\delta _{a,b},$$
(5.36)
which is indeed the spectrum of $`\mathrm{exp}(2\pi iR)`$ for the $`𝒩=2`$ minimal model. This is morally the analog of the fact that in rational conformal field theory the modular transformation matrix corresponding to $`\tau \frac{1}{\tau }`$ diagonalizes the fusion algebra $`N_{ij}^k`$ . Namely in this case the intersection matrix $`S`$ is related to the $`N_{ij}^k`$ coefficients, as already shown, and the $`M`$ is given by the overlap of Ishibashi states with chiral fields eq. (5.31) which is given in terms of the modular transformation matrix of the rational conformal theory.
## 6 Boundary Linear Sigma Models
In this section, we study $`(2,2)`$ supersymmetric gauge theories formulated on a worldsheet with boundary. We seek for boundary conditions that preserve B-type supersymmetry and study its relation to non-linear sigma model to which the theory reduces at low energies. We also analyze how these boundary conditions are described in the dual description that was found in . We include a brief review of the analysis of .
### 6.1 Supersymmetric Boundary Conditions
Let us consider a supersymmetric $`U(1)`$ gauge theory with chiral multiplets $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_N`$ of charge $`Q_1,\mathrm{},Q_N`$. We formulate the theory on the strip $`\mathrm{\Sigma }=𝐑\times I`$ where $`I`$ is a finite interval parametrized by $`x^1[0,\pi ]`$ and $`𝐑`$ is parametrized by the time coordinate $`x^0`$. The action of the system is given by
$$S=\frac{1}{2\pi }\underset{\mathrm{\Sigma }}{}\left(L_{\mathrm{𝑘𝑖𝑛}}+L_{\mathrm{𝑔𝑎𝑢𝑔𝑒}}+L_{\mathrm{𝐹𝐼},\theta }\right)\mathrm{d}^2x.$$
(6.1)
The terms in the integrand are respectively the matter kinetic term, gauge kinetic term and the Fayet-Iliopoulos-Theta term, which are given by
$`L_{\mathrm{𝑘𝑖𝑛}}=D^\mu \overline{\varphi }D_\mu \varphi +{\displaystyle \frac{i}{2}}\overline{\psi }_{}(\stackrel{}{D}\stackrel{}{}_0+\stackrel{}{D}\stackrel{}{}_1)\psi _{}+{\displaystyle \frac{i}{2}}\overline{\psi }_+(\stackrel{}{D}\stackrel{}{}_0\stackrel{}{D}\stackrel{}{}_1)\psi _++D|\varphi |^2+|F|^2`$
$`|\sigma |^2|\varphi |^2\overline{\psi }_{}\sigma \psi _+\overline{\psi }_+\overline{\sigma }\psi _{}i\overline{\varphi }\lambda _{}\psi _++i\overline{\varphi }\lambda _+\psi _{}+i\overline{\psi }_+\overline{\lambda }_{}\varphi i\overline{\psi }_{}\overline{\lambda }_+\varphi ,`$
(6.2)
$`L_{\mathrm{𝑔𝑎𝑢𝑔𝑒}}={\displaystyle \frac{1}{2e^2}}\left[^\mu \overline{\sigma }_\mu \sigma +{\displaystyle \frac{i}{2}}\overline{\lambda }_{}(\stackrel{}{}\stackrel{}{}_0+\stackrel{}{}\stackrel{}{}_1)\lambda _{}+{\displaystyle \frac{i}{2}}\overline{\lambda }_+(\stackrel{}{}\stackrel{}{}_0\stackrel{}{}\stackrel{}{}_1)\lambda _++v_{01}^2+D^2\right],`$
(6.3)
$`L_{\mathrm{𝐹𝐼},\theta }=rD+\theta v_{01}.`$ (6.4)
In the above expressions, the notation $`\overline{\psi }\stackrel{}{D}\stackrel{}{}_\mu \psi =\overline{\psi }(D_\mu \psi )(D_\mu \overline{\psi })\psi `$ is used. Also, we have written the Lagrangian only in the case of single matter field of unit charge ($`N=1`$, $`Q_1=1`$) to avoid complicated expressions, but the generalization is obvious.
If there were no boundary of $`\mathrm{\Sigma }`$, the system would be invariant under $`(2,2)`$ supersymmetry whose transformation laws are given by
$`\delta v_\pm =i\overline{ϵ}_\pm \lambda _\pm +iϵ_\pm \overline{\lambda }_\pm ,`$ (6.5)
$`\delta \sigma =i\overline{ϵ}_+\lambda _{}iϵ_{}\overline{\lambda }_+,`$
$`\delta D={\displaystyle \frac{1}{2}}\left(\overline{ϵ}_+(_0_1)\lambda _+\overline{ϵ}_{}(_0+_1)\lambda _{}+ϵ_+(_0_1)\overline{\lambda }_++ϵ_{}(_0+_1)\overline{\lambda }_{}\right).`$
$`\delta \lambda _+=iϵ_+(D+iv_{01})+ϵ_{}(_0+_1)\overline{\sigma },`$
$`\delta \lambda _{}=iϵ_{}(Div_{01})+ϵ_+(_0_1)\sigma ,`$
and
$`\delta \varphi =ϵ_+\psi _{}ϵ_{}\psi _+,`$
$`\delta \psi _+=i\overline{ϵ}_{}(D_0+D_1)\varphi +ϵ_+F\overline{ϵ}_+\overline{\sigma }\varphi ,`$
$`\delta \psi _{}=i\overline{ϵ}_+(D_0D_1)\varphi +ϵ_{}F+\overline{ϵ}_{}\sigma \varphi ,`$
$`\delta F=i\overline{ϵ}_+(D_0D_1)\psi _+i\overline{ϵ}_{}(D_0+D_1)\psi _{}+(\overline{ϵ}_+\overline{\sigma }\psi _{}+\overline{ϵ}_{}\sigma \psi _+)+i(\overline{ϵ}_{}\overline{\lambda }_+\overline{ϵ}_+\overline{\lambda }_{})\varphi .`$
In the case where $`\mathrm{\Sigma }`$ has a boundary (where we now consider the strip $`\mathrm{\Sigma }=𝐑\times I`$), the action transforms under (6.5) and (LABEL:chiSUSY) as
$`\delta S={\displaystyle \frac{1}{4\pi }}{\displaystyle \underset{\mathrm{\Sigma }}{}}\mathrm{d}x^0\{`$ $`ϵ_+\left[T\overline{\lambda }_+(D_0+D_1)\overline{\varphi }\psi _{}+i\overline{\varphi }\sigma \psi _++{\displaystyle \frac{i}{2e^2}}\overline{\lambda }_{}(_0+_1)\sigma +i\overline{\psi }_+F\right]`$
$`+ϵ_{}\left[\overline{T}\overline{\lambda }_{}(D_0D_1)\overline{\varphi }\psi _++i\overline{\varphi }\overline{\sigma }\psi _{}{\displaystyle \frac{i}{2e^2}}\overline{\lambda }_+(_0_1)\overline{\sigma }i\overline{\psi }_{}F\right]`$
$`+\overline{ϵ}_+\left[\overline{T}\lambda _++\overline{\psi }_{}(D_0+D_1)\varphi +i\overline{\psi }_+\overline{\sigma }\varphi +{\displaystyle \frac{i}{2e^2}}(_0+_1)\overline{\sigma }\lambda _{}+i\overline{F}\psi _+\right]`$
$`+\overline{ϵ}_{}[T\lambda _{}+\overline{\psi }_+(D_0D_1)\varphi +i\overline{\psi }_{}\sigma \varphi {\displaystyle \frac{i}{2e^2}}(_0_1)\sigma \lambda _+i\overline{F}\psi _{}]\},`$
where $`T`$ is defined by
$$T=\left(r|\varphi |^2\frac{D}{2e^2}\right)i\left(\theta +\frac{v_{01}}{2e^2}\right).$$
(6.8)
We look for a boundary condition such that B-type supersymmetry generated by $`Q=\overline{Q}_++\mathrm{e}^{i\beta }\overline{Q}_{}`$ and $`Q^{}`$ is unbroken.
Before discussing the detail, we note that the locality of equation motion for the gauge fields requires the boundary condition
$$\frac{v_{01}}{e^2}=\theta ,\text{at}\mathrm{\Sigma }.$$
(6.9)
Also, the auxiliary fields are solved by $`F=0`$ and
$$\frac{D}{e^2}=r|\varphi |^2.$$
(6.10)
If we use these relations we have $`T=\frac{1}{2}(r|\varphi |^2i\theta )=\frac{1}{2e^2}(D+iv_{01})`$ at $`\mathrm{\Sigma }`$.
We also make an important remark. Since $`v_{01}=_0v_1_1v_0`$ is a total derivative, the Theta term of the action can naïvely be written as the boundary term
$$\frac{\theta }{2\pi }\underset{\mathrm{\Sigma }}{}v_{01}\mathrm{d}^2x\stackrel{\mathrm{?}}{=}\frac{\theta }{2\pi }\underset{\mathrm{\Sigma }}{}v_0dx^0.$$
(6.11)
However, $`v_0`$ is not gauge invariant whereas $`v_{01}`$ is. In particular, when the boundary components are compactified on circles, the right hand side changes by integer multiples of $`\theta `$ under gauge transformations. Thus, for a generic $`\theta `$, (6.11) is not an allowed thing to do. If $`\theta `$ is an integer multiple of $`2\pi `$, however, the right hand side of (6.11) is gauge invariant up to $`2\pi `$ shifts so that $`\mathrm{exp}\left(i\frac{\theta }{2\pi }v\right)`$ is well-defined. Thus, only for those cases, the manipulation (6.11) is allowed. More generally, for a general $`\theta `$ one can shift $`\theta \theta 2\pi n`$, with $`n`$ integer, provided the boundary term $`n_\mathrm{\Sigma }v_0dx^0`$ is added to the action.
### Pure Maxwell Theory
We start the study with the simplest case; without matter. In this case, the theory has a single twisted chiral (gauge) multiplet $`\mathrm{\Sigma }`$ with the twisted superpotential
$$\stackrel{~}{W}=t\mathrm{\Sigma },$$
(6.12)
where $`t`$ is the complex combination of the FI and Theta parameters
$$t=ri\theta .$$
(6.13)
B-type boundary condition for twisted chiral multiplet fields is like A-type boundary condition for chiral multiplet fields. In particular, the world sheet boundary must end on a middle dimensional Lagrangian submanifold whose image in the $`\stackrel{~}{W}`$-plane is a straight line. Since (6.12) is linear in $`\mathrm{\Sigma }`$, this means that the worldsheet boundary must end on a straight line in $`\mathrm{\Sigma }`$. The Lagrangian condition is trivially satisfied. Thus, if we denote the phase of the FI-Theta parameter $`t`$ as
$$t=|t|\mathrm{e}^{i\gamma },$$
(6.14)
the B-type supersymmetric D-brane is the straight line in the $`\sigma `$-plane whose slope is given by $`\gamma `$;
$$\mathrm{Im}\left(\mathrm{e}^{i\gamma }\sigma \right)=\mathrm{constant}\text{at}\mathrm{\Sigma }.$$
(6.15)
The boundary condition on the component fields is given by
$$\begin{array}{c}\mathrm{e}^{i\gamma }(_0+_1)\sigma =\mathrm{e}^{i\gamma }(_0_1)\overline{\sigma },\hfill \\ \mathrm{e}^{i\gamma }\lambda _++\mathrm{e}^{i\gamma }\lambda _{}=0,\hfill \\ \mathrm{e}^{i\gamma }\overline{\lambda }_++\mathrm{e}^{i\gamma }\overline{\lambda }_{}=0,\hfill \end{array}\text{at}\mathrm{\Sigma }.$$
(6.16)
It is indeed easy to check that the variation (LABEL:vari) vanishes for B-type supersymmetry with $`ϵ_{}=ϵ_+`$. For supersymmetry with $`ϵ_{}=\mathrm{e}^{i\beta }ϵ_+`$, we only have to make the replacement $`\sigma \mathrm{e}^{i\beta }\sigma `$, $`\lambda _\pm \mathrm{e}^{\pm i\beta /2}\lambda _\pm `$ in the above expressions.
The zero point energies of $`\sigma `$ and $`\lambda _\pm `$ cancel against each other and the vacuum energy of the system comes purely from the gauge field sector. By the equation of motion (or a Gauss law constraint) $`_1v_{01}=0`$, the field strength $`v_{01}`$ is a constant and by the boundary condition (6.9) it is given by $`v_{01}=e^2\theta `$. The vacuum energy is then given by
$$E_0=\pi \frac{e^2|t|^2}{2}.$$
(6.17)
In particular, the supersymmetry is spontaneously broken if $`t0`$ as can be seen also by the supersymmetry transformation of $`\lambda _\pm `$ in (6.5) (where $`D\pm iv_{01}=e^2(ri\theta )`$ by the constraint). All these are the same as the elements of the standard story in the bulk theory.
### The General Case
Let us now consider the case with matters. It is known that, under certain conditions, the bulk theory can be identified as a non-linear sigma model at low enough energies compared to $`e\sqrt{r}`$ (see for example ). The target space $`X`$ is a toric manifold defined as the solution space to $`_{i=1}^NQ_i|\varphi _i|^2=r`$ modded out by the $`U(1)`$ gauge transformations. We look for the boundary conditions corresponding to D-branes wrapping totally on $`X`$ (with or without coupling to gauge fields on $`X`$).
Theta angle in the gauge theory is usually identified as the $`B`$-field. In non-linear sigma models, as we have seen in section 3, $`B`$-field modifies the boundary condition on the coordinate fields as (3.21), from pure Neumann to mixed Dirichlet-Neumann condition. However, in the gauge theory with action (6.1), the condition remains pure Neumann type $`D_1\varphi =0`$ even if we turn on $`\theta `$. Thus, there appears to be a discrepancy between the gauge theory and the non-linear sigma model when formulated on a worldsheet with boundary. This mismatch can be cured by adding the boundary term
$$S_{\mathrm{𝑏𝑜𝑢𝑛𝑑𝑎𝑟𝑦}}=\frac{\theta }{4\pi r}\underset{\mathrm{\Sigma }}{}\left(iD_0\overline{\varphi }\varphi i\overline{\varphi }D_0\varphi \right)dx^0$$
(6.18)
to the action (6.1). Then, the boundary condition required from the locality of equation of motion becomes
$$\mathrm{cos}(\gamma )D_1\varphi i\mathrm{sin}(\gamma )D_0\varphi =0,\text{at}\mathrm{\Sigma },$$
(6.19)
where $`t=ri\theta =|t|\mathrm{e}^{i\gamma }`$. This corresponds to the mixed Dirichlet-Neumann boundary condition of (3.21). We note that the addition of (6.18) also alters the boundary condition (6.9) for the gauge field as
$$\frac{v_{01}}{e^2}=\theta +\frac{|\varphi |^2}{r}\theta .$$
(6.20)
The total action
$$S_{\mathrm{tot}}=S+S_{\mathrm{𝑏𝑜𝑢𝑛𝑑𝑎𝑟𝑦}}$$
(6.21)
transforms under (6.5) and (LABEL:chiSUSY) as
$`\delta S_{\mathrm{tot}}={\displaystyle \frac{1}{4\pi }}{\displaystyle \underset{\mathrm{\Sigma }}{}}\mathrm{d}x^0\{`$ $`ϵ_+\left[\stackrel{~}{T}\overline{\lambda }_+\left((1\frac{2i\theta }{r})D_0+D_1\right)\overline{\varphi }\psi _{}+i\overline{\varphi }\sigma \psi _++{\displaystyle \frac{i}{2e^2}}\overline{\lambda }_{}(_0+_1)\sigma \right]`$
$`+ϵ_{}\left[\overline{\stackrel{~}{T}}\overline{\lambda }_{}\left((1+\frac{2i\theta }{r})D_0D_1\right)\overline{\varphi }\psi _++i\overline{\varphi }\overline{\sigma }\psi _{}{\displaystyle \frac{i}{2e^2}}\overline{\lambda }_+(_0_1)\overline{\sigma }\right]`$
$`+\overline{ϵ}_+\left[\overline{\stackrel{~}{T}}\lambda _++\overline{\psi }_{}\left((1+\frac{2i\theta }{r})D_0+D_1\right)\varphi +i\overline{\psi }_+\overline{\sigma }\varphi +{\displaystyle \frac{i}{2e^2}}(_0+_1)\overline{\sigma }\lambda _{}\right]`$
$`+\overline{ϵ}_{}[\stackrel{~}{T}\lambda _{}+\overline{\psi }_+((1\frac{2i\theta }{r})D_0D_1)\varphi +i\overline{\psi }_{}\sigma \varphi {\displaystyle \frac{i}{2e^2}}(_0_1)\sigma \lambda _+]\}.`$
Here $`\stackrel{~}{T}`$ is given by
$`\stackrel{~}{T}`$ $`=`$ $`\left(r|\varphi |^2{\displaystyle \frac{D}{2e^2}}\right)i\left(\theta \left(1{\displaystyle \frac{|\varphi |^2}{r}}\right)+{\displaystyle \frac{v_{01}}{2e^2}}\right)`$ (6.23)
$`=`$ $`{\displaystyle \frac{r|\varphi |^2}{2r}}\left(ri\theta \right),`$
where we have used (6.10) and the new boundary condition (6.20) in the second equality. Since $`\stackrel{~}{T}`$ is proportional to $`t=ri\theta `$ as in the pure Maxwell theory, it is obvious that $`\mathrm{\Phi }`$ independent part of the variation (LABEL:variation) vanishes for $`ϵ_{}=ϵ_+`$ under the same condition (6.16) as in the Maxwell theory. We are now left with the following terms (for $`ϵ_{}=ϵ_+`$)
$`ϵ_+\left[\left((1\frac{2i\theta }{r})D_0+D_1\right)\overline{\varphi }\psi _{}+\left((1+\frac{2i\theta }{r})D_0D_1\right)\overline{\varphi }\psi _++i\overline{\varphi }(\sigma \psi _+\overline{\sigma }\psi _{})\right]`$
$`+\overline{ϵ}_+\left[\overline{\psi }_{}\left((1+\frac{2i\theta }{r})D_0+D_1\right)\varphi \overline{\psi }_+\left((1\frac{2i\theta }{r})D_0D_1\right)\varphi +i(\overline{\psi }_+\overline{\sigma }\overline{\psi }_{}\sigma )\varphi \right].`$
The non-derivative terms vanish if the straight line of $`\sigma `$ is of the type:
$$\mathrm{Im}\left(\mathrm{e}^{i\gamma }\sigma \right)=0\text{at}\mathrm{\Sigma },$$
(6.24)
and the matter fermions satisfy the boundary condition
$$\begin{array}{c}\mathrm{e}^{i\gamma }\psi _+=\mathrm{e}^{i\gamma }\psi _{},\hfill \\ \mathrm{e}^{i\gamma }\overline{\psi }_+=\mathrm{e}^{i\gamma }\overline{\psi }_{},\hfill \end{array}\text{at}\mathrm{\Sigma }.$$
(6.25)
It is now straightforward to see that the derivative terms also vanish under the boundary conditions (6.19) and (6.25). It is also easy to see that these boundary conditions (including (6.24)) are invariant under the B-type supersymmetry.
To summarize, the total action $`S_{\mathrm{tot}}`$ has B-type supersymmetry with $`ϵ_{}=ϵ_+`$ under the boundary conditions (6.19) and (6.25) for the matter fields and (6.16), (6.24), and (6.20) for the gauge multiplet fields. These conditions reduce to the ordinary mixed Dirichlet-Neumann boundary conditions (3.21) and (3.22) of the non-linear sigma model on $`X`$. To recover the phase, $`ϵ_{}=\mathrm{e}^{i\beta }ϵ_+`$, it is enough to make the replacement $`\psi _\pm \mathrm{e}^{\pm i\beta /2}\psi _\pm `$, $`\sigma \mathrm{e}^{i\beta }\sigma `$ and $`\lambda _\pm \mathrm{e}^{\pm i\beta /2}\lambda _\pm `$.
So far we have been analyzing the boundary condition of the classical theory. There are two important quantum effects of the theory with $`_{i=1}^NQ_i0`$; the running of the FI parameter $`r`$ and the anomaly of the axial $`U(1)`$ R-symmetry. From the running of $`r`$, $`r_0=_{i=1}^NQ_i\mathrm{log}(\mathrm{\Lambda }_{\mathrm{UV}}/\mathrm{\Lambda })`$, the phase $`\mathrm{e}^{i\gamma }=t/|t|`$ which enters in the boundary condition changes along the renormalization group flow. In particular, if $`_{i=1}^NQ_i>0`$ (which corresponds to an asymptotic free sigma model), the “bare phase” becomes trivial $`\mathrm{e}^{i\gamma _0}1`$ in the continuum limit $`\mathrm{\Lambda }_{\mathrm{UV}}/\mathrm{\Lambda }\mathrm{}`$. Also, by the axial anomaly, the axial rotation can be done not just by the replacement $`\psi _\pm \mathrm{e}^{\pm i\beta /2}\psi _\pm `$, $`\sigma \mathrm{e}^{i\beta }\sigma `$ and $`\lambda _\pm \mathrm{e}^{\pm i\beta /2}\lambda _\pm `$ but together with the shift of the Theta angle $`\theta \theta +_{i=1}^NQ_i\beta `$. These effects should be visible in a quantum effective description. Here we look at the effective action in terms of $`\mathrm{\Sigma }`$-field whose scalar component is chosen to have large expectation values. This is obtained by integrating out the charged matter fields and is given (for $`Q_i=1`$ case) by
$$\stackrel{~}{W}=N\mathrm{\Sigma }(\mathrm{log}\mathrm{\Sigma }1)t\mathrm{\Sigma }.$$
(6.26)
This yields the following effective FI-Theta parameter
$$t_{\mathrm{𝑒𝑓𝑓}}=t+N\mathrm{log}\mathrm{\Sigma },$$
(6.27)
where the energy scale is set by the value of $`\mathrm{\Sigma }`$. This effective theory is essentially the LG model with the superpotential (6.26)<sup>1</sup><sup>1</sup>1 Strictly speaking, the theory involves a gauge field. However, in the absence of light or tachyonic charged matter field, the effect of the gauge field is simply to create the vacuum energy $`e^2(\mathrm{Im}t_{\mathrm{𝑒𝑓𝑓}})^2/2`$, as the standard auxiliary field does. There is actually a (minor) subtlety; If the theory is formulated on $`𝐑^2`$, the physics is periodic in $`\theta `$ which is identified as the constant electric field (divided by $`e^2`$). This is because of the pair creation of the electron and positron which run away to opposite infinity in the space. However, if the theory is formulated on a strip, $`𝐑\times [0,\pi ]`$, the electron positron pair, even if they are pair-created, can never run away to infinity. Thus, the physics is not periodic in $`\theta `$. which has $`N`$ non-degenerate critical points $`\mathrm{\Sigma }_a=\mathrm{e}^{t/N+2\pi ai/N}`$ ($`a=0,\mathrm{},N1`$). As we have seen, a D-brane preserving the B-type supercharge $`Q=\overline{Q}_++\overline{Q}_{}`$ is the preimage of the straight line in the $`\stackrel{~}{W}`$-plane. The equation is given by
$$\mathrm{Im}\left(\mathrm{e}^{i\gamma _{\mathrm{𝑒𝑓𝑓}}(\sigma )}\sigma \right)=\mathrm{constant},$$
(6.28)
where $`t_{\mathrm{𝑒𝑓𝑓}}N=\mathrm{e}^{i\gamma _{\mathrm{𝑒𝑓𝑓}}}|t_{\mathrm{𝑒𝑓𝑓}}N|`$. If we insist the straight line to pass through a critical value $`\stackrel{~}{W}(\mathrm{\Sigma }_a)N\mathrm{e}^{t/N}`$, the constant in the r.h.s. is of order $`\mathrm{e}^{t/N}`$ and can be considered as the correction to the condition (6.24). It is in general a non-trivial task to find the explicit solution to the straight line equation. However, there is a trivial one if the Theta angle vanishes $`\theta =0`$. In this case $`\sigma =|\sigma |`$ is a solution to the straight line equation with the zero slope $`\mathrm{e}^{i\beta }=1`$. By the axial rotation $`\sigma \mathrm{e}^{i\beta }\sigma `$, $`\lambda _\pm \mathrm{e}^{\pm i\beta /2}\lambda _\pm `$, we obtain the solution $`\sigma =\mathrm{e}^{i\beta }|\sigma |`$ with the slope $`\beta `$. However, we should note that this axial rotation shifts the Theta angle from zero to $`\theta =N\beta `$. Indeed the image of $`\sigma =\mathrm{e}^{i\beta }|\sigma |`$ in the $`\stackrel{~}{W}`$-plane is a straight line only when this shift is made. Thus, we have seen that there is a one parameter family of explicit solutions
$$\begin{array}{c}\sigma =\mathrm{e}^{i\theta /N}|\sigma |,\hfill \\ \mathrm{e}^{i\theta /2N}\lambda _++\mathrm{e}^{i\theta /2N}\lambda _{}=0,\hfill \\ \mathrm{e}^{i\theta /2N}\overline{\lambda }_++\mathrm{e}^{i\theta /2N}\overline{\lambda }_{}=0,\hfill \end{array}\text{at}\mathrm{\Sigma },$$
(6.29)
parametrized by the worldsheet Theta angle $`\theta `$. This preserves the supercharge
$$Q=\overline{Q}_++\mathrm{e}^{i\theta /N}\overline{Q}_{}$$
(6.30)
and $`Q^{}`$. There are of course other solutions (especially those with $`\beta \theta /N`$) but the quantum correction is non-trivial and it is not easy to determine them explicitly. It is easy to extend the above solutions to the general $`Q_i`$’s: replace $`N`$ in these formulae by $`_{i=1}^NQ_i`$.
There is actually a better quantum effective description of the bulk theory found in , using the dual variables $`Y_i`$ of the charged fields $`\mathrm{\Phi }_i`$. Later in this section and further in the next section, we will see how the boundary condition is described in that theory. This will lead to the map of D-branes under mirror symmetry.
### Coupling to Gauge Fields on $`X`$
So far, we have been considering a gauge theory that corresponds to the non-linear sigma model on $`X`$ with a $`B`$-field, but not including coupling the worldsheet boundary to the target space gauge fields. Now it is useful to observe that the $`B`$-field obeying a certain quantization condition can be considered as the curvature of a gauge field $`A_I`$ on $`X`$. In such a case, as noted in section 3, the coupling to $`B`$ field is equal to the boundary coupling to the gauge field $`A_I`$. The quantized $`B`$ field corresponds to the case where the worldsheet Theta angle becomes an integer multiple of $`2\pi `$, $`\theta =2\pi n`$. We now recall that in such a case (and only in such a case) the worldsheet Theta term can be converted into a boundary term (6.11). Then, the total boundary term becomes
$`S_{\mathrm{boundary}}^{}`$ $`=`$ $`{\displaystyle \frac{n}{2r}}{\displaystyle \underset{\mathrm{\Sigma }}{}}\left(iD_0\overline{\varphi }\varphi i\overline{\varphi }D_0\varphi \right)dx^0n{\displaystyle \underset{\mathrm{\Sigma }}{}}v_0dx^0`$ (6.31)
$`=`$ $`{\displaystyle \frac{n}{2r}}{\displaystyle \underset{\mathrm{\Sigma }}{}}\left(i_0\overline{\varphi }\varphi i\overline{\varphi }_0\varphi +2v_0(|\varphi |^2r)\right)dx^0`$
In the sigma model limit $`e\sqrt{r}\mathrm{}`$, the constraint $`|\varphi |^2=r`$ is strictly imposed. Then, the boundary term is given by
$$S_{\mathrm{boundary}}^{}=n\underset{\mathrm{\Sigma }}{}A_I_0\varphi ^I\mathrm{d}x^0$$
(6.32)
where
$$A_I\mathrm{d}\varphi ^I=\frac{i}{2}\frac{{\displaystyle {}_{i=1}{}^{N}\overline{\varphi }_{i}^{}\stackrel{}{\mathrm{d}}\stackrel{}{}\varphi _i}}{{\displaystyle {}_{i=1}{}^{N}Q_{i}^{}|\varphi _i|^2}}.$$
(6.33)
In this expression, we have recovered all the $`N`$ matter fields of charge $`Q_1,\mathrm{},Q_N`$ where the constraint is $`_{i=1}^NQ_i|\varphi _i|^2=r`$.
The gauge field $`A_I`$ in (6.33) is nothing but the hermitian connection of the natural holomorphic line bundle $`𝒪_X(1)`$ on the toric manifold $`X`$ (where $`\varphi _i`$’s represent the sections) with respect to the natural hermitian metric. To see this, let us make a gauge transformation $`\varphi _i\mathrm{e}^{iQ_i\lambda }\varphi _i`$. Then, the gauge field transforms as
$$A_I\mathrm{d}\varphi ^IA_I\mathrm{d}\varphi ^I\mathrm{d}\lambda .$$
(6.34)
This is indeed the transformation property of a connection form of the bundle $`𝒪_X(1)`$. For example, let us consider the simplest case $`X=𝐂\mathrm{P}^1`$ where the gauge theory has two matters $`\mathrm{\Phi }_1`$, $`\mathrm{\Phi }_2`$ of charge $`1`$. In the gauge where $`\varphi _1=1`$ and $`\varphi _2=z`$, the gauge field (6.33) is given by
$$A=\frac{i}{2}\frac{\overline{z}\mathrm{d}zz\mathrm{d}\overline{z}}{1+|z|^2}.$$
(6.35)
This is the gauge field of the line bundle $`𝒪(1)`$ of $`𝐂\mathrm{P}^1`$. Indeed, the first Chern class is represented by the curvature $`\frac{i}{2\pi }i\mathrm{d}A=\frac{i}{2\pi }\mathrm{d}z\mathrm{d}\overline{z}/(1+|z|^2)^2`$ which is the positive unit volume form of $`𝐂\mathrm{P}^1`$.
Thus, we indeed see that the boundary term (6.31) corresponds to the boundary coupling to the natural gauge fields of the bundle $`𝒪_X(n)`$. Here we have to bear in mind that the boundary condition should be given by (6.19)-(6.25) and (6.16)-(6.24)-(6.20) where it is understood that $`\gamma =\mathrm{arg}(r2\pi ni)`$. If we turn on the bulk $`\theta `$-term anew, the angle is given by $`\gamma =\mathrm{arg}(t2\pi ni)`$ where $`t=ri\theta `$.
### Alternative Formulation
In the non-linear sigma model, we have seen that there is an alternative formulation for coupling to target space gauge fields where we do not change the boundary condition but add a fermion bilinear boundary term. This can also be done in the gauge theory. The relevant boundary term for the gauge field of the bundle $`𝒪_X(n)`$ is given by
$$S_{\mathrm{boundary}}^{\prime \prime }=\frac{n}{2r}\underset{\mathrm{\Sigma }}{}\left(iD_0\overline{\varphi }\varphi i\overline{\varphi }D_0\varphi +(\psi _++\psi _{})(\overline{\psi }_++\overline{\psi }_{})(\sigma +\overline{\sigma })|\varphi |^2\right)dx^0.$$
(6.36)
It is straightforward to check that this is by itself invariant under the B-type supersymmetry with $`ϵ_{}=ϵ_+`$. Thus, one can add this to the total action $`S_{\mathrm{tot}}`$ without changing the boundary condition. We note that, as in non-linear sigma models, the equations of motion for the worldsheet fields have boundary contributions in this formulation.
### 6.2 A Review of a Derivation of Mirror Symmetry
We now briefly review the dual description of the gauge theory found in .
Let us consider the $`U(1)`$ gauge theory on $`\mathrm{\Sigma }=𝐑^2`$ with matters of charge $`Q_1,\mathrm{},Q_N`$. The action is given by (6.1) (in the case $`N=1,Q_1=1`$). In the superfield notation, the Lagrangian is expressed as
$$L=\mathrm{d}^4\theta (\underset{i=1}{\overset{N}{}}\overline{\mathrm{\Phi }}_i\mathrm{e}^{2Q_iV}\mathrm{\Phi }_i\frac{1}{2e^2}\overline{\mathrm{\Sigma }}\mathrm{\Sigma })+\frac{1}{2}(\mathrm{d}^2\stackrel{~}{\theta }t\mathrm{\Sigma }+c.c.).$$
(6.37)
If we dualize the phase of the charged chiral superfield $`\mathrm{\Phi }_i`$, we obtain a neutral twisted chiral superfield $`Y_i`$ that is periodic with periodicity $`2\pi i`$, $`Y_iY_i+2\pi i`$. The fields $`Y_i`$ are related to the original charged chiral superfields $`\mathrm{\Phi }_i`$ by
$$Y_i+\overline{Y}_i=2\overline{\mathrm{\Phi }_i}\mathrm{e}^{2Q_iV}\mathrm{\Phi }_i,$$
(6.38)
or in components, $`Y_i=y_i+\sqrt{2}\theta ^+\overline{\chi }_{i+}+\sqrt{2}\overline{\theta }^{}\chi _i+\mathrm{}`$,
$`y_i=\varrho _ii\vartheta _i,\{\begin{array}{c}\varrho _i=|\varphi _i|^2,\hfill \\ _\pm \vartheta _i=\pm 2\left(|\varphi _i|^2(_\pm \phi _i+Q_iv_\pm )+\overline{\psi }_{i\pm }\psi _{i\pm }\right),\hfill \end{array}`$ (6.41)
$`\chi _{i+}=2\overline{\psi }_{i+}\varphi _i,\chi _i=2\overline{\psi }_i\varphi _i,`$ (6.42)
$`\overline{\chi }_{i+}=2\overline{\varphi }_i\psi _{i+},\overline{\chi }_i=2\overline{\varphi }_i\psi _i,`$
where $`\phi _i`$ is the phase of $`\varphi _i`$, $`\varphi _i=|\varphi _i|\mathrm{e}^{i\phi _i}`$.
The fields $`Y_i`$ couple to the gauge field as dynamical Theta angle. Thus, at the level of dualization we have the twisted superpotential $`\stackrel{~}{W}=\mathrm{\Sigma }(_{i=1}^NQ_iY_{i0}t_0)`$ where subscript $`0`$ stands for the bare parameters and fields. The FI parameter runs as $`r_0=b_1\mathrm{log}(\mathrm{\Lambda }_{UV}/\mathrm{\Lambda })`$ with
$$b_1:=\underset{i=1}{\overset{N}{}}Q_i,$$
(6.43)
but one can make the superpotential finite by renormalizing the bare fields $`\varrho _{i0}`$ as $`\varrho _{i0}=\varrho _i+\mathrm{log}(\mathrm{\Lambda }_{UV}/\mu )`$, where $`\mu `$ is the renormalization point. The Kähler metric of the $`y_i`$ variables is given classically by
$$\mathrm{d}s^2=\underset{i=1}{\overset{N}{}}\frac{|\mathrm{d}y_i|^2}{2(2r_0/b_1+y_i+\overline{y}_i)}\frac{b_1}{4r_0}\underset{i=1}{\overset{N}{}}|\mathrm{d}y_i|^2.$$
(6.44)
This superpotential is corrected by instanton effect where the instantons are the vortices of the gauge theory. The correction is of the form $`\mathrm{e}^{Y_i}`$ and the exact twisted superpotential is given by
$$\stackrel{~}{W}=\mathrm{\Sigma }\left(\underset{i=1}{\overset{N}{}}Q_iY_it(\mu )\right)+\underset{i=1}{\overset{N}{}}\mu \mathrm{e}^{Y_i}.$$
(6.45)
In the case $`b_10`$, the FI parameter is renormalized and $`\mathrm{\Lambda }=\mu \mathrm{e}^{t/b_1}`$ is renormalization group invariant, as can be seen also from (6.45) using the shifts of $`Y_i`$’s. In the conformal case $`b_1=0`$, $`t`$ is the dimensionless parameter of the theory, and $`\mu `$ can be simply absorbed by the shifts of $`Y_i`$’s. In what follows, we omit the scale $`\mu `$.
In the sigma model limit $`e\sqrt{r}\mathrm{}`$, the gauge multiplet fields becomes infinitely heavy and can be integrated out. Then, this yields a constraint
$$\underset{i=1}{\overset{N}{}}Q_iY_i=t.$$
(6.46)
Thus, we obtain a theory of $`N`$ periodic fields $`Y_i`$ with one constraint (6.46) which has a twisted superpotential
$$\stackrel{~}{W}=\underset{i=1}{\overset{N}{}}\mathrm{e}^{Y_i}.$$
(6.47)
In other words, we obtain a LG model on $`(𝐂^\times )^{N1}`$. Since the original gauge theory becomes non-linear sigma model on the toric manifold $`X`$ in the limit $`e\sqrt{r}\mathrm{}`$, the above LG model is a dual description of the non-linear sigma model on $`X`$. Since it is described by twisted chiral fields, it is the mirror of the sigma model on $`X`$.
It is easy to find the critical point of the superpotential (6.47) under the constraint (6.46). There are $`b_1`$ critical points $`p_0,\mathrm{},p_{b_11}`$, where at the $`a`$-th critical point $`\mathrm{e}^{y_i}(p_a)=Q_i\mathrm{e}^{t/b_1+2\pi ai/b_1}_{j=1}^NQ_j^{Q_j}`$ with the critical value
$$\stackrel{~}{w}_a=b_1\mathrm{e}^{t/b_1+2\pi ai/b_1}\underset{j=1}{\overset{N}{}}Q_j^{Q_j}.$$
(6.48)
All these are massive vacua at which the $`𝐙_{2b_1}`$ axial R-symmetry is spontaneously broken to $`𝐙_2`$.
### 6.3 D-branes and Mirror Symmetry: First Example
We would like to see how the boundary conditions for the linear sigma model can be described in the quantum effective theory in terms of the dual variables. We consider the model with
$$b_1=\underset{i=1}{\overset{N}{}}Q_i>0,$$
(6.49)
that corresponds to an asymptotic free non-linear sigma model. Since the dual theory is a LG model described in terms of twisted chiral superfields, B-type supersymmetry looks like A-type supersymmetry for chiral superfields. In particular, the worldsheet boundary must end on a middle dimensional Lagrangian submanifold of $`(𝐂^\times )^{N1}`$ that is mapped to a straight line in the $`\stackrel{~}{W}`$-plane.
We focus on the family of boundary conditions (6.29) parametrized by the worldsheet Theta angle $`\theta `$. The boundary conditions on the matter fields are
$$\begin{array}{c}\mathrm{cos}(\gamma _0)D_1\varphi _ii\mathrm{sin}(\gamma _0)D_0\varphi _i=0,\hfill \\ \mathrm{e}^{i\gamma _0+i\theta /2b_1}\psi _{i+}=\mathrm{e}^{i\gamma _0i\theta /2b_1}\psi _i,\hfill \\ \mathrm{e}^{i\gamma _0i\theta /2b_1}\overline{\psi }_{i+}=\mathrm{e}^{i\gamma _0+i\theta /2b_1}\overline{\psi }_i,\hfill \end{array}\text{at}\mathrm{\Sigma }.$$
(6.50)
We note that the phase $`\mathrm{e}^{i\gamma _0}`$ defined by $`r_0i\theta =\mathrm{e}^{i\gamma _0}|r_0i\theta |`$ becomes trivial
$$\gamma _00,$$
(6.51)
in the continuum limit $`\mathrm{\Lambda }_{\mathrm{UV}}\mathrm{}`$ where $`r_0=b_1\mathrm{log}(\mathrm{\Lambda }_{\mathrm{UV}}/\mathrm{\Lambda })\mathrm{}`$. This boundary condition preserves the supercharge
$$Q=\overline{Q}_++\mathrm{e}^{i\theta /b_1}\overline{Q}_{}$$
(6.52)
and its conjugate $`Q^{}`$.
We recall that there is a boundary term in the action
$`S_{\mathrm{𝑏𝑜𝑢𝑛𝑑𝑎𝑟𝑦}}`$ $`=`$ $`{\displaystyle \frac{\theta }{4\pi r_0}}{\displaystyle \underset{\mathrm{\Sigma }}{}}{\displaystyle \underset{i=1}{\overset{N}{}}}\left(iD_0\overline{\varphi }_i\varphi _ii\overline{\varphi }_iD_0\varphi _i\right)\mathrm{d}x^0`$ (6.53)
$`=`$ $`{\displaystyle \frac{\theta }{2\pi r_0}}{\displaystyle \underset{\mathrm{\Sigma }}{}}{\displaystyle \underset{i=1}{\overset{N}{}}}|\varphi _i|^2(_0\phi _i+Q_iv_0)\mathrm{d}x^0.`$
Note that in terms of the renormalized dual fields we have $`|\varphi _i|^2=r_0/b_1+\varrho _i`$. Then, in the continuum limit the boundary term can be written as
$$S_{\mathrm{𝑏𝑜𝑢𝑛𝑑𝑎𝑟𝑦}}=\frac{\theta }{2\pi }\underset{\mathrm{\Sigma }}{}\left(\frac{1}{b_1}\underset{i=1}{\overset{N}{}}_0\phi _i+v_0\right)dx^0.$$
(6.54)
Now the relevant part of the action in the dualization is
$$S_\phi =\frac{1}{2\pi }\underset{\mathrm{\Sigma }}{}\underset{i=1}{\overset{N}{}}r_0^2|\mathrm{d}\phi _i+Q_iv|^2\frac{i\theta }{2\pi }\underset{\mathrm{\Sigma }}{}\left(\frac{1}{b_1}\underset{i=1}{\overset{N}{}}\mathrm{d}\phi _i+v\right)$$
(6.55)
where we consider Euclidean signature (for simplicity) and we ignore the fermionic components which are not essential in this part of the argument. We consider another system involving one-form fields $`_i=_{i\mu }\mathrm{d}x^\mu `$ with the action given by
$$S^{}=\underset{i=1}{\overset{N}{}}[\frac{1}{8\pi r_0^2}\underset{\mathrm{\Sigma }}{}_i_i+\frac{i}{2\pi }\underset{\mathrm{\Sigma }}{}_i(\mathrm{d}\phi _i+Q_iv)]\frac{i\theta }{2\pi }\underset{\mathrm{\Sigma }}{}(\frac{1}{b_1}\underset{i=1}{\overset{N}{}}\mathrm{d}\phi _i+v).$$
(6.56)
We require the boundary condition on the one-form fields $`_i`$ that they vanish against the tangent vectors of the boundary
$$_i|_\mathrm{\Sigma }=0.$$
(6.57)
If we first integrate out the one-form field $`_i`$, we obtain the constraint $`_i=i2r_0^2(\mathrm{d}\phi _i+Q_iv)`$ (which is consistent with the boundary condition (6.50) in the continuum limit) and we obtain the original action (6.55). Instead, if we first integrate out the variables $`\phi _i`$, we obtain the constraint
$$_i=\mathrm{d}\vartheta _i$$
(6.58)
where $`\vartheta _i`$ are periodic variables of period $`2\pi `$. By the boundary condition (6.57), we see that $`\vartheta _i`$ are constants along the boundary of $`\mathrm{\Sigma }`$. By the terms $`i(\theta /2\pi b_1)_\mathrm{\Sigma }d\phi _i`$ in the action (6.56), we see that the constants are
$$\vartheta _i=\theta /b_1\text{at}\mathrm{\Sigma },$$
(6.59)
for all $`i`$. Now, if we plug the constraint (6.58) back into (6.56) we obtain the action
$`S_\vartheta `$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\left[{\displaystyle \frac{1}{8\pi r_0^2}}{\displaystyle \underset{\mathrm{\Sigma }}{}}|\mathrm{d}\vartheta _i|^2+{\displaystyle \frac{i}{2\pi }}{\displaystyle \underset{\mathrm{\Sigma }}{}}d\vartheta _iQ_iv\right]{\displaystyle \frac{i\theta }{2\pi }}{\displaystyle \underset{\mathrm{\Sigma }}{}}v`$ (6.60)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\left[{\displaystyle \frac{1}{8\pi r_0^2}}{\displaystyle \underset{\mathrm{\Sigma }}{}}|\mathrm{d}\vartheta _i|^2{\displaystyle \frac{i}{2\pi }}{\displaystyle \underset{\mathrm{\Sigma }}{}}Q_i\vartheta _idv\right]+{\displaystyle \frac{i}{2\pi }}{\displaystyle \underset{\mathrm{\Sigma }}{}}\left({\displaystyle \underset{i=1}{\overset{N}{}}}Q_i\vartheta _i\theta \right)v,`$
where we have performed partial integrations in the last step. Note that the boundary term in the right hand side vanishes if we use the boundary condition (6.59). Then, the dualization proceeds precisely as in the bulk theory and we will obtain the superpotential (6.45). In the sigma model limit $`e^2\sqrt{r}\mathrm{}`$, after integrating out the $`\mathrm{\Sigma }`$-field, we obtain the constraint (6.46) and the twisted superpotential (6.47).
The boundary condition (6.59) means that $`\mathrm{e}^{y_i}`$’s at the boundary have a fixed common phase $`\theta /b_1`$. Namely, the worldsheet boundary $`\mathrm{\Sigma }`$ is mapped by $`(\mathrm{e}^{y_i})`$ to a real $`(N1)`$-dimensional cycle $`\gamma _\theta `$ in the algebraic torus $`(𝐂^\times )^{N1}`$ defined by
$$(\mathrm{e}^{y_1},\mathrm{},\mathrm{e}^{y_N})=(\mathrm{e}^{\varrho _1+i\theta /b_1},\mathrm{},\mathrm{e}^{\varrho _N+i\theta /b_1}),$$
(6.61)
where $`(\varrho _1,\mathrm{},\varrho _N)`$ are the real coordinates constrained by $`_{i=1}^NQ_i\varrho _i=0`$. By the boundary condition for $`\varphi _i`$ in (6.50), we see that the tangent coordinates $`\varrho _i`$ obey the Neumann boundary condition
$$_1\varrho _i=0\text{at}\mathrm{\Sigma },$$
(6.62)
in the continuum limit. The boundary condition on the fermionic components can be read from (6.42) and (6.50) and is given by
$$\begin{array}{c}\mathrm{e}^{i\theta /2b_1}\chi _{i+}+\mathrm{e}^{i\theta /2b_1}\chi _i=0,\hfill \\ \mathrm{e}^{i\theta /2b_1}\overline{\chi }_{i+}+\mathrm{e}^{i\theta /2b_1}\overline{\chi }_i=0,\hfill \end{array}\text{at}\mathrm{\Sigma }.$$
(6.63)
These are the standard boundary condition on the worldsheet fields corresponding to the D-brane wrapped on $`\gamma _\theta `$. The phases $`\mathrm{e}^{i\theta /2b_1}`$ in the condition for the fermionic components shows that we are performing an R-rotation.
The cycle $`\gamma _\theta `$ is a Lagrangian submanifold of $`(𝐂^\times )^{N1}`$ with respect to the flat cylinder metric (6.44). The image of the cycle $`\gamma _\theta `$ in the $`\stackrel{~}{W}`$ plane is
$$\stackrel{~}{W}=\mathrm{e}^{i\theta /b_1}\underset{i=1}{\overset{N}{}}|\mathrm{e}^{y_i}|,$$
(6.64)
and is indeed a straight line (see Figure 18).
Moreover, the cycle passes through the critical point $`p_0`$ and the straight line in the $`\stackrel{~}{W}`$-plane emanates outward from the critical value $`\stackrel{~}{w}_0`$. Thus, $`\gamma _\theta `$ is the wavefront trajectory emanating from the critical point $`p_0`$. Since the image in the $`\stackrel{~}{W}`$-plane has the slope $`\theta /b_1`$ and the boundary condition of the fermions is rotated as (6.63), the boundary condition indeed preserves the supersymmetry $`Q`$ and $`Q^{}`$ in (6.52).
When $`\theta `$ is an integer multiple of $`2\pi `$, as we have noted before, in the non-linear sigma model limit the corresponding $`B`$-field is integral and the coupling of the worldsheet to such a $`B`$-field can be identified as the boundary coupling to the gauge field on $`X`$. In particular, for $`\theta =2\pi n`$, it is the gauge field of the bundle $`𝒪_X(n)`$. Thus, the D-brane wrapped on $`\gamma _{2\pi n}`$ in the LG model can be considered as the mirror of the D-brane on $`X`$ which supports the bundle $`𝒪_X(n)`$ in the sigma model with trivial $`B`$-field. Note that the image of $`\gamma _{2\pi n}`$ in the $`\stackrel{~}{W}`$-plane is the straight line emanating from $`\stackrel{~}{w}_n`$ in the radial direction (where the labeling of the critical values is made in the theory with $`\theta =0`$). More generally, the D-brane which supports the bundle $`𝒪_X(n)`$ in the sigma model with the $`B`$-field corresponding to $`\theta 0`$ is mirror to the D-brane wrapped on $`\gamma _{\theta +2\pi n}`$ whose image in the $`\stackrel{~}{W}`$-plane is the straight line emanating from $`\stackrel{~}{w}_n`$ in the radial direction.
In the next section, we make use of the connection explained in this section to determine the relation of the D-branes in the non-linear sigma model on the toric manifold $`X`$ (including more general cases corresponding to the gauge group $`U(1)^k`$ in the linear sigma model) and the D-branes of the mirror Landau-Ginzburg model.
## 7 D-Branes and Mirror Symmetry: Massive Theories
In this section we will study how D-branes transform under mirror symmetry for sigma models on Kähler manifolds with $`c_1>0`$. We will mainly concentrate on the case where the theory has only massive vacua, and discuss the mirror of D-branes wrapped on holomorphic cycles on the target Kähler manifold $`X`$ in terms of Lagrangian submanifolds of the mirror LG models. In particular we concentrate on D-branes which corresponds to exceptional bundles on $`X`$ (to be defined below). It turns out that this connection explains the observations of Kontsevich noting a formal correspondence between the properties of Helices and exceptional bundles on the Kähler manifolds and the soliton numbers of an associated LG model.
In order to do this, we will first review what exceptional bundles and Helices are. Afterwards we discuss how mirror symmetry acts in this context.
### 7.1 D-branes, Exceptional bundles and Helices
The similarities between the structures appearing in the classification of $`𝒩=2`$ theories and Helix theory was observed by Kontsevich . This observation was the starting point of in which the mysterious correspondence between the soliton numbers of a non-linear sigma model with $`\text{I}\text{P}^N`$ target space and exceptional bundles was further explored. We will explain this correspondence in this section as a consequence of mirror symmetry.
### 7.2 Exceptional bundles and mutations
A vector bundle or a sheaf $`E`$ on an $`N`$-dimensional variety $`X`$ with $`c_1>0`$ is called exceptional if
$$\text{Ext}^0(E,E)=𝐂𝖨,\text{Ext}^i(E,E)=0,i1,$$
(7.1)
where $`\text{Ext}^i`$ is the sheaf theory generalization of cohomology groups $`H^i`$, i.e., for vector bundles $`E`$ and $`F`$, $`\text{Ext}^i(E,F)=H^i(X,E^{}F)`$ which in tern equals to the Dolbeault cohomology $`H^{0,i}(X,E^{}F)`$. An exceptional collection is a collection of exceptional sheaves $`\{E_1,\mathrm{},E_n\}`$ such that if $`a<b`$ then
$`\text{Ext}^i(E_a,E_b)`$ $`=`$ $`0,ii_0\text{for some}i_0,`$ (7.2)
$`\text{Ext}^i(E_b,E_a)`$ $`=`$ $`0,i0.`$
Note that the above condition leaves $`dim\text{Ext}^{i_0}(E_a,E_b)`$ undetermined for $`a<b`$ and that could in principle be any integer. The alternating sum of dimensions of the groups $`\text{Ext}^i`$ defines a bilinear product ,
$$\chi (E,F)=\underset{i=0}{\overset{N}{}}(1)^i\text{dim}_{𝐂𝖨}\text{Ext}^i(E,F)=_X\text{ch}(E^{}F)\text{Td}(X).$$
(7.3)
An exceptional sheaf $`E`$ has the property that $`\chi (E,E)=1`$. An important property of an exceptional collection is that they can be transformed into new exceptional collections by transformations called mutations. For an exceptional collection of sheaves $`\{E_1,\mathrm{},E_n\}`$ we can sometimes define two transformations, left mutation and right mutation. Given a neighboring pair of sheaves $`(E_a,E_{a+1})`$ in an exceptional collection, the transformations $`L_{E_a}`$ and $`R_{E_a}`$ are such that
$`L_{E_a}(E_a,E_{a+1})=(L_{E_a}(E_{a+1}),E_a),R_{a+1}(E_a,E_{a+1})=(E_{a+1},R_{E_{a+1}}(E_a)),`$ (7.4)
The transformed sheaf $`L_{E_a}(E_{a+1})`$ is defined through an exact sequence. The exact sequence used to define the mutated sheaf depends on the $`\text{Ext}^i`$ groups of the pair $`(E_a,E_{a+1})`$ ,
$``$ If $`\text{Ext}^0(E_a,E_{a+1})0`$ and $`\text{Ext}^0(E_a,E_{a+1})E_aE_{a+1}`$ is surjective then $`L_{E_a}(E_{a+1})`$ is defined by the exact sequence,
$$0L_{E_a}(E_{a+1})\text{Ext}^0(E_a,E_{a+1})E_aE_{a+1}0,$$
(7.5)
$``$ If $`\text{Ext}^0(E_a,E_{a+1})0`$ and $`\text{Ext}^0(E_a,E_{a+1})E_aE_{a+1}`$ is injective then $`L_{E_a}(E_{a+1})`$ is defined by the exact sequence,
$$0\text{Ext}^0(E_a,E_{a+1})E_aE_{a+1}L_{E_a}(E_{a+1})0,$$
(7.6)
$``$ If $`\text{Ext}^1(E_a,E_{a+1})0`$ then $`L_{E_a}(E_{a+1})`$ is defined by exact sequence
$$0E_{a+1}L_{E_a}(E_{a+1})\text{Ext}^1(E_a,E_{a+1})E_a0.$$
(7.7)
Similarly we can define the right mutated sheaf $`R_{E_{a+1}}(E_a)`$ by an exact sequence,
$``$ If $`\text{Ext}^0(E_a,E_{a+1})0`$ and $`\text{Ext}^0(E_a,E_{a+1})E_aE_{a+1}`$ is surjective then $`R_{E_{a+1}}(E_a)`$ is defined by the exact sequence,
$$0E_a\text{Ext}^0(E_a,E_{a+1})^{}E_{a+1}R_{E_{a+1}}(E_a)0,$$
(7.8)
$``$ If $`\text{Ext}^0(E_a,E_{a+1})0`$ and $`\text{Ext}^0(E_a,E_{a+1})E_aE_{a+1}`$ is injective then $`R_{E_{a+1}}(E_a)`$ is defined by the exact sequence,
$$0R_{E_{a+1}}(E_a)E_a\text{Ext}^0(E_a,E_{a+1})^{}E_{a+1}E_{a+1}0,$$
(7.9)
$``$ If $`\text{Ext}^1(E_a,E_{a+1})0`$ then $`L_{E_a}(E_{a+1})`$ is defined by exact sequence
$$0\text{Ext}^1(E_a,E_{a+1})^{}E_{a+1}R_{E_{a+1}}(E_a)E_a0.$$
(7.10)
As far as the Chern characters are concerned the new sheaves $`L_{E_a}(E_{a+1})`$ and $`R_{E_{a+1}}(E_a)`$ are such that
$`\pm \text{ch}(L_{E_a}(E_{a+1}))`$ $`=`$ $`\text{ch}(E_{a+1})\chi (E_a,E_{a+1})\text{ch}(E_a),`$
$`\pm \text{ch}(R_{E_{a+1}}(E_a))`$ $`=`$ $`\text{ch}(E_a)\chi (E_a,E_{a+1})\text{ch}(E_{a+1}).`$ (7.11)
where this follows from the exact sequences used in the definition of the mutation and $`\pm `$ depends on which mutation one uses. The left and the right mutations are inverse of each other and satisfy the braid group relations <sup>1</sup><sup>1</sup>1 $`L_aL__a,R_aR__a`$.
$`L_aL_b`$ $`=`$ $`L_bL_a,R_aR_b=R_bR_a,\text{if}|ab|>1,`$
$`L_aL_{a+1}L_a`$ $`=`$ $`L_{a+1}L_aL_{a+1},R_aR_{a+1}R_a=R_{a+1}R_aR_{a+1}.`$ (7.12)
These transformations implement the braid group action on the collection of exceptional sheaves .
A helix of period $`n`$, $`\{E_i|i𝖹𝖹\}`$ is a collection of infinitely many exceptional sheaves such that
$`\{E_{i+1},`$ $`\mathrm{}`$ $`,E_{i+n}\},\text{is an exceptional collection for all }i𝖹𝖹,`$ (7.13)
$`E_{i+n}`$ $`=`$ $`R_{E_{i+n1}}\mathrm{}R_{E_{i+2}}R_{E_{i+1}}(E_i).`$ (7.14)
Thus given any exceptional collection $`\{E_1,\mathrm{},E_n\}`$ we can define a helix by extending the exceptional collection periodically i.e., $`E_{i+n}=R_{E_{i+n1}}\mathrm{}R_{E_{i+1}}(E_i)`$ and $`E_{i+n}=L_{E_{i+n1}}\mathrm{}L_{E_{i1}}(E_i)`$ for $`0in`$. The exceptional collection defining a helix is called the foundation of a helix. Such an exceptional collection generates the derived category of $`X`$ . For any exceptional collection $`\{E_i|i=1,\mathrm{},n\}`$ which is the foundation of a helix
$$R_{E_{i+n}}\mathrm{}R_{E_{i+1}}(E_i)=E_i\omega _X,$$
(7.15)
where $`\omega _X`$ is the canonical line bundle of $`X`$. The collection of line bundles $`\{𝒪(0),𝒪(1),\mathrm{},𝒪(n)\}`$ on $`\text{I}\text{P}^n`$ provides an important example of an exceptional collection which is also the foundation of a helix of period $`n+1`$. Their Chern character is given by $`\text{ch}(𝒪(k))=e^{kx}`$ where $`_{\text{I}\text{P}^n}x^n=1`$. In this case
$$R_{𝒪(n)}\mathrm{}R_{𝒪(1)}(𝒪(0))=𝒪(0)\omega _{\text{I}\text{P}^n}=𝒪(n+1).$$
(7.16)
The bilinear form for the exceptional collection on $`\text{I}\text{P}^n`$ is given by
$`\chi (𝒪(a),𝒪(b))`$ $`=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{n+ab}{ab}}\right),ab`$ (7.17)
$`=`$ $`0,a>b.`$ (7.18)
If we consider D-branes corresponding to exceptional sheaves on the Kähler manifold we can couple them to sigma models. In this context, as discussed earlier in this paper (3.46), $`\mathrm{Ext}^i(E,F)`$ is interpreted as the ground states in the open string sector stretched between $`E`$ and $`F`$ with fermion number $`i`$.
The similarities of the objects defined above and the D-branes we have studied in the context of LG models is striking and as we will discuss below not accidental: The D-branes we have constructed in the LG model turn out to be the mirror of the exceptional bundles. In particular the property that $`\text{dim}\text{Ext}^i(E,E)=\delta _{i,0}`$ is the statement we discussed before, namely the open string sector of a string stretched between the same D-brane has only one vacuum with fermion number 0 (the function $`W`$ gives a Morse function on it with exactly one critical point corresponding to an absolute minimum). Also the system of exceptional collection of sheaves has a natural parallel in the LG system: If we consider the D-branes on the cycles $`\gamma _i`$ that we constructed, ordered with decreasing value of $`\mathrm{Im}W`$, then the fact that for $`i<j`$ the open string stretched between $`\gamma _i`$ and $`\gamma _j`$ has no Ramond ground states and that for $`i>j`$ there can only be zero modes in this sector at a fixed fermion number, as discussed before, is exactly the conditions imposed on a collection of exceptional sheaves. Moreover the braiding with left and right mutations has also a natural parallel: If we change the combination of left and right supercharges that we are preserving the image of the $`\gamma _i`$ in the $`W`$-plane will rotate by the corresponding angle. Moreover as we change the angle two neighboring $`\gamma _i`$ and $`\gamma _{i+1}`$ might switch order. In this case the switched $`\gamma _i`$ define a different basis for $`H_n(𝐂^n,\mathrm{Re}We^{i\theta }>0)`$ related by Picard-Lefshetz action as discussed before. This is exactly the same as the change in the Chern characters of transmuted exceptional sheaves, up to the $`\pm `$ sign, which one can interpret as the orientation of the corresponding LG D-brane. Moreover the left versus right mutation corresponds to the reversal of the direction of change of $`\theta `$. Finally if we consider a massive sigma model with $`N`$ isolated vacua, taking any $`\gamma _i`$ around in the $`W`$ plane, as discussed before, is equivalent to changing the $`\theta _ik_i`$ B-fields of the sigma model by $`c_1`$ of the manifolds, which on the D-brane is realized as a tensoring with a $`U(1)`$ connection with curvature given by $`c_1`$, i.e. tensoring the D-brane with $`\omega _X`$. This is exactly the condition of having a helix of period $`N`$. Moreover once we discuss why $`\gamma _i`$ are the mirrors of the corresponding D-branes it becomes clear why, up to braidings, the number of solitons in Fano varieties (with only massive vacua) are given by the index of the $`\overline{}`$ operator coupled to $`E_i^{}E_j`$ where $`E_i,E_j`$ belong to a collection of exceptional bundles on the Fano variety.
Below we will present many examples of this connection. The discussions are aimed at giving a sample rather than an exhaustive search through examples.
### 7.3 $`\text{I}\text{P}^n`$
The $`\text{I}\text{P}^n`$ sigma model is realized as the $`U(1)`$ gauge theory with $`N=n+1`$ matters of charge $`1`$. The mirror is the $`A_n`$ affine Toda field theory with the superpotential
$$W=\mathrm{e}^{Y_1}+\mathrm{}+\mathrm{e}^{Y_n}+\lambda \mathrm{e}^{Y_1+\mathrm{}+Y_n},$$
(7.19)
where $`\lambda =\mathrm{e}^t=\mathrm{e}^{r+i\theta }`$. There are $`n+1`$ critical points $`p_0,\mathrm{},p_n`$ given by $`\mathrm{e}^{Y_i}(p_k)=\lambda ^{\frac{1}{n+1}}\mathrm{e}^{\frac{2\pi ki}{n+1}}`$ with the critical value $`w_k=(n+1)\lambda ^{\frac{1}{n+1}}\mathrm{e}^{\frac{2\pi ki}{n+1}}`$.
As we saw in section 6, the mirror of the trivial bundle $`𝒪(0)`$ for the sigma model with $`\theta =0`$ is the middle dimensional cycle whose image in the $`W`$-plane is a straight line starting at the critical value $`w_0`$ and extending along the real axis. We have also seen that the mirror of the bundle $`𝒪(k)`$ for the sigma model with $`\theta =0`$ is the Lagrangian submanifold whose image in the $`W`$-plane is a straight line emanating from $`w_k`$ in the radial direction. This has been obtained by the shift $`\theta =0\theta =2\pi k`$, i.e. the rotation $`\lambda ^{\frac{1}{n+1}}\lambda ^{\frac{1}{n+1}}\mathrm{e}^{\frac{2\pi ki}{n+1}}`$, and interpreting the result as the $`𝒪(k)`$ bundle for the sigma model with $`\theta =0`$. In this case, the unbroken supercharge is of the A-type combination<sup>1</sup><sup>1</sup>1Note that we have switched back to the standard convention of chirality: $`Y_i`$’s are chiral superfield rather than twisted chiral superfield. That is why the unbroken supercharges are A-type rather than B-type. $`Q=\overline{Q}_++\mathrm{e}^{\frac{2\pi ki}{n+1}}Q_{}`$ and its conjugate $`Q^{}`$. At this stage, one can rotate the Lagrangian cycle (without touching $`\theta `$) so that the image in the $`W`$-plane is parallel to the real axis. In such a case, the corresponding D-brane preserves the standard combination $`Q=\overline{Q}_++Q_{}`$ and its conjugate $`Q^{}`$.
These are depicted in the example of $`\text{I}\text{P}^5`$ in Figure 19.
For the bundles $`𝒪(2)`$, $`𝒪(3)`$,… or $`𝒪(2)`$, $`𝒪(3)`$,…, it is impossible for rotating the cycle so that the images in the $`W`$-plane are parallel to the real axis without passing through other critical values. One can avoid this cross-over by bending the branes although it results in the breaking of the supersymmetry. Bending in the clockwise direction, we obtain the collection of bundles $`\{𝒪(0),\mathrm{},𝒪(n+1)\}`$ which are exceptional collections as we have seen above. By partially changing the direction of bending, we can obtain other exceptional collections as shown in Fig. 20 for the case of $`\text{I}\text{P}^5`$. If we order the lines in terms of decreasing asymptotic imaginary part then the exceptional collection in Fig. 20(b) is $`\{𝒪(3),𝒪(2),𝒪(1),𝒪(0),𝒪(1),𝒪(2)\}`$ and the exceptional collection in Fig. 20(c) is $`\{𝒪(2),𝒪(1),𝒪(0),𝒪(1),𝒪(2),𝒪(3)\}`$.
The observation of can now be understood as a consequence of mirror symmetry. The soliton numbers between different vacua are given by intersection number of middle dimensional cycles determined by the superpotential as shown in section 2. Mirror symmetry relates these cycles and their intersection form to bundles on $`\text{I}\text{P}^n`$ and the bilinear form $`\chi (,)`$ respectively. To find the soliton numbers from this data, as reviewed in section 2, we need to choose suitable classes of cycles. This configuration of cycles is related to the D-branes we have by some Picard-Lefshetz action, which is the mirror realization of Left/Right mutations discussed in the case of exceptional bundles. Let us first discuss LG analog of mutation and then return to the computation of soliton numbers using the $`\chi (,)`$.
Let us denote the D-brane corresponding to the $`i`$-th critical point by $`C_i`$. Note that from the mirror symmetry map we have for $`i>j`$, $`C_jC_i=\chi (𝒪(j),𝒪(i))=(n+ij)!/n!(ij)!`$. Consider as an example the case of $`\text{I}\text{P}^2`$ and the bundle $`𝒪(2)`$ shown in Fig. 21. Making the middle dimensional $`C_2`$ cycle pass through the the critical value $`w_2`$ and using the Picard Lefshetz formula
$$C_2^{}=C_2(C_1C_2)C_1=C_23C_1.$$
(7.20)
which is exactly the same as how the left mutation acts on Chern character upon left mutation of the $`𝒪(2)`$ bundle over $`𝒪(1)`$. We thus identify $`C_2^{}`$ as the mirror of the $`𝐋_\mathrm{𝟏}𝒪(2)`$ (with the opposite orientation). Moreover in order for charges not to change, we see that we create three new D-branes $`+3C_1`$. If we again make the cycle $`C_2^{}`$ pass through $`w_0`$ we see using the Picard Lefshetz formula that,
$$C_2^{\prime \prime }=C_2^{}(C_0C_2^{})C_0=C_23C_1+3C_0.$$
(7.21)
Thus we see that we have created $`3C_0`$ branes, for charge conservation. Again we see that we can identify $`C_2^{\prime \prime }`$ with the mirror of $`𝐋_\mathrm{𝟎}𝐋_\mathrm{𝟏}(𝒪(2))`$. In other words the above process viewed in terms of the mutation of $`𝒪(2)`$ in the exceptional collection $`\{𝒪(0),𝒪(1),𝒪(2)\}`$, is given as
$`\{𝒪(0),𝒪(1),𝒪(\mathrm{𝟐})\}`$ $``$ $`\{𝒪(0),𝐋_\mathrm{𝟏}𝒪(\mathrm{𝟐}),𝒪(1)\}\{𝐋_\mathrm{𝟎}𝐋_\mathrm{𝟏}𝒪(\mathrm{𝟐}),𝒪(0),𝒪(1)\}`$
$`\text{ch}(L_1𝒪(2))`$ $`=`$ $`\text{ch}(𝒪(2))3\text{ch}(𝒪(1)),`$ (7.22)
$`\text{ch}(L_0L_1𝒪(2))`$ $`=`$ $`\text{ch}(L_0𝒪(2))3\text{ch}(L_0𝒪(1))`$
$`=`$ $`\text{ch}(𝒪(2))3\text{ch}(𝒪(1))+3\text{ch}(𝒪(0))=\text{ch}(𝒪(1))`$
The fact that mutating $`𝒪(2)`$ through all the other critical points gives $`𝒪(1)=𝒪(2)𝒪(3)`$ and that $`𝒪(3)`$ is the inverse of $`c_1`$ of the canonical bundles, is related to the axial anomaly of the $`\text{I}\text{P}^2`$ sigma model.
The left mutated bundle $`L_1(𝒪(2))`$ is shown in the figure below.
Fig. 22(a) and Fig. 22(b) represent two different representatives of the homology class mirror to $`L_1(𝒪(2))`$. In the case of Fig. 22(a) the representative is not a supersymmetric cycle since its image in the $`W`$-plane is not a straight line. The representative shown in Fig. 22(b), however, is supersymmetric and preserves A-model supercharge $`\overline{Q}_++e^{i\alpha }Q_{}`$, where $`\alpha `$ is the angle that the straight line makes with the real axis. This comment also applies to the branes depicted in figure 19, and other branes that will be discussed in this section.
As another interesting example of mutation consider the right mutation of the pair $`\{𝒪(0),𝒪(1)\}`$ on $`\text{I}\text{P}^n`$ as shown in Fig. 23.
Since $`\text{Ext}^0(𝒪(0),𝒪(1))=H^0(𝒪(0),𝒪(1))0`$, we can use the Euler sequence
$$0𝒪(0)H^0(𝒪(0),𝒪(1))^{}𝒪(1)T0$$
(7.23)
In fact $`R_{𝒪(1)}(𝒪(0))`$ is $`T`$, the tangent bundle, and its mirror is identified in the LG theory with the brane $`\gamma _0^{}`$ , shown in Fig. 23(b), which is obtained by Picard-Lefshetz monodromy action on $`\gamma _0`$.<sup>2</sup><sup>2</sup>2 It is natural to conjecture that the mirror of all bundles on $`\text{I}\text{P}^n`$ is given on the LG mirror by the D-branes corresponding to exceptional bundles with multiplicities given by the decomposition of its Chern character.
Now we come back to the question of counting the soliton numbers. The soliton number between $`w_0`$ and $`w_k`$ can be computed by first left mutating the $`𝒪(k)`$ brane through $`𝒪(k1),\mathrm{},𝒪(1)`$ and then taking its inner product with $`𝒪(0)`$. Since $`\text{ch}(L_1\mathrm{}L_{k1}𝒪(k))=_{i=0}^{k1}(1)^i\left(\genfrac{}{}{0pt}{}{n+1}{i}\right)\text{ch}(𝒪(ki))`$, hence soliton numbers between two different vacua $`w_0`$ and $`w_k`$ is equal to
$`\mu _{i,ik}=\mu _{0,k}`$ $`=`$ $`\chi (𝒪(0),L_1\mathrm{}L_{k1}𝒪(k))`$
$`=`$ $`{\displaystyle \underset{i=0}{\overset{k1}{}}}(1)^i\left({\displaystyle \genfrac{}{}{0pt}{}{n+1}{i}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{n+ki}{ki}}\right)=(1)^{k1}\left({\displaystyle \genfrac{}{}{0pt}{}{n+1}{k}}\right)`$
in agreement with what we had obtained before for the soliton numbers.
### 7.4 Toric del Pezzo Surfaces
We next consider the non-linear sigma model with two dimensional toric Fano varieties as the target space. We will show that the supersymmetric cycles of the mirror LG theory, which are the preimages of the straight lines in the $`W`$-plane defined by the superpotential, are related to an exceptional collection of bundles on the target space.
From the classification toric Fano varieties it is known that there are five toric Fano surfaces. The toric diagram of these surfaces is captured by a dual lattice shown in Fig. 24 (see and for a detailed discussion) which is obtained naturally from the mirror symmetry description derived in . The first four diagram are that of $`\text{I}\text{P}^2`$ and its three blow ups respectively. Fig. 24(e) is the toric diagram of $`F_0=\text{I}\text{P}^1\times \text{I}\text{P}^1`$.
The superpotential of the mirror LG theory can be written from the toric data given in the figure. Let $`\{v^{(a)}=(v_1^{(a)},v_2^{(a)})|a=1,\mathrm{},N\}`$ be the set of vertices. Then
$$W(X)=\underset{a=1}{\overset{N}{}}C_aX_1^{v_1^{(a)}}X_2^{v_2^{(a)}},$$
(7.25)
where $`C_a`$ are complex numbers. Only a subset of $`C_a`$ are actually physical since some of them can be absorbed by rescaling $`X_i`$.
#### 7.4.1 $`\text{I}\text{P}^2`$
From the discussion in the previous section we know that the lines shown in the Fig. 25(a) correspond to the exceptional collection $`\{𝒪(1),𝒪(0),𝒪(1)\}`$. The exceptional collection shown in Fig. 25(b) is $`\{𝒪(1),V,𝒪(0)\}`$ where $`V=L_0𝒪(1)`$ is such that $`\text{ch}(V)=(2,1,\frac{1}{2})`$. Here we are using the notation such that $`\text{ch}(V)=(c_0(V),c_1(V),_X\text{ch}_2(V))`$. What we actually mean by the negative number for $`c_0`$ is that if we reverse the orientation of the D-brane we obtain the corresponding mirror of the bundle. In other words the Chern character is multiplied by a minus sign, when comparing with the charges of the LG D-brane. We shall be somewhat implicit about this in this section, but it should be clear from the context what we mean– namely $`c_0`$ of the bundle should always be positive.
The soliton counting matrix can be determined from the exceptional collection shown in Fig. 25(b) and is given by $`\chi (E_i,E_j)`$ where $`E_i\{𝒪(1),V,𝒪(0)\}`$,
$$S_{ij}=\chi (E_i,E_j)=\left(\begin{array}{ccc}1& 3& 3\\ 0& 1& 3\\ 0& 0& 1\end{array}\right).$$
(7.26)
#### 7.4.2 $`_1`$
The superpotential is given by
$$W(X_1,X_2)=X_1+X_2+\frac{e^t}{X_1X_2}+\frac{e^{t_{E_1}}}{X_1},$$
(7.27)
where $`t_{E_1}=A(E)i\theta (E)`$ is the complexified Kähler parameter of the exceptional cycle, $`E_1`$. After rescaling the coordinates we can express the superpotential in the following form which will be more useful for the later discussion,
$$W(X_1,X_2)=e^{\frac{t}{3}}(X_1+X_2+\frac{1}{X_1X_2}+\frac{e^{t_{E_1}+\frac{2}{3}t}}{X_1}).$$
(7.28)
We define $`\mu _1=e^{t_{E_1}+\frac{2}{3}t}`$, then the critical values of the superpotential given above are
$`w=e^{\frac{t}{3}}(3y_{}+2\mu _1y_{}^2),\text{where}\mu _1y_{}^4+y_{}^31=0.`$ (7.29)
For $`|\mu _1|<<1`$ we see that
$`y_{}`$ $``$ $`\{e^{2\pi ik/3}+O(\mu _1),\frac{1}{\mu _1}+O(\mu _1)|k=0,1,2\},`$ (7.30)
$`w`$ $``$ $`e^{\frac{t}{3}}\{3e^{2\pi ik/3}+O(\mu _1),\frac{1}{\mu _1}+O(\mu _1)|k=0,1,2\}.`$ (7.31)
We will denote the first three critical values as $`\{w_k|k=0,1,2\}`$ and the fourth one as $`\widehat{w}_1`$. Thus we see that for $`|\mu _1|`$ very small the transformation $`\mu \mu e^{i\theta }`$ does not change the three symmetrically located critical values $`w_k`$ much but the critical value $`\widehat{w}_1`$ undergoes a clockwise rotation by an angle $`\theta `$ as shown in Fig. 26.
Thus to determine the bundles associated with the semi-infinite lines we first consider $`\mu _10`$. The critical value $`\widehat{w}_1`$ goes to $`\mathrm{}`$ and we are left with the case of $`\text{I}\text{P}^2`$ for which we know the correspondence. Thus the three bundles $`V_1,V_2`$ and $`V_3`$ shown in Fig. 27 correspond to the pull back of $`𝒪(1),𝒪(0)`$ and $`𝒪(1)`$ from $`\text{I}\text{P}^2`$ to $`_1`$ respectively.
To determine the bundle associated with the critical value $`\widehat{w}_1`$ we take $`\mu _1`$ to be very small and positive so that $`\widehat{w}_1`$ is far from other critical values. In this case the critical value is on the negative real axis as shown in the Fig. 27. Consider the case, as shown in Fig. 28, when there is the D-brane mirror to $`𝒪(0)`$, represented by the straight line starting from $`w_0`$, present. As mentioned before under the transformation $`\mu _1\mu _1e^{2\pi i}`$ the critical value $`\widehat{w}_1`$ makes a clockwise rotation around the origin, Fig. 28. As it passes through the D-brane associated with the bundle $`𝒪(0)`$, due to brane creation effect discussed earlier, it acquires a D-brane charge consistent with the homology class of the cycle associated with this critical point.
To determine the bundle corresponding to this new D-brane we use charge conservation. Recall that $`\mu _1`$ was defined in terms of $`t_{E_1}`$ and $`t`$, the complexified Kähler parameters of $`_1`$. We have kept $`t`$ fixed in above discussion therefore since the imaginary part of $`t_{E_1}`$ is minus the B-field integrated over the exceptional curve $`E_1`$ the transformation $`\mu _1\mu _1e^{2\pi i}`$ corresponds to turning on the B-field through $`E_1`$. Thus we interpret Fig. 28(e) as the $`𝒪(0)`$ bundle in the B-field background. This implies that if we denote the cohomology class dual to the exceptional curve $`E_1`$ as $`[E_1]`$, the first Chern class of this bundle ($`𝒪(0)`$ in the B-field background) is $`c_1(𝒪(0))[E_1]`$. Denoting by $`\widehat{V}_1`$ the bundle mirror to the line starting at $`\widehat{w}_1`$ we get,
$`\text{ch}(𝒪(0)\widehat{V}_1)`$ $`=`$ $`e^{c_1(𝒪(0))[E_1]},`$ (7.32)
$`\text{ch}(\widehat{V}_1)`$ $`=`$ $`e^{c_1(𝒪(0)[E_1]}\text{ch}(𝒪(0))`$
$`=`$ $`(1,[E]_1,\frac{1}{2})(1,0,0)=(0,[E_1],\frac{1}{2}).`$
Thus $`\widehat{V}_1`$ is actually a sheaf with support on the exceptional divisor $`E_1`$. The collection of sheaves $`\{V_1,V_2,V_3,\widehat{V}_1\}`$ is an exceptional collection since $`\chi (\widehat{V}_1,\widehat{V}_1)=1`$ and $`\chi (\widehat{V}_1,V_i)=0`$ for $`i=1,2,3`$.
It is interesting to consider the effect of transformation $`e^te^{t+2\pi i}`$. Since $`\mu _1=e^{t_{E_1}+\frac{2}{3}t}`$ we see that under the above transformation $`\mu _1\mu _1e^{\frac{4\pi i}{3}}`$. From the previous discussion it would seem that the critical point $`\widehat{w}_1`$ undergoes a rotation by an angle $`\frac{4\pi }{3}`$. This, however, is not the case. From eq. (7.28) it is clear that because of the overall factor $`e^{\frac{t}{3}}`$ a phase transformation $`e^te^{t+2\pi i}`$ rotates all the critical values by an angle $`\frac{2\pi }{3}`$. Thus the critical value $`\widehat{w}_1`$ gets rotated by $`\frac{4\pi }{3}+\frac{2\pi }{3}=2\pi `$,
$$e^te^{t+2\pi i}:w_1w_2w_3w_1,\widehat{w}_1\widehat{w}_1.$$
(7.33)
The effect on the bundles, however, is more non-trivial and is shown in Fig. 30.
$`V_1^{},V_2^{},V_3^{}`$ and $`\widehat{V}_1^{}`$ are $`V_1,V_2,V_3`$ and $`\widehat{V}_1`$ bundles in the background where the B-field through the cycle $`l`$ has been turned on, where $`l`$ along with $`E_1`$ forms a basis of $`H_2(_1,𝖹𝖹)`$ such that the self-intersection number of $`l`$ is plus one, $`l^2=1`$. Denoting the cohomology class dual to $`l`$ by $`[l]`$ we get
$`\text{ch}(V_k^{})=e^{c_1(V_k)[l]}=e^{(k3)[l]}=(1,(k3)[l],\frac{1}{2}(k3)^2)`$ (7.34)
Thus we see that $`V_k^{}`$ is the pull back of $`𝒪(k3)`$ bundle from $`\text{I}\text{P}^2`$ to $`_1`$ and can be written as $`\{V_1^{},V_2^{},V_3^{}\}=\{V_1V_1,V_2V_1,V_3V_1\}`$. The bundle $`\widehat{V}_1^{}`$ is easy to determine since under the transformation $`t_{E_1}t_{E_1}+2\pi i`$ the critical value $`\widehat{w}_1`$ rotates in the counter clockwise direction therefore argument similar to the one used to determine $`\widehat{V}_1`$ shows that $`\widehat{V}_1^{}`$ is such that
$$\text{ch}(V_2\widehat{V}_1^{})=e^{c_1(V_2)+[E_1]}\text{ch}(\widehat{V}_1^{})=(0,[E_1],\frac{1}{2}).$$
(7.35)
Using this exceptional collection we can calculate the the number of solitons between various vacua. The soliton number between two vacua is given by $`\chi (E,F)`$ where $`E`$ and $`F`$ are the bundle corresponding to the semi-infinite straight lines starting from the two vacua we are studying. The semi-infinite lines must be such that they together do not enclose another critical value. Otherwise the intersection will get contribution from the critical value enclosed by the lines. Thus we first transform to accomplish the exceptional collection given above into one for which no two lines enclose another critical value. This exceptional collection is shown in Fig. 31 and is obtained from $`\{V_1,V_2,V_3,\widehat{V}_1,\}`$ by successive left mutations shown in Fig. 31,
$`\{V_1,V_2,V_3,\widehat{V}_1\}`$ $``$ $`\{V_1,V_2,L_{V_3}(\widehat{V}_1),V_3\}`$
$``$ $`\{V_1,L_{V_2}L_{V_3}(\widehat{V}_1),V_2,V_3\}`$
$``$ $`\{V_1,L_{V_2}L_{V_3}(\widehat{V}_1),L_{V_2}(V_3),V_2\}=:\{E_1,E_2,E_3,E_4\}.`$
As far as the Chern characters of $`E_i`$ are concerned we have,
$`\text{ch}(E_1)=(1,1,\frac{1}{2}),\text{ch}(E_2)=(1,l+[E_1],0),\text{ch}(E_3)=(2,l,\frac{1}{2}),\text{ch}(E_4)=(1,0,0).`$
The soliton number between the vacua are now given by $`\chi (E_i,E_j)`$,
$$S_{ij}=\chi (E_i,E_j)=\left(\begin{array}{cccc}1& 1& 3& +3\\ 0& 1& +1& 2\\ 0& 0& 1& 3\\ 0& 0& 0& 1\end{array}\right)$$
(7.37)
In other words there is one soliton between vacua $`2,1`$ and $`2,3`$, three solitons between $`1,3`$ and $`1,4`$ and $`3,4`$ and two solitons between vacua $`2,4`$. Of course, as we change the Kähler parameters of the manifold the position of the vacua change and the number of solitons change, as reviewed in section 2.
#### 7.4.3 $`_2`$
The superpotential of the LG theory mirror to non-linear sigma model with $`_2`$ is given by
$$W(X_1,X_2)=X_1+X_2+\frac{e^t}{X_1X_2}+\frac{e^{t_{E_1}}}{X_1}+\frac{e^{t_{E_2}}}{X_2}.$$
(7.38)
$`t_{E_1}`$ and $`t_{E_2}`$ are the complexified Kähler parameters of the exceptional curves $`E_1`$ and $`E_2`$ respectively. As in the case of $`_1`$ we rescale the coordinates $`X_1`$ and $`X_2`$ to obtain the following form for the superpotential which will be useful for later discussion,
$$W(X_1,X_2)=e^{\frac{t}{3}}(X_1+X_2+\frac{1}{X_1X_2}+\frac{e^{t_{E_1}+\frac{2}{3}t}}{X_1}+\frac{e^{t_{E_2}+\frac{2}{3}t}}{X_2}).$$
(7.39)
We define $`\mu _i=e^{t_{E_i}+\frac{2}{3}t}`$ for $`i=1,2`$. The critical values of the superpotential are given by
$$w=\frac{1}{y_{}^2\mu _2}+2y_{}+\mu _1(y_{}^2\mu _2),\text{where}(y_{}^2\mu _2)^2(\mu _1y_{}+1)y_{}=0.$$
(7.40)
For $`|\mu _1|,\mu _2|<<1`$ we can see that leading order terms for the critical points and critical values are,
$`y_{}`$ $`=`$ $`\{e^{\frac{2\pi ik}{3}}+O(\mu _1,\mu _2),\frac{1}{\mu _1}+O(\mu _1,\mu _2),\mu _2+O(\mu _1,\mu _2^2)|k=0,1,2\},`$
$`w`$ $`=`$ $`\{3e^{\frac{2\pi ik}{3}}+O(\mu _1,\mu _2),\frac{1}{\mu _1}+O(\mu _1,\mu _2),\frac{1}{\mu _2}+O(\mu _1,\mu _2)|k=0,1,2\}.`$
We will denote the above critical values by $`w_0,w_1,w_2,\widehat{w}_1`$ and $`\widehat{w}_2`$ respectively.
To determine the bundle corresponding to $`\widehat{w}_1`$ we use the same argument as for the case of $`_1`$. As $`\mu _1,\mu _20`$ we recover the $`\text{I}\text{P}^2`$ configuration and therefore the the three bundles corresponding to the semi-infinite lines starting at $`w_1,w_2`$ and $`w_3`$ are the pull backs of the $`𝒪(1),𝒪(0)`$ and $`𝒪(1)`$ bundles from $`\text{I}\text{P}^2`$ to $`_2`$ respectively. We will continue to denote these bundle as $`V_1,V_2`$ and $`V_3`$ respectively as before even though they are different bundles than the ones considered in the last section. However, the Chern classes of these bundles are the same as before. Since in the limit $`\mu _20`$ we recover the $`_1`$ configuration therefore the bundle corresponding to the line starting at $`\widehat{w}_1`$ is the sheaf with support on the exceptional curve $`E_1`$. We will denote it, as before, by $`\widehat{V}_1`$.
As shown in Fig. 32 the critical value $`\widehat{w}_2`$ rotates in a clockwise direction as $`\mu _2\mu _2e^{2\pi i}`$. In the presence of D-brane corresponding to $`V_2`$ such a rotation of $`\widehat{w}_2`$ creates a D-brane whose image in the W-plane is the line starting at $`\widehat{w}_2`$ as shown in Fig. 32. We will denote the corresponding bundle by $`\widehat{V}_2`$. The transformation $`\mu _2\mu _2e^{2\pi i}`$ corresponds to turning on the B-field through the exceptional curve $`E_2`$. Denoting the cohomology class dual to $`E_2`$ by $`[E_2]`$, charge conservation implies that
$`\text{ch}(V_2\widehat{V}_2)=e^{c_1(V_2)[E_2]},`$
$`\text{ch}(\widehat{V}_2)=(0,[E_2],\frac{1}{2}).`$ (7.41)
The collection of bundles $`\{V_1,V_2,V_3,\widehat{V}_1,\widehat{V}_2,\}`$ is an exceptional collection. As explained before to calculate the soliton numbers we have to mutate this exceptional collection into an exceptional collections for which the corresponding semi-infinite lines are such that any two of them do not enclose a critical point. To obtain such an exceptional collection we consider the right mutations shown in Fig. 33.
$`\{V_1,V_2,V_3,\widehat{V}_1,\widehat{V}_2\}`$ $``$ $`\{V_1,V_2,L_{V_3}(\widehat{V}_1),L_{V_3}(\widehat{V}_2),V_3\}`$
$``$ $`\{V_1,L_{V_2}L_{V_3}(\widehat{V}_1),L_{V_2}L_{V_3}(\widehat{V}_2),V_2,V_3\}`$
$``$ $`\{V_1,L_{V_2}L_{V_3}(\widehat{V}_1),L_{V_2}L_{V_2}(\widehat{V}_2),L_{V_2}(V_3),V_2\}=:\{E_1,E_2,E_3,E_4,E_5\}.`$
The soliton numbers between different vacua are now given by $`\chi (E_i,E_j)`$ where the vacua are labeled in the counter clockwise direction as shown in Fig. 33,
$$S_{ij}=\chi (E_i,E_j)=\left(\begin{array}{ccccc}1& 1& 1& 3& +3\\ 0& 1& 0& 1& 2\\ 0& 0& 1& 1& 2\\ 0& 0& 0& 1& 3\\ 0& 0& 0& 0& 1\end{array}\right).$$
(7.43)
We can read off the number of solitons from the above matrix. Note that the matrix $`S`$ will change by a braid transformation under the change of parameters. Thus the matrix $`S`$ given above is the soliton counting matrix as long as the convexity of the critical values shown in Fig. 33 is maintained.
#### 7.4.4 $`_3`$
Now we will consider the case of $`\text{I}\text{P}^2`$ blown up at three points, $`_3`$. The superpotential of the mirror LG theory is
$$W__3(X_1,X_2)=X_1+X_2+\frac{e^t}{X_1X_2}+\frac{e^{t_{E_1}}}{X_1}+\frac{e^{t_{E_2}}}{X_2}+e^{t_{E_3}}X_1X_2.$$
(7.44)
$`t_{E_i}`$ are the complexified Kähler parameters of the exceptional curves $`E_i`$. After rescaling the coordinates we can write the above superpotential as,
$$W__3(X_1,X_2)=e^{\frac{t}{3}}(X_1+X_2+\frac{1}{X_1X_2}+\frac{e^{t_{E_1}+\frac{2}{3}t}}{X_1}+\frac{e^{t_{E_2}+\frac{2}{3}t}}{X_2}+e^{t_{E_3}\frac{1}{3}t}X_1X_2),$$
(7.45)
We denote by $`\mu _1,\mu _2`$ and $`\mu _3`$ the three parameters $`e^{t_{E_1}+\frac{2}{3}t},e^{t_{E_2}+\frac{2}{3}t}`$ and $`e^{t_{E_3}\frac{1}{3}t}`$ respectively. The positions of the critical values are determined by the $`\mu _i`$ and $`e^{\frac{t}{3}}`$ determines the overall scale. To determine the critical values let $`\mu _1=\mu _2=\mu _3=\mu `$. Then the critical points and the critical values are
$`(X_1,X_2)`$ $`=`$ $`\{(e^{\frac{2\pi ik}{3}},e^{\frac{2\pi ik}{3}}),(\mu ^2,\frac{1}{\mu })|k=0,1,2\},`$
$`W__3(X_1,X_2)`$ $`=`$ $`\{3e^{\frac{t}{3}}(e^{\frac{2\pi ik}{3}}+\mu e^{\frac{2\pi ik}{3}}),e^{\frac{t}{3}}(\frac{1}{\mu }+\mu ^2)|k=0,1,2\}.`$ (7.46)
The critical value $`w_4`$ is degenerate and the three critical values at this point can be separated by taking $`\mu _i\mu _j`$. For $`\mu `$ close to zero we see that as $`\mu \mu e^{2\pi i}`$ the critical values $`w_1,w_2,w_3`$ move in a closed path close to the original critical value. On the other hand the critical value $`w_4`$ moves around the origin in the clockwise direction. The solution given above for the critical values in terms of $`\mu `$ is an exact solution. For $`\mu _i\mu _j`$ we can construct approximate solutions that specify the behavior of critical values for $`|\mu _i|`$ small. This was done for the case of $`_1`$ and $`_2`$ and the result here is similar to that case. We denote the three degenerate critical values by $`\widehat{w}_i`$ such that $`\mu _i0`$ implies $`\widehat{w}_i\mathrm{}`$. Under the transformation $`\mu _3\mu _3e^{2\pi i}`$, $`\widehat{w}_3`$ undergoes a clockwise rotation around the origin as shown in Fig. 34. Since taking $`\mu _1,\mu _20`$ has no effect on the three symmetrically located critical values therefore we can identify the corresponding bundles as the pull backs of $`𝒪(1),𝒪(0)`$ and $`𝒪(1)`$ from $`\text{I}\text{P}^2`$ to $`_3`$, we will denote these bundles by $`V_1,V_2`$ and $`V_3`$ respectively. Similarly from Fig. 34 and charge conservation we see that the bundle associated with $`\widehat{w}_i`$ denoted by $`\widehat{V}_i`$ is such that,
$$\text{ch}(\widehat{V}_3)=(0,E_3,\frac{1}{2}).$$
(7.47)
It is easy to check that $`\{V_1,V_2,V_3,\widehat{V}_1,\widehat{V}_2,\widehat{V}_3\}`$ is an exceptional collection. The semi-infinite lines corresponding to this collection are shown in Fig. 35.
To calculate the number of solitons between the vacua we need to transform this collection into the one shown in Fig. 36.
We see that the exceptional collection shown in Fig. 35 is obtained from $`\{V_1,V_2,V_3,\widehat{V}_1,\widehat{V}_2,\widehat{V}_3\}`$ by successive left mutations. Note that since $`\chi (\widehat{V}_i,\widehat{V}_j)=\delta _{ij}`$ therefore we can move the corresponding branes through each other without generating any new branes,
$`\{V_1,V_2,V_3,\widehat{V}_1,\widehat{V}_2,\widehat{V}_3\}`$ $``$ $`\{V_1,V_2,V_3,\widehat{V}_2,\widehat{V}_3,\widehat{V}_1\}`$
$``$ $`\{V_1,V_2,L_{V_3}(\widehat{V}_2),L_{V_3}(\widehat{V}_3),L_{V_3}(\widehat{V}_1),V_3\}`$
$``$ $`\{V_1,L_{V_2}L_{V_3}(\widehat{V}_2),L_{V_2}L_{V_3}(\widehat{V}_3),L_{V_2}L_{V_3}(\widehat{V}_1),V_2,V_3\}`$
$``$ $`\{V_2,L_{V_2}L_{V_3}(\widehat{V}_2),L_{V_2}L_{V_3}(\widehat{V}_3),L_{V_2}L_{V_3}(\widehat{V}_1),L_{V_2}(V_3),V_2\}`$
$`=:`$ $`\{E_1,E_2,E_3,E_4,E_5,E_6\}.`$
The soliton counting matrix is given by
$$S_{ij}=\chi (E_i,E_j)=\left(\begin{array}{cccccc}1& 1& 1& 1& 3& +3\\ 0& 1& 0& 0& +1& 2\\ 0& 0& 1& 0& 1& 2\\ 0& 0& 0& 1& 1& 2\\ 0& 0& 0& 0& 1& 3\\ 0& 0& 0& 0& 0& 1\end{array}\right)$$
(7.49)
Note that this matrix gives the number of solitons only for those values of the parameters for which we can have the convex configuration shown in Fig. 36.
### 7.5 $`F_0=\text{I}\text{P}^1\times \text{I}\text{P}^1`$
Since $`F_0`$ is a product manifold the superpotential of the mirror LG theory is just two copies of the superpotential of LG theory mirror to $`\text{I}\text{P}^1`$,
$`W_{F_0}(X_1,X_2)=X_1+X_2+{\displaystyle \frac{e^{t_1}}{X_1}}+{\displaystyle \frac{e^{t_2}}{X_2}}.`$
$`t_1`$ and $`t_2`$ are complexified Kähler parameters of the two $`\text{I}\text{P}^1`$’s. After rescaling the variables we can write the above superpotential as
$$W_{F_0}(X_1,X_2)=e^{\frac{t_1}{2}}(X_1+\frac{1}{X_1})+e^{\frac{t_2}{2}}(X_2+\frac{1}{X_2}).$$
(7.50)
The critical points and the corresponding critical values of above superpotential are,
$`(X_1,X_2)`$ $`=`$ $`\{(1,1),(1,1),(1,1),(1,1)\},`$
$`W(X_1,X_2)`$ $`=`$ $`w_i=\{2e^{\frac{t_1}{2}}+2e^{\frac{t_2}{2}},\mathrm{\hspace{0.17em}2}e^{\frac{t_1}{2}}2e^{\frac{t_2}{2}},2e^{\frac{t_1}{2}}+2e^{\frac{t_2}{2}},2e^{\frac{t_1}{2}}2e^{\frac{t_2}{2}}\}.`$
Without loss of generality we assume that $`|e^{\frac{t_1}{2}}||e^{\frac{t_2}{2}}|`$. To determine the bundles associated with the cycles $`𝒞_i`$ (which are the preimages of the semi-infinite lines in the W-plane) we consider the configuration with $`e^{\frac{t_1}{2}}`$ and $`e^{\frac{t_2}{2}}`$ real. The critical values in the W-plane in this case lie on the real axis and are non-degenerate as long as $`e^{\frac{t_1}{2}}e^{\frac{t_2}{2}}`$ as shown in Fig. 37.
We know that the cycle $`𝒞_1`$ (whose image in the W-plane is the semi-infinite line starting at $`w_1`$) is mirror to the trivial bundle $`𝒪(0,0)`$ <sup>1</sup><sup>1</sup>1A bundle $`𝒪(a,b)`$ on $`F_0`$ is a rank one bundle with first Chern class $`c_1=al_1+bl_2`$, where $`l_1`$ and $`l_2`$ are the generators of $`H_2(F_0)`$ such that $`l_1l_1=l_2l_2=0`$ and $`l_1l_2=1`$.. Note that the transformation $`e^{t_i}e^{t_i}e^{2\pi i}`$ has the following effect on the critical values,
$`e^{t_1}e^{t_1}e^{2\pi i}w_1w_3,w_2w_4,`$
$`e^{t_2}e^{t_2}e^{2\pi i}w_1w_2,w_3w_4.`$ (7.51)
To determine the bundle mirror to the cycle $`𝒞_2`$ (whose image in the W-plane is the semi-infinite line starting at $`w_2`$ and going to infinity along $`e^{iϵ}`$ ($`ϵ<<1`$) we perform the second transformation given above. The effect of this transformation on the image of cycle $`𝒞_1`$ in the W-plane is shown in Fig. 38.
Thus the bundle,$`V_2`$, mirror to $`𝒞_2`$ is such that,
$$\text{ch}(V_2)=\text{ch}(𝒪(0,0))+\text{B-field}=e^{c_1(𝒪(0,0))l_2}=1l_2=(1,l_2,0).$$
(7.52)
Thus we can identify $`V_2`$ with $`𝒪(0,1)`$. The effect of the first transformation of eq. (7.51) on the cycle $`𝒞_1`$ is shown in Fig. 39.
Thus by charge conservation we see that the bundle,$`V_3`$, mirror to $`𝒞_3`$ is such that
$$\text{ch}(V_3)=\text{ch}(𝒪(0,0))+\text{B-field}=e^{c_1(𝒪(0,0))l_1}=1l_1=(1,l_1,0).$$
(7.53)
Thus we can identify $`V_3`$ with $`𝒪(1,0)`$. Now consider the transformation
$`(e^{t_1},e^{t_2})(e^{i\theta }e^{t_1},e^{i\theta }e^{t_2}),\theta [0,2\pi ].`$ (7.54)
The effect of this transformation on the critical values and the cycle $`𝒞_1`$ is shown in Fig. 40.
Thus charge conservation implies that the bundle $`V_4`$ mirror to the cycle $`𝒞_4`$ is such that
$`\text{ch}(V_4)=\text{ch}(𝒪(0,0))+\text{B-field}=e^{c_1(𝒪(0,0))l_1l_2}=(1,l_1l_2,1).`$ (7.55)
Thus we can identify $`V_4`$ with $`𝒪(1,1)`$.
The set of bundles $`\{𝒪(1,1),𝒪(1,0),𝒪(0,1),𝒪(0,0)\}`$ is an exceptional collection and using it we can calculate the soliton counting matrix. But first since we do not want any three vacua to be collinear we deform the configuration using eq. (7.54) for $`|\theta |<<1`$. And also we need to transform this collection of exceptional bundles into the one shown in Fig. 41 by left or right mutations.
$`\{𝒪(1,1),𝒪(1,0),𝒪(0,1),𝒪(0,0)\}`$
$`\{𝒪(1,0),R_{𝒪(1,0)}(𝒪(1,1)),𝒪(0,1),𝒪(0,0)\}=:\{E_1,E_2,E_3,E_4\}.`$
The soliton counting matrix $`S_{ij}`$ is then given by
$`S_{ij}=\chi (E_i,E_j)=\left(\begin{array}{cccc}1& 2& 0& 2\\ 0& 1& 2& 0\\ 0& 0& 1& 2\\ 0& 0& 0& 1\end{array}\right).`$ (7.56)
### 7.6 Higher dimensional toric Fano varieties
As an example of higher dimensional toric Fano varieties we consider the blow ups of projective spaces. Blow up of $`\text{I}\text{P}^{n1}`$ upto $`n`$ points is a toric Fano variety . Each blow up corresponds to replacing a point by $`\text{I}\text{P}^{n2}`$ and thus each blow up introduces $`n2`$ new cohomology elements. We consider the case of maximal blow ups since others can be obtained from this one as we saw for the case of two dimensional del Pezzo surfaces.
The linear sigma model is a $`U(1)^n`$ gauge theory. The $`2n`$ chiral superfields have the following charge assignment under the $`U(1)^n`$ gauge group,
$`Q_1`$ $`=`$ $`(1,1,1\mathrm{}1,1;0,0,0,\mathrm{}0,0),`$ (7.57)
$`Q_2`$ $`=`$ $`(0,1,1\mathrm{}1,1;1,0,0,\mathrm{}0,0),`$
$`Q_3`$ $`=`$ $`(1,0,1\mathrm{}1,1;0,1,0,\mathrm{}0,0),`$
$`\mathrm{}`$
$`Q_n`$ $`=`$ $`\underset{n}{\underset{}{(1,1,1\mathrm{}1,0}};\underset{n}{\underset{}{0,0,0,\mathrm{},0,1}}).`$
The LG superpotential is then given by
$`W(X)={\displaystyle \underset{i=1}{\overset{n1}{}}}X_i+{\displaystyle \frac{e^t}{X_1\mathrm{}X_{n1}}}+{\displaystyle \underset{i=1}{\overset{n1}{}}}{\displaystyle \frac{e^{t_i}}{X_i}}+e^{t_{n1}}X_1\mathrm{}X_{n1}.`$
After rescaling $`X_i`$ we can write the above superpotential as
$$W(X)=e^{\frac{t}{n}}(\underset{i=1}{\overset{n1}{}}X_i+\frac{1}{X_1\mathrm{}X_{n1}}+\underset{i=1}{\overset{n1}{}}\frac{e^{t_i+\frac{2}{n}t}}{X_i}+e^{t_{n1}+\frac{n2}{n}t}X_1\mathrm{}X_{n1}).$$
(7.58)
Let $`\mu _i=e^{t_i+\frac{2}{n}t}`$ for $`i=1,\mathrm{},n2`$ and $`\mu _{n1}=e^{t_{n1}+\frac{n2}{n}t}`$. Consider the case when $`\{\mu _i=\mu |i=1,\mathrm{},n1\}`$, in this case there are $`2(n1)`$ critical points given by
$$X_i=f,(\mu f^{n2}+1)(f^n1)=0.$$
(7.59)
$`n`$ of these critical points are also the critical point of the $`\text{I}\text{P}^{n1}`$ superpotential. The new critical points and critical values are
$`X_i^{(k)}`$ $`=`$ $`\mu ^{\frac{1}{n2}}e^{\frac{i\pi (2k+1)}{n2}},i=1,\mathrm{},n1,k=1,\mathrm{},n2,`$
$`\widehat{w}_k`$ $`=`$ $`W(X^{(k)})=n(\mu ^{\frac{1}{n2}}e^{\frac{i\pi (2k+1)}{n2}}+\mu ^{n1}e^{\frac{i\pi (2k+1)}{n2}})`$ (7.60)
For $`\mu _i\mu _j`$ each of the above new critical points splits up into $`n1`$ critical points. Thus the critical value $`\widehat{w}_k`$ is degenerate with multiplicity $`n1`$. If $`\mu _i0`$ then $`n2`$ critical values go to infinity and the multiplicity of $`\widehat{w}_k`$ reduces to $`n2`$. To determine the bundles corresponding to the lines ending on these new critical values we only need to consider the case when $`\mu _10`$ and $`\mu _i=0`$ for $`i=2,\mathrm{},n1`$. This is the case of $`\text{I}\text{P}^{n1}`$ blown up at one point. In this case there are $`2n2`$ non-degenerate critical points given by the following equation
$$X_i=f,\mu _1f^{2n2}+f^n1=0,$$
For $`\mu _1`$ very small we can write the leading terms in the solution as
$`X_i^{(k)}`$ $`=`$ $`f_k=e^{\frac{2\pi ik}{n}}+O(\mu _1),X_i^{(k^{})}=f_k^{}=\mu _1^{\frac{1}{n2}}e^{\frac{i\pi (2k^{}1)}{n2}}+O(\mu _1),`$
$`w_k`$ $`=`$ $`W(X_i^{(k)})=ne^{\frac{2\pi ik}{n}}+O(\mu _1),\widehat{w}_k^{}=W(X_i^{(k^{})})=\mu _1^{\frac{1}{n2}}e^{\frac{i\pi (2k^{}1)}{n2}}+O(\mu _1),`$
where $`k\{0,\mathrm{},n1\}`$ and $`k^{}\{1,\mathrm{},n2\}`$. Thus we see that as $`\mu _1\mu _1e^{2\pi i}`$ the critical values $`\widehat{w}_k^{}`$ are rotated by $`e^{\frac{2\pi i}{n2}}`$. Hence $`w_k^{}^{}`$ is mapped to $`w_{k^{}1}^{}`$ by this transformation as shown in Fig. 42 for the case of $`\text{I}\text{P}^5`$.
Consider the case of Fig. 43 where we have the D-brane corresponding to the bundle $`𝒪(0)`$ ending on the critical value $`w_0`$ on the positive real axis. After the transformation $`\mu _1\mu _1e^{2\pi i}`$ we create another D-brane whose image in the W-plane is the semi-infinite line starting at $`w_{n2}^{}`$.
Denote by $`B`$ the cohomology class dual to the exceptional divisor and by $`V_1`$ the new bundle created then by charge conservation it follows that,
$`\text{ch}(𝒪(0)V)=e^{c_1(𝒪(0))B}\text{ch}(V)=e^B1,`$ (7.61)
$`c_0(V)=0,c_1(V)=B,c_i(V)=0,i=2,\mathrm{},n.`$ (7.62)
Thus $`V`$ is the $`[𝒪(0)]`$ bundle on $`\text{I}\text{P}^{n2}`$. A transformation $`\mu _1\mu _1e^{2\pi im}`$ maps it to $`\widehat{w}_{n2m}`$, thus the bundle which corresponds to the semi-infinite line starting at $`\widehat{w}_k^{}`$ is the line bundle $`𝒪(n2k^{})`$ on $`\text{I}\text{P}^{n2}`$, a sheaf on $`\text{I}\text{P}^{n1}`$ with support on the exceptional divisor. The case of $`\text{I}\text{P}^5`$ is shown in Fig. 44.
## 8 D-Brane in String Theory and Mirror Symmetry
It is natural to ask how the map between D-branes corresponding to sheaves and the Lagrangian submanifolds in the Landau-Ginzburg model works in the case of conformal theories. In the context of Gepner models, the structure of Cardy states and their sigma model interpretation have been studied . our constructions of Cardy states in terms of Lagrangian submanifolds of LG models lead to a deeper geometric insight in this regard. Since we have considered the case of minimal models in detail, and Gepner model is an orbifold of their tensor products, it is straightforward to identify the relevant D-branes in the orbifold LG model.
There is however, another case of the conformal theory we can consider namely the non-compact CY manifolds. These are the cases of most interest in the context of geometric engineering of QFT’s. As an example of this class consider the total space of the canonical line bundle over a compact Fano variety. This space is a non-compact Calabi-Yau manifold. We find relations between the LG theories mirror to the superconformal sigma model on the non-compact Calabi-Yau and the sigma model on the Fano variety.
For simplicity let us assume that the Fano variety is given by the weighted projective space with weights ($`q_i>0`$). Then the total space of the canonical bundle is captured by a linear sigma model with a single $`U(1)`$ gauge theory with matter fields with charges $`(_iq_i,q_1,q_2,\mathrm{},q_n)`$. The mirror of this is an LG theory with of $`n+1`$ variables with superpotential,
$$W=\underset{i=0}{\overset{n}{}}x_i\mathrm{subject}\mathrm{to}\underset{i=1}{\overset{n}{}}x_i^{q_i}=e^tx_0^{{\scriptscriptstyle q_i}}.$$
(8.1)
Recall that the correct field variables are $`Y_i`$ where $`x_i=e^{Y_i}`$. For simplicity let us assume one of the charges say $`q_n=1`$ (the more general case can also be done with the additional complication of introducing orbifold groups). Then we can also write the above superpotential as
$$W=x_0[1+\underset{i=1}{\overset{n1}{}}\widehat{x}_i+\frac{e^t}{_{i=1}^{n1}\widehat{x}_i^{q_i}}].$$
(8.2)
Where $`\widehat{x}_i=e^{\widehat{Y}_i}`$ and $`\widehat{Y}_i=Y_iY_0`$. This change of fields is linear and introduces no Jacobians in the field measure. Now, as far as periods and BPS states which are sensitive only to period integrals $`d\varphi _je^W`$ are concerned this LG theory is equivalent to the LG theory given by
$$W=x_0[1+\underset{i=1}{\overset{n1}{}}\widehat{x}_i+\frac{e^t}{_{i=1}^{n1}\widehat{x}_i^{q_i}}uv],$$
(8.3)
where now $`x_0`$ is the right field variable (i.e. $`x_0𝐂`$ rather than $`𝐂^{}`$), and $`u,v`$ are chiral fields also taking value in $`𝐂`$. To see this note that in the BPS computations integrating over the $`u,v`$ fields leads to a $`\frac{1}{x_0}`$ in the measure which combined with $`dx_0`$ converts it back to the measure appropriate for $`x_0`$ taking values in $`𝐂^{}`$, and leading to the previous LG periods. Having established their equivalence (at least in the weak sense discussed in ), in the period integrals we can integrate out $`x_0`$ in the new version of the LG theory, and obtain a $`\delta (1+_{i=1}^{n1}\widehat{x}_i+\frac{e^t}{_{i=1}^{n1}\widehat{x}_i^{q_i}}uv)`$. Thus we see that as far as the BPS data is concerned the mirror of the sigma model on the non-compact CY, which is originally the LG model, is equivalent to the sigma model on another non-compact Calabi-Yau given by
$$f(\widehat{x}_i)=uv\mathrm{where}f(\widehat{x}_i)=1+\underset{i=1}{\overset{n1}{}}\widehat{x}_i+\frac{e^t}{_{i=1}^{n1}\widehat{x}_i^{q_i}},$$
(8.4)
where $`\widehat{x}_i`$ take values in $`𝐂^{}`$ but $`u,v`$ are variables in $`𝐂`$. Note that this non-compact mirror CY has dimension $`n`$ which is the same as the dimension of the original non-compact CY. The holomorphic $`n`$ form can be viewed as
$$\mathrm{\Omega }=\frac{\underset{i=1}{\overset{n}{}}\frac{d\widehat{x}_i}{\widehat{x}_i}dudv}{df}=\frac{\underset{i=1}{\overset{n}{}}\frac{d\widehat{x}_i}{\widehat{x}_i}du}{u}.$$
(8.5)
It is this version of local mirror symmetry that was first discovered in the literature .
We see a striking resemblance between the non-conformal sigma model on the Fano variety and the conformal sigma model on the total space of the canonical line bundle over the Fano variety. In particular the $`f(\widehat{x}_i)`$ appearing in the above formula is precisely the superpotential for the LG theory mirror to the non-conformal sigma model on the Fano variety.
Even though we have presented the above discussion in the context of a linear sigma model with a single $`U(1)`$, it can be easily generalized to the case with more $`U(1)`$’s as well as with extra superpotentials corresponding to complete intersections.
### BPS states and local mirror symmetry
In computing the BPS states in such local contexts the idea developed in was to consider supersymmetric mid-dimensional cycles on the mirror. Moreover one could simplify the counting of such cycles by considering fibration structure of the non-compact CY and studying the supersymmetric cycles on the fibers and consider the effective tension of the branes as one varies over the base. In this way it was shown in , in the context of local mirror of $`SU(N)`$ gauge theories, how the problem is translated to finding minimal energy string configurations (with varying tensions) on a Riemann surface given by $`f(x_1,x_2)=0`$. This was implemented in detail for the $`SU(2)`$ case where various expected properties of BPS states in the corresponding $`𝒩=2`$ gauge theory in 4 dimensions was recovered including the decay of certain BPS states. Further applications along these lines have been considered .
We can also connect the above description to the BPS states for the probes in the context of F-theory, which is the subject of the next section.
### 8.1 Local mirror symmetry and F-theory
In this section we will see that the W-plane geometry of LG theory mirror to sigma model with certain non-compact CY threefolds as target space is closely related to some F-theory backgrounds . The link we find is as follows: The BPS states on the non-compact CY threefold side are D-branes wrapped on compact even dimensional cycles. As discussed above these get transformed on the mirror side to certain 3-cycles in a non-compact CY 3-fold. For the particular backgrounds of interest the non-compact CY 3-fold itself has a simple $`𝐂^{}`$ fibration structure over CY 2-fold (a local description of elliptic $`K3`$). The image of the closed 3-cycles get mapped to minimal 2-cycles in this geometry, which can possibly have boundaries where the $`𝐂^{}`$ fibration degenerates. This in turn can be viewed as computation of BPS state in a certain F-theory background with a 3-brane probe (placed where the $`𝐂^{}`$ fibration degenerates).
We will first review the probe theory description for F-theory and its BPS states and then give some examples of non-compact CY manifolds and the corresponding probe theory.
#### 8.1.1 Probe theory and BPS states
Consider a manifold $`𝒳`$ which is an elliptic fibration over the complex plane $`B`$
$$y^2=x^3+f(z)x+g(z),zB,$$
(8.6)
provided with a non-vanishing holomorphic 2-form $`\mathrm{\Omega }`$
$$\mathrm{\Omega }=\lambda dz,$$
(8.7)
where $`\lambda =\frac{dx}{y}`$ is the holomorphic 1-form on the elliptic fibers. F-theory compactification on $`𝒳`$ is equivalent to type IIB compactification on the base $`B`$ with a varying coupling constant $`\tau `$ (defined up to $`SL(2,𝖹𝖹)`$ transformations) given by the complex structure of the elliptic fiber. The position of the degenerate elliptic fibers on the base is given by the zeroes of the discriminant, $`\mathrm{\Delta }(z)`$, of the elliptic fibration (8.6),
$$\mathrm{\Delta }(z)=4f(z)^3+27g(z)^2.$$
(8.8)
From Picard Lefshetz theory we know that as we go around the position of a degenerate fiber in the base, the complex structure parameter $`\tau `$ undergoes an $`SL(2,𝖹𝖹)`$ transformation. As mentioned before this complex structure parameter is identified with the coupling constant of type IIB. Since in type IIB monodromies associated with 7-branes transform $`\tau `$ by $`SL(2,𝖹𝖹)`$ transformations, the position of a degenerate fiber on the base, in type IIB, is associated with a 7-brane. $`SL(2,𝖹𝖹)`$ symmetry of type IIB then implies the existence of a family of 7-branes labelled by two relatively prime integers, $`(p,q)`$. As for the case of 7-brane, a $`(p,q)`$ 7-brane at a point $`z_{}`$ can be associated, in F-theory, with an elliptic fiber over $`z_{}`$, $`T_z_{}^{\mathrm{\hspace{0.17em}2}}`$, whose degenerating 1-cycle is $`p\alpha +q\beta H_1(T_z_{}^{\mathrm{\hspace{0.17em}2}},𝖹𝖹)`$.
In ref., $`𝒩`$=2 $`SU(2)`$ Seiberg-Witten theory was interpreted as the worldvolume theory of a D3-brane in the presence of mutually non-local 7-branes. A BPS state of charge $`(p,q)`$ in the D3-brane theory is a BPS string or a BPS string junction of total asymptotic charge $`(p,q)`$ with support on 7-branes and ending on the D3-brane. In the F-theory picture D3-brane lifts to a regular elliptic curve of $`𝒳`$. Strings or string junctions stretched between the 7-branes and the D3-brane are, in F-theory, two real dimensional curves in the manifold $`𝒳`$ with or without boundary depending on whether the string junction ends on the D3-brane or not<sup>1</sup><sup>1</sup>1 This description follows from the connection between F-theory compactified on a circle and the M-theory in one lower dimension. In the M-theory description the D3-brane probe gets mapped to an M5 brane wrapped over the corresponding elliptic fiber and the BPS states are M2 branes wrapped over 2-cycles of the $`K3`$ geometry, possibly ending on the M5 brane. The image of the M2 brane projected on the $`z`$-plane gives the string junction description in the type IIB setup.. If the string junction ends on a D3-brane the corresponding curve in F-theory has a boundary on the elliptic curve above the position of the D3-brane. The homology cycle of the boundary is determined by the $`(p,q)`$ charge of the string junction ending on the D3-brane. BPS string junctions correspond to curves holomorphic in the complex structure whose kähler form is $`\mathrm{\Omega }`$. The mass of a BPS state of charge $`(p,q)`$ is given by the area of the corresponding curve $`𝒞_{p,q}`$,
$$M_{p,q}=|_{𝒞_{p,q}}\mathrm{\Omega }|.$$
(8.9)
#### 8.1.2 Superpotentials and F-theory backgrounds
We will see in this section that the non-compact CY 3-folds with an equation of the form $`f(x_1,x_2)=uv`$ where $`x_1,x_2`$ are $`𝐂^{}`$ variables and $`u,v`$ are $`𝐂`$ variables get related to the F-theory probe description. We first discuss the structure of the BPS D3-branes in the non-compact local Calabi-Yau description and then relate it to the F-theory description.
Instead of being general, we consider a concrete example. The general case is similar. Consider the case of $`𝒪(3)`$ over $`\text{I}\text{P}^2`$. This non-compact CY threefold, which is the total space of $`𝒪(3)`$ bundle over $`\text{I}\text{P}^2`$, will be denoted by $``$. This has linear sigma model description in terms of a single $`U(1)`$ gauge theory with charges of the matter fields $`(3,1,1,1)`$. The LG superpotential of the mirror theory is,
$$W(x)=x_0+x_1+x_2+e^t\frac{x_0^3}{x_1x_2}.$$
(8.10)
As we discussed before as far as the BPS data is concerned the non-compact CY defined by eq. (8.10) is equivalent to another non-compact CY, $`\widehat{}`$, defined by,
$`1+x_1+x_2+{\displaystyle \frac{e^t}{x_1x_2}}=z,z=uv.`$ (8.11)
where $`x_1,x_2`$ are $`𝐂^{}`$ variables and $`u,v,z`$ are variables in $`𝐂`$. In particular the relevant holomorphic 3-form is given by $`\frac{dx_1dx_2du}{x_1x_2u}`$ (by eliminating $`z,v`$ and noting that the denominator has $`uv/v`$). To better understand the geometry of $`\widehat{}`$ we rewrite the defining equation eq. (8.11) in the following form,
$`h(x_1,x_2,z)`$ $`:=`$ $`x_1^2x_2+x_1x_2^2+1+zx_1x_2=0`$ (8.12)
$`ze^{t/3}`$ $`=`$ $`uv.`$ (8.13)
where we have rescaled variables and shifted $`z`$. In this form the holomorphic 3-form becomes
$$\mathrm{\Omega }=\frac{dx_1dx_2du}{x_1x_2u}=\frac{dx_1dx_2}{h/z}\frac{du}{u}=\mathrm{\Omega }_2\frac{du}{u}$$
(8.14)
The first equation in eq. (8.12) defines an elliptic fibration (in terms of the $`x_1,x_2`$ variables over the complex plane with coordinate $`z`$). Moreover the corresponding two form $`\mathrm{\Omega }_2`$ is the same as would be for a K3 geometry where $`x_1,x_2`$ are now viewed as variables in $`𝐂`$. It is more convenient to homogenize the above elliptic curve by introducing an extra variable $`x_0`$:
$$x_1^2x_2+x_1x_2^2+x_0^3+zx_0x_1x_2=0$$
(8.15)
We can convert this into the Weierstrass form by the following coordinate transformation,
$`x_1=Y+\frac{U}{2}+z\frac{X}{2},x_2=Y+\frac{U}{2}+z\frac{X}{2},x_0=X,`$ (8.16)
$`UY^2=X^3+(\frac{z}{2})^2UX^2+\frac{z}{2}U^2X+\frac{1}{4}U^3,`$ (8.17)
where now $`U`$ is a scaling variable and we can set it to $`1`$. Thus the BPS data of $`\widehat{}`$ is the same as that of
$$Y^2=X^3+(\frac{z}{2})^2X^2+\frac{z}{2}X+\frac{1}{4},ze^{t/3}=uv,$$
(8.18)
with the holomorphic 3-form
$$\mathrm{\Omega }=\frac{dX}{Y}dz\frac{du}{u}.$$
(8.19)
Thus we see that $`\widehat{}`$ can be viewed as the product of an elliptic fibration times a $`𝐂^{}`$ fibration over $`z`$. The elliptic fibration has a discriminant given by $`\mathrm{\Delta }(z)=\frac{1}{16}(27z^3)`$, the three degenerate fibers are located symmetrically at $`(27)^{\frac{1}{3}}\{1,e^{2\pi i/3},e^{4\pi i/3}\}=:\{z_1,z_2,z_3\}`$. The $`𝐂^{}`$ fibration has one degenerate fiber given at $`z=e^{t/3}`$.
To determine supersymmetric 3-cycles in $`\widehat{}`$ we use the circle of $`𝐂^{}`$ fibration as one cycle times an additional 2-cycle. We note that there are only two closed 2-cycles in this fibration. One is the elliptic fiber of the fibration itself and the other closed 2-cycle is formed by taking a path in the base that encloses $`z_i`$ and the $`(1,0)`$ cycle of the elliptic fiber above this path.<sup>2</sup><sup>2</sup>2 The existence of the 2nd type of closed 2-cycle is precisely the reason this fibration can be used to construct 5-brane web description of 5D $`E_0`$ field theory which can also be obtained via geometric engineering from the CY threefold $``$ . The fact that there are no other closed 2-cycles is not obvious even though we do not have degenerating cycles of the same charge, since a cycle starting from a degenerate fiber can undergo monodromy transformations when it goes around other degenerating fibers. One can, however, construct 2-cycles with boundaries such that the boundary is a 1-cycle of an elliptic fiber above the point $`z_{}`$. We can construct closed 3-cycles using the 1-cycle of the $`𝐂𝖨^{}`$ fibration and the 2-cycles (with boundaries) if we choose the position of the boundary carefully. If $`z_{}=e^{\frac{t}{3}}`$ then the boundary of the 2-cycle is exactly at the point on the base where the $`𝐂𝖨^{}`$ fibration degenerates. In this case the 2-cycle and the 1-cycle of the $`𝐂𝖨^{}`$ fibration together define a closed 3-cycle . Since the 1-cycle of the $`𝐂𝖨^{}`$ fibration is always present the essential geometry of the 3-cycle is captured by the 2-cycle with boundary at $`z_{}=e^{\frac{t}{3}}`$. This gets translated to finding minimal 2 surfaces in the corresponding $`K3`$ geometry with boundary being a circle on a particular elliptic fiber.
The connection with F-theory is now rather clear. In fact this elliptic fibration defines an F-theory background studied before in the context of non-BPS stable states in F-theory and compactification of 5D $`E_n`$ field theories on a circle . As usual in the F-theory description, we can assign $`(p,q)`$ charges to the 7-branes which correspond to $`(p,q)`$ degenerating cycle of $`T^2`$ ( which can be defined by choosing paths to a base point, $`z_0`$). In this case the charges are as shown in Fig. 45 .Note that the $`(p,q)`$ charge of the three degenerating cycles can be cyclically transformed by the $`SL(2,𝖹𝖹)`$ matrix $`ST`$, $`(ST)^3=1`$ .
The relation with the probe theory and its BPS states is now clear. The degenerating fibers at $`z_i`$ define a 7-brane background and the degenerating 1-cycle at $`z_{}=e^{\frac{t}{3}}`$ of the $`𝐂𝖨^{}`$ fibration defines the position of the D3-brane. The 2-cycles with boundary are strings and string junctions stretched between the 7-branes and the D3-brane. BPS states in the D3-brane theory correspond to BPS string junctions which are the projections of holomorphic 2-cycles. Thus we see that D-branes wrapped on even dimensional cycles of $``$ are mirror to states in the D3-brane worldvolume theory. This connection between sheaves on a non-compact CY and states in the field theory realized on a D3-brane in the presence of 7-branes was also studied in .
#### 8.1.3 Soliton Numbers for $`\text{I}\text{P}^2`$ and its Blowups
As noted before, the elliptic fibration defined by eq. (8.12) is exactly the W-plane geometry of the massive LG theory mirror to sigma model with $`\text{I}\text{P}^2`$ target space. The vanishing cycles are just the cycles of the $`T^2`$ fiber, and so we can use the knowledge of the degeneration types to find the intersection number of vanishing cycles, and thus the soliton numbers of this theory. From Fig. 45 it follows that,
$$\gamma _i\gamma _j=\left(\begin{array}{ccc}0& 3& 3\\ 0& 0& 3\\ 0& 0& 0\end{array}\right).$$
(8.20)
Where we have written the intersection matrix as an upper triangular matrix. Similarly, as mentioned before, we can consider other non-compact CY threefolds $`_n`$ which are the total space of the canonical line bundle over the toric del Pezzo, $`_n`$. As for the case of $`\text{I}\text{P}^2`$, in this case as well the mirror is an elliptic fibration and a $`𝐂𝖨^{}`$ fibration over the z-plane. The elliptic fibration is defined by the superpotential of the corresponding massive LG theory. The corresponding F-theory backgrounds and the D3-brane theory were studied in . Since the charges of the vanishing cycles are known for these cases we can compute the soliton counting matrix and compare with the matrices obtained from the collection of exceptional bundles. In the following we will denote the superpotential of the LG theory mirror to sigma model on $`X`$ as $`W_X`$.
$`_1`$: There are four degenerate fibers in this case as shown in Fig. 46.
Three out of four cycles are the same as before. The new degenerate fiber has charge $`(1,1)`$. The intersection matrix is,
$$\gamma _i\gamma _j=\left(\begin{array}{cccc}0& 1& 3& 3\\ 0& 0& 1& 2\\ 0& 0& 0& 3\\ 0& 0& 0& 0\end{array}\right).$$
(8.21)
One can check that this matrix produces correct Ramond charges for the chiral fields.
$`_2`$: In this case there are five degenerate fibers as shown in Fig. 47. There are two mutually local (with the same charge) fibers.
The intersection matrix is given by
$$\gamma _i\gamma _j=\left(\begin{array}{ccccc}0& 1& 1& 3& 3\\ 0& 0& 0& 1& 2\\ 0& 0& 0& 1& 2\\ 0& 0& 0& 0& 3\\ 0& 0& 0& 0& 0\end{array}\right).$$
(8.22)
$`_3`$: The six degenerate fibers in this case are shown in Fig. 48. In this case there are three mutually local fibers.
The intersection matrix is given by
$$\gamma _i\gamma _j=\left(\begin{array}{cccccc}0& 1& 1& 1& 3& 3\\ 0& 0& 0& 0& 1& 2\\ 0& 0& 0& 0& 1& 2\\ 0& 0& 0& 0& 1& 2\\ 0& 0& 0& 0& 0& 3\\ 0& 0& 0& 0& 0& 0\end{array}\right).$$
(8.23)
$`F_0`$: The four degenerate fibers in this case are shown in Fig. 49.
The intersection matrix is given by
$$\gamma _i\gamma _j=\left(\begin{array}{cccc}0& 2& 0& 2\\ 0& 0& 2& 0\\ 0& 0& 0& 2\\ 0& 0& 0& 0\end{array}\right).$$
(8.24)
One can check that these matrices give the correct Ramond charge for the chiral fields.
## Acknowledgement
We would like to thank M. Aganagic, C. Bachas, A. Chari, S. Coleman, M. Gutperle, R. Myers, A. Polishchuk, R. Thomas, S.T. Yau and E. Zaslow for valuable discussions. AI would also like to thank Asad Naqvi for useful discussions. K.H. would like to thank McGill University for hospitality. The research of K.H. is supported in part by NSF-DMS 9709694. The research of A.I. is supported in part by the US Department of Energy under contract #DE-FC02-94ER40818. The research of C.V. is supported in part by NSF grant PHY-98-02709.
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# References
Universality and individuality in a neural code
Elad Schneidman,<sup>1,2</sup> Naama Brenner,<sup>3</sup> Naftali Tishby,<sup>1,3</sup>
Rob R. de Ruyter van Steveninck,<sup>3</sup> and William Bialek<sup>3</sup>
<sup>1</sup>Institute for Computer Science,
Center for Neural Computation and
<sup>2</sup>Department of Neurobiology
Hebrew University
Jerusalem 91904, Israel
<sup>3</sup>NEC Research Institute
4 Independence Way
Princeton, New Jersey 08540, USA
elads/tishby@cs.huji.ac.il
brenner/ruyter/bialek@research.nj.nec.com
The problem of neural coding is to understand how sequences of action potentials (spikes) are related to sensory stimuli, motor outputs, or (ultimately) thoughts and intentions. One clear question is whether the same coding rules are used by different neurons, or by corresponding neurons in different individuals. We present a quantitative formulation of this problem using ideas from information theory, and apply this approach to the analysis of experiments in the fly visual system. We find significant individual differences in the structure of the code, particularly in the way that temporal patterns of spikes are used to convey information beyond that available from variations in spike rate. On the other hand, all the flies in our ensemble exhibit a high coding efficiency, so that every spike carries the same amount of information in all the individuals. Thus the neural code has a quantifiable mixture of individuality and universality.
Introduction
When two people look at the same scene, do they see the same things? This basic question in the theory of knowledge seems to be beyond the scope of experimental investigation. An accessible version of this question is whether different observers of the same sense data have the same neural representation of these data: how much of the neural code is universal, and how much is individual? To approach this problem we must give a quantitative definition of similarity or distance among neural codes.
The problem of comparing neural codes has some analogies to the problem of comparing among amino acid sequences in proteins. In the neural code, sequences of action potentials stand for sensory inputs or motor outputs; in the genetic code sequences of nucleotides stand for amino acids which in turn encode the three dimensional structure of proteins. Just as motifs of several amino acids collectively can encode structural elements in proteins, patterns of several action potentials across time or across a population of cells can have a special meaning. For proteins, sequences can be similar because they have a recent common ancestor; alternatively, sequences can be equivalent functionally even if they are distinguishable, since convergent evolution has led to very different sequences that encode proteins with similar structures and functions. For the neural code we show that these notions of distinguishability and functional equivalence can be quantified using ideas from information theory . In particular this approach does not require a metric either on the space of stimuli or on the space of neural responses; all notions of similarity emerge from the statistical structure of the neural responses.
We apply these methods to analyze experiments on an identified motion sensitive neuron in the fly’s visual system, the cell H1 . Many invertebrate nervous systems have cells that can be named and numbered , and in many cases the total number of neurons involved in representing a portion of the sensory world is quite small, so that destruction of individual neurons can have a substantial impact on behavior (see, for example, Ref. ). In these cases the neural representation of sensory information is especially accessible, precisely because it is localized to a small set of identified cells. On the other hand, if a large fraction of neurons is identifiable it might seem that the question of whether different individuals share the same neural representation of the visual world would have a trivial answer.
Far from trivial, we shall see that the neural code even for identified neurons in flies has components which are common among flies and significant components which are individual to each fly. The existence of identified neurons thus does not preclude the expression of individuality in neural representations; we should expect that all neural circuits, both vertebrate and invertebrate, express a degree of universality and a degree of individuality. For H1 we quantify these ideas, and we hope that the methods we introduce will be applicable more generally. In the interest of making the discussion more accessible we have confined mathematical arguments to the Methods section; in the main text we make use of analogies to provide some intuition for what the information theoretic quantities are measuring.
Distinguishing among flies
We place our discussion in the context of the experiments shown in Fig. 1. Nine different flies are shown precisely the same movie, which is repeated many times for each fly; as we show the movie we record the action potentials from the H1 neuron. The details of the stimulus movie should not have a qualitative impact on the results, provided that the movie that is sufficiently long and rich to drive the system through a reasonable and natural range of responses. Figure 1b makes clear that qualitative features of the neural response on long time scales ($`100\mathrm{ms}`$) are common to almost all the flies, and that some aspects of the response are reproducible on a (few) millisecond time scale across multiple presentations of the movie to each fly. Nonetheless the responses are not identical in the different flies, nor are they perfectly reproduced from trial to trial in the same fly. The most obvious difference among the flies is in the average spike rate, which varies from 22 to 63 spikes/s among our ensemble of flies. But beyond that, how should we quantify the similarity or difference among the neural responses?
One way is to imagine that each spike train is a point in an abstract space, and that there is a metric on this space. Considerable effort has gone into the definition of metrics that are plausible biologically and tractable computationally , and these methods have been used widely (for applications in neurobiology (see ), but all approaches based on metric spaces have several problems. First, the metric is imposed by the investigator and does not emerge from the data. Second, even within a plausible class of metrics there are arbitrary parameters, such as the relative distance cost of moving vs. deleting a spike. Finally, it is not clear that our intuitive notion of similarity among neural responses (or amino acid sequences) is captured by the mathematical concept of a metric.
In contrast, information theory provides a method for quantifying directly the differences among the sources of the spike trains. Imagine that we record multiple speakers reading from the same text, in the same way that we record the activity of neurons from different flies responding to the same sensory inputs. There are many possible speakers, and we are shown a small sample of the speech signal: how well can we identify the speaker? If we can collect enough data to characterize the distribution of speech sounds made by each speaker then we can quantify, in bits, the average amount of information that a segment of speech gives about the identity of the speaker. Further we can decompose this information into components carried by different features of the sounds. Following this analogy, we will measure the information that a segment of the neural response provides about the identity of the fly, and we will ask how this individuality is distributed across different features of the spike train. Apparently large differences in spike rate are surprisingly uninformative, but the temporal patterns of spikes allow for much more efficient discrimination among individuals. The largest differences among individuals are in the way that these patterns are associated with specific stimulus features.
We discretize the neural response into time bins of size $`\mathrm{\Delta }t=2\mathrm{ms}`$, which is also the time required to draw one frame of the stimulus movie. At this resolution there are almost never two spikes in a single bin, so we can think of the neural response as a binary string, as in Fig. 1c-d. We examine the response in blocks or windows of time having length $`T`$, so that an individual neural response becomes a binary ‘word’ $`W`$ with $`T/\mathrm{\Delta }t`$ ’letters’. Clearly, any fixed choice of $`T`$ and $`\mathrm{\Delta }t`$ is arbitrary, and so we explore a range of these parameters.
The distribution of words used by a particular fly’s H1 in response to the stimulus movie, $`P^\mathrm{i}(W)`$ for the $`\mathrm{i}^{\mathrm{th}}`$ fly, tells us about the ‘vocabulary’ of that cell. Figure 1f shows that different flies ‘speak’ with similar but distinct vocabularies. From these distributions $`P^\mathrm{i}(W)`$ we can quantify the average information that a single word of length $`T`$ gives about the identity of the fly, $`I(W\mathrm{identity};T)`$ \[see Eq. (1) in Methods\]. Thus, we measure how well we can discriminate between one individual and a mixture of all the other individuals in the ensemble, or effectively how ‘far’ each individual is from the mean of her conspecifics.
The finite size of our data set prevents us from exploring arbitrarily long words, but happily we find that information about identity is accumulating at more or less constant rate $`R`$ well before the undersampling limits of the experiment are reached (Fig. 2a). Thus $`I(W\mathrm{identity};T)R(W\mathrm{identity})T`$; $`R(W\mathrm{identity})5`$ bits/s, with a very weak dependence on the time resolution $`\mathrm{\Delta }t`$. Since the mean spike rate can be measured by counting the number of 1s in each word $`W`$, this information includes the differences in firing rate among the different flies.
Even if flies use very similar ‘vocabularies,’ they may differ substantially in the way that they associate words with particular stimulus features. In our experiments the stimulus runs continuously in a loop, so that we can specify the stimulus precisely by giving the time relative to the start of the loop; in this way we don’t need to make any assumptions about which features of the stimulus are important for the neuron, nor do we need a metric in the space of stimuli. We can therefore consider the word $`W`$ that the $`\mathrm{i}^{\mathrm{th}}`$ fly will generate at time $`t`$. This word is drawn from the distribution $`P^\mathrm{i}(W|t)`$ which we can sample, as in Fig. 1c–e, by looking across multiple presentations of the same stimulus movie. In parallel with the discussion above, we can now ask for the average information that a word $`W`$ provides about identity given that it was observed at a particular time $`t`$. This depends on the time $`t`$ because some moments in the stimulus are more informative than others, as is obvious from Fig. 1. The more natural quantity is an average over all times $`t`$, which is the average information that we can gain about the identity of the fly by observing a word $`W`$ at a known time $`t`$ relative to the stimulus, $`I(\{W,t\}\mathrm{identity};T)`$ \[see Eq. (5) in Methods\].
Figure 2b shows a plot of $`I(\{W,t\}\mathrm{identity};T)/T`$ as a function of the observation time window of size $`T`$. Observing both the spike train and the stimulus together provides $`32\pm 1\mathrm{bits}/\mathrm{s}`$ about the identity of the fly. This is more than six times as much information as we can gain by observing the spike train alone, and corresponds to gaining one bit in $`30\mathrm{ms}`$. Correspondingly, a typical pair of flies in our ensemble can be distinguished reliably in $`30\mathrm{ms}`$. This is the time scale on which flies actually use their estimates of visual motion to guide their flight during chasing behavior , so that the neural codes of different individuals are distinguishable on the time scales relevant to behavior.
Spike rates and information rates
Having seen that we can distinguish reliably among individual flies using relatively short samples of the neural response, it is natural to ask about the origins and implications of these individual differences. Perhaps the most obvious question is whether the substantial differences in the code among the different neurons have an impact on the ability of these cells to convey information about the visual stimulus. As discussed in Refs. , the rate at which the neural response provides information about the visual stimulus, $`R^\mathrm{i}(Ws(t);T)`$, is determined by the same probability distributions $`P^\mathrm{i}(W|t)`$ as before \[see Eq’s. (69) in Methods\]. Again we note that our estimate of the information rate itself is independent of any metric in the space of stimuli, nor does it depend on assumptions about which stimulus features are most important in the code.
Figure 3a shows that the flies in our ensemble span a range of information rates from $`R^\mathrm{i}(Ws(t))50`$ to $`150\mathrm{bits}/\mathrm{s}`$. This threefold range of information rates is correlated with the range of spike rates, so that each of the cells transmits nearly a constant amount of information per spike, $`2.39\pm 0.24\mathrm{bits}/\mathrm{spike}`$. The error bar in this case ($`\pm 0.24\mathrm{bits}/\mathrm{spike}`$) is a standard deviation across the ensemble of flies, not a standard error of the mean: the number of bits per spike transmitted by H1 (under these stimulus conditions) is constant from fly to fly within $`10\%`$, despite three fold variations in total spike rate.
Although information rates are correlated with spike rates, this does not mean that information is carried by a “rate code” alone. Rate coding usually is distinguished from “timing codes” in which the detailed temporal structure of the spike train plays a crucial role. In particular, our computation of the information carried by the spike train includes automatically any contribution from temporal patterns, but to demonstrate that these patterns are important we must show that this information is more than we expect just by summing the contributions of the individual spikes that make up the patterns. This ‘single spike information’ can also be thought of as the information conveyed by temporal modulations in the spike rate; see Ref. and Eq. (11) in Methods. For all the flies in our ensemble, the total rate at which the spike train carries information is substantially larger than the ‘single spike’ information—$`2.39`$ vs. $`1.64`$ bits/spike, on average. This extra information, as defined in Eq. (12) and illustrated in Fig. 3b, is carried in the temporal patterns of spikes.
The fact that the information per spike is constant across the ensemble of flies means that cells with higher spike rates are not generating extra spikes at random, but rather each extra spike is equally informative about the visual stimulus. The capacity of the code to carry information is quantified by the entropy rate $`𝒮_{\mathrm{total}}^\mathrm{i}`$ of the distribution of neural responses \[see Eq’s. (6,10) in Methods\], and is different in different flies. It is natural to define the efficiency of the code as the fraction of this capacity which is used to convey information about the visual stimulus, $`ϵ^\mathrm{i}=R^\mathrm{i}(Ws(t))/𝒮_{\mathrm{total}}^\mathrm{i}`$. Like the information per spike, this efficiency is nearly constant across the ensemble of flies, $`ϵ=0.59\pm 0.05`$, at $`\mathrm{\Delta }t=2\mathrm{ms}`$, with a very weak dependence on $`\mathrm{\Delta }t`$ .
A universal codebook?
Even though flies differ in the structures of their neural responses, distinguishable responses could be functionally equivalent, as with distinct amino acid sequences that fold to the same protein structure. It might therefore be that all flies could be endowed (genetically?) with a universal or consensus codebook that allows each individual to make sense of her own spike trains, despite the differences from her conspecifics. Thus we want to ask how much information we lose if the identity of the flies is hidden from us, or equivalently how much each fly can gain by knowing its own individual code.
The codebook for any individual fly can be thought of as a probabilistic mapping from neural responses or words back into the space of visual stimuli . The information conveyed by the spike train quantifies the specificity of this mapping: the ‘tighter’ the distribution of stimuli consistent with a given response the more information is conveyed. If the neural codes used by different flies are different, then these conditional distributions in stimulus space are also different. If we don’t know the identity of the fly, all we can do is to associate each neural response with a distribution of stimuli that corresponds to an average over the individuals, and this distribution necessarily is broader than any of the individual distributions. As a result, we have less information about the visual stimulus, as summarized by Eq. (14) in the Methods.
Intuitively, the greater the differences among the neural responses of different flies, the more visual information we will lose if we don’t know the identity of the individual. On the other hand, these differences mean that the neural response provides information about individual identity, so that information gained about identity is information lost about the stimulus if we use a universal codebook. This intuitive connection is made precise by Eq. (15) in the Methods. As a practical matter, this means that the answer to our question about the efficacy of a universal decoder is contained in the results of Fig. 2. The result is that, on average, not knowing the identity of the fly limits us to extracting only $`64`$ bits/s of information about the visual stimulus. This should be compared with the average information rate of $`92.3`$ bits/s in our ensemble of flies: knowing her own identity allows the average fly to extract $`44\%`$ more information from H1. Further analysis shows that each individual fly gains approximately the same relative amount of information from knowing its personal codebook.
The nature of the ‘personal’ bits
Thus far we have analyzed the differences among the neural codes of different flies, and how much extra information a fly can extract by knowing it’s individual codebook. It is natural to ask what is being “said” by these extra bits, characterizing more explicitly the mapping from neural responses back to stimulus space for the different flies.
For each neural response $`W`$ we can look back through the entire experiment and accumulate the motion trajectories that lead up to the response, and these provide samples from the distribution of stimuli conditional on the response as described above. Because the space of trajectories has many dimensions, this distribution is difficult to visualize, and so we focus here on the means of these distributions. This is a generalization of the reverse correlation or spike triggered average method : rather than looking at the average stimulus that leads to a single spike, we look at the average stimulus that leads to the responses $`W`$, which can consist of a pattern of spikes and empty intervals .
In Fig. 4 we show the average waveforms of the stimulus velocity preceding a specific binary word in the spike trains of flies 1 and 6. As fly 6 spike trains convey almost 3 times more information about the stimulus, one might have speculated that the same word was used in completely different stimulus contexts for the two flies. In fact, the differences are in the details and not in the general picture: spikes stand for pulses of positive velocity (as in Fig. 4b), long silent intervals stand for negative velocities (as in Figs. 4a&c), and the largest differences among the flies are in the widths, latencies and amplitudes of the pulses; combinations of spikes and intervals then lead to very different trajectories (as in Fig. 4d). For the fly which conveys less information, spikes are associated with larger positive velocities (Fig. 4b) and silences are associated with (slightly) larger negative velocities (Fig. 4a); thus, these elementary responses come closer to exhausting the dynamic range of the inputs. Conversely, the more informative spike train covers the dynamic range of inputs with a greater variety of composite responses.
Discussion
One obvious difference between invertebrate and vertebrate nervous systems is the existence of identified neurons in invertebrates. The identifiability of invertebrate neurons sometimes has been interpreted to mean that these smaller nervous systems are hard wired automata; indeed the optomotor system of flies has been held up as a clear example of this extreme view. In this view, individuality plays no role, and it should even be possible to average the results of experiments on corresponding neurons in different individuals. For vertebrates, substantial individuality arises through development and learning, and there are few if any identified neurons; at best vertebrates have identifiable modules consisting of hundreds or thousands of neurons, such as the columns in visual cortex. Against this clear dichotomy it is worth remembering that even genetically identical single celled organisms exhibit individuality in their sensory–motor behavior .
In the present work we have tried to the quantify the individuality of the neural code used by a single neuron in the fly visual system. On the one hand, this individuality is sufficient to allow discrimination among individuals on time scales of relevance to behavior. Correspondingly, each individual fly would lose a significant amount ($`30\%`$) of the visual information carried by this neuron if it ‘knew’ only the codebook appropriate to the whole ensemble of flies. On the other hand, these differences among the codebooks of different flies seem to be matters of detail. Although different flies extract very different amounts of information from the same visual inputs, all the flies achieve a high and constant efficiency in their encoding of this information. From previous work it is known that the visual system of an individual fly exhibits substantial changes in coding strategy as it adapts to different ensembles of inputs. Rather than converging on the same information rates in different flies, these adaptation processes seem to converge on codes of uniformly high efficiency, supporting the idea that efficiency of representation is a ‘design principle’ for the system .
On average the flies in our ensemble have neural codes in which a substantial amount of information is carried by patterns of spikes. This antiredundancy or synergy among spikes is reduced substantially if we are forced to use a universal codebook. Mathematically this loss of synergy in the universal codebook is related to the fact that the rate at which we gain information about the identity of the fly (Fig. 2b) increases with window size out to $`T_\mathrm{c}10\mathrm{ms}`$: discrimination among flies is enhanced by being able to see patterns of spikes in windows of size $`T_\mathrm{c}`$, implying that the way these patterns are used to encode visual information is unique to each individual. Each individual fly thus gains nearly $`50\%`$ more information through the use of a code in which patterns of spikes carry extra information, and more than half of this is lost if the fly does not have knowledge of its own identity. Not only is spike timing important for the neural code, but the way in which timing is used is specific to each individual.
Methods
Flies, neural recording and stimulus generation
Recordings were made from the H1 neuron using standard methods: the fly was immobilized in wax, a tungsten microelectrode was inserted through a small hole at the back of the fly’s head, and H1 was identified through its response properties; spikes were detected with a window discriminator. The stimulus was a rigidly moving pattern of vertical bars, randomly dark or bright, with average intensity $`\overline{I}100\mathrm{m}\mathrm{W}/(\mathrm{m}^2\mathrm{sr})`$, displayed on a Tektronix 608 high brightness display; bar widths were set equal to the horizontal lattice spacing (interommatidial angle) of the compound eye. The fly viewed the display through a round diaphragm, showing approximately 30 bars. Frames of the stimulus pattern were refreshed every 2 ms, and with each new frame the pattern was displayed at a new position. This resulted in an apparent horizontal motion of the bar pattern, which is suitable to excite the H1 neuron. The pattern position was defined by a pseudorandom sequence, simulating a diffusive motion or random walk. We draw attention to three points relevant for the present analysis: (1) The flies are freshly caught female Calliphora, so that our ‘ensemble of flies’ approaches a natural ensemble and is not restricted to a highly inbred laboratory stock. (2) In each fly we identify the H1 cell as the unique spiking neuron in the lobula plate that has a combination of wide field sensitivity, inward directional selectivity for horizontal motion, and contralateral projection. (3) Recordings are rejected only if raw electrode signals are excessively noisy or unstable; in particular we do not select for flies that exhibit mean spike rates (spontaneous or driven) in a predefined range.
Definition of $`I(W\mathrm{identity};T)`$
Imagine that we record the response of the H1 neuron from one fly, but we don’t know which one. A priori there are $`N`$ equally likely possibilities. Once we observe the spike train for some time $`T`$ our uncertainty is reduced, and hence we gain information about the identity of the fly from which we are recording. The average information that an individual word provides about the fly’s identity is
$$I(W\mathrm{identity};T)=\underset{\mathrm{i}=1}{\overset{N}{}}P_\mathrm{i}\underset{W}{}P^\mathrm{i}(W)\mathrm{log}_2\left[\frac{P_\mathrm{i}P^\mathrm{i}(W)}{P^{\mathrm{ens}}(W)}\right]\mathrm{bits},$$
(1)
where $`P_\mathrm{i}=1/N`$ is the a priori probability that we are recording from fly $`\mathrm{i}`$ and $`P^{\mathrm{ens}}(W)`$ is the probability that any fly in the whole ensemble of flies would generate the word $`W`$,
$$P^{\mathrm{ens}}(W)=\underset{\mathrm{i}=1}{\overset{N}{}}P_\mathrm{i}P^\mathrm{i}(W).$$
(2)
The measure $`I(W\mathrm{identity};T)`$ has been discussed by Lin as the ‘Jensen–Shannon divergence’ $`D_{\mathrm{JS}}`$ among the distributions $`P^\mathrm{i}(W)`$. We recall that the problem of finding a measure of similarity among distributions is not simple; obvious choices such as the Kullback–Leibler divergence are not symmetric, and may have spurious technical requirements such as absolute continuity of one distribution with respect to the others. Lin proposed $`D_{\mathrm{JS}}`$ as a way of getting around these difficulties, and he showed that $`D_{\mathrm{JS}}`$ can be used to bound other measures of similarity, such as the optimal or Bayesian probability of identifying correctly the origin of a sample (as in forced choice psychophysical discrimination experiments). Here $`D_{\mathrm{JS}}`$ arises not just as an interesting possible measure of similarity (see also ), but as the unique answer to the question of how much information a sample provides about its source.
Definition of $`I(\{W,t\}\mathrm{identity};T)`$
By analogy with Eq. (2) we define the distribution of words used at time $`t`$ by the whole ensemble of flies.
$$P^{\mathrm{ens}}(W|t)=\underset{\mathrm{i}=1}{\overset{N}{}}P_\mathrm{i}P^\mathrm{i}(W|t),$$
(3)
and by analogy with Eq. (1) we can measure the information that the word $`W`$ observed at known $`t`$ gives us about the identity of the fly,
$$I(W\mathrm{identity}|t;T)=\underset{\mathrm{i}=1}{\overset{N}{}}P_\mathrm{i}\underset{W}{}P^\mathrm{i}(W|t)\mathrm{log}_2\left[\frac{P_\mathrm{i}P^\mathrm{i}(W|t)}{P^{\mathrm{ens}}(W|t)}\right].$$
(4)
The more natural quantity is an average over all times $`t`$,
$$I(\{W,s(t)\}\mathrm{identity};T)=I(W\mathrm{identity}|t;T)_t\mathrm{bits},$$
(5)
where $`\mathrm{}_t`$ denotes an average over $`t`$.
Definition of $`I^\mathrm{i}(Ws(t);T)`$
The discussion here follows Refs. . The entropy of the distribution of words,
$$S_{\mathrm{total}}^\mathrm{i}(T)=\underset{W}{}P^\mathrm{i}(W)\mathrm{log}_2P^\mathrm{i}(W)\mathrm{bits},$$
(6)
measures the capacity of the neuron to transmit information. At each time $`t`$ we can define the entropy of the conditional distribution $`P^\mathrm{i}(W|t)`$, which measures the noise in the response to repeated presentations of the same movie,
$$S_{\mathrm{noise}}^\mathrm{i}(t;T)=\underset{W}{}P^\mathrm{i}(W|t)\mathrm{log}_2P^\mathrm{i}(W|t)\mathrm{bits}.$$
(7)
The information is the difference between the total entropy of the cell’s vocabulary and the average noise entropy,
$$I^\mathrm{i}(Ws(t);T)=S_{\mathrm{total}}^\mathrm{i}(T)S_{\mathrm{noise}}^\mathrm{i}(t;T)_t.$$
(8)
For sufficiently large window of time $`T`$, we expect that to gain information in proportion to the duration of our observations,
$$I^\mathrm{i}(Ws(t);T)R^\mathrm{i}(Ws(t))T,$$
(9)
so that there is a well defined information rate $`R^\mathrm{i}(Ws(t))`$. This asymptotic behavior is observed in the data for values of $`T`$ that are relevant to fly behavior. Similar behavior is observed for the total entropy,
$$S_{\mathrm{total}}^\mathrm{i}(T)𝒮_{\mathrm{total}}^\mathrm{i}T,$$
(10)
leading to the definition of coding efficiency discussed in the text .
Information from single spikes
The discussion here follows Ref. . In the experiments the stimulus runs continuously for a time $`T_{\mathrm{loop}}`$ and then repeats. When we observe a single spike at time $`t`$ we learn something about the stimulus in the neighborhood of this time, and if we average this information over all possible arrival times we obtain the average information carried by a single spike. Because information is mutual, we can relate the information that the spike provides about the stimulus to the information that the stimulus provides about the occurrence of a spike, but this is contained in the time dependent firing rate or post–stimulus time histogram for cell $`\mathrm{i}`$, $`r_\mathrm{i}(t)`$. After some algebra , the single spike information takes the form of an integral that depends only on $`r_\mathrm{i}(t)`$,
$$I_{\mathrm{one}\mathrm{spike}}^\mathrm{i}=\frac{1}{T_{\mathrm{loop}}}_0^{T_{\mathrm{loop}}}𝑑t\frac{r_\mathrm{i}(t)}{\overline{r}_\mathrm{i}}\mathrm{log}_2\left[\frac{r_\mathrm{i}(t)}{\overline{r}_\mathrm{i}}\right]\mathrm{bits},$$
(11)
where $`\overline{r}_\mathrm{i}`$ is the average spike rate in cell $`\mathrm{i}`$. If spikes were to carry information independently, then each cell would transmit $`R_{\mathrm{ind}}^\mathrm{i}=\overline{r}_\mathrm{i}I_{\mathrm{one}\mathrm{spike}}^\mathrm{i}`$ bits per second. If the total information rate $`R^\mathrm{i}(Ws(t))`$ is smaller than this then spikes are redundant (on average) while if the total information rate is larger then there is synergy among the spikes and the extra information must be carried in the temporal patterns of spikes. We can quantify this extra information as a fraction,
$$F_{\mathrm{extra}}^\mathrm{i}=[R^\mathrm{i}(Ws(t))R_{\mathrm{ind}}^\mathrm{i}]/R_{\mathrm{ind}}^\mathrm{i},$$
(12)
as illustrated in Fig. 3b.
Information loss with universal decoding
If we observe the response of a neuron but don’t know the identity of the individual generating this response, then we are observing responses drawn from the ensemble distributions defined above, $`P^{\mathrm{ens}}(W|t)`$ and $`P^{\mathrm{ens}}(W)`$. The information that words provide about the visual stimulus then is
$$I^{\mathrm{mix}}(Ws(t);T)=\underset{W}{}P^{\mathrm{ens}}(W|t)\mathrm{log}_2\left[\frac{P^{\mathrm{ens}}(W|t)}{P^{\mathrm{ens}}(W)}\right]_t\mathrm{bits}.$$
(13)
On the other hand, if we know the identity of the fly to be $`\mathrm{i}`$, we gain the information $`I^\mathrm{i}(Ws(t);T)`$ from above \[Eq’s. (68)\]. The average information loss is then
$$I_{\mathrm{loss}}^{\mathrm{avg}}(Ws(t);T)=\underset{\mathrm{i}=1}{\overset{N}{}}P_\mathrm{i}I^\mathrm{i}(Ws(t);T)I^{\mathrm{mix}}(Ws(t);T).$$
(14)
After some algebra it can be shown that this average information loss is related to the information that the neural responses give about the identity of the individuals, as defined above:
$`I_{\mathrm{loss}}^{\mathrm{avg}}(Ws(t);T)`$ $`=`$ $`I(\{W,t\}\mathrm{identity};T)`$ (15)
$`I(W\mathrm{identity};T).`$
Figure 1. Different flies’ spike trains and word statistics. (a) All flies view the same random vertical bar pattern moving across their visual field with a time dependent velocity, part of which is shown (see Methods section for details). In the experiment, a 40 sec waveform is presented repeatedly, 90 times. (b) A set of 45 response traces to the part of the stimulus shown in (a) from each of the 9 flies. The traces are taken from the segment of the experiment where the transient responses have decayed. Spike trains from flies 1 and 6 are colored by red and blue, respectively, which we will use as a color code for the other parts of the figure.(c) Example of construction of the local word distributions. Zooming in on a segment of the repeated responses of fly 1 to the visual stimuli (see green rectangle in (b)), the fly’s spike trains are divided into contiguous 2 ms bins, and the spikes in each of the bins are counted. E.g., we get the 6 letter words that the fly used at time 3306 ms into the input trace. (d) Similar as (c) for fly 6. (e) The distributions of words that flies 1 and 6 used at time $`t=3306\mathrm{ms}`$ from the beginning of the stimulus. The time dependent distributions, $`P^1(W|t=3306\mathrm{ms})`$ and $`P^6(W|t=3306\mathrm{ms})`$ are presented as a function of the binary value of the actual ’word’, e.g., binary word value $`{}_{}{}^{}17_{}^{}`$ stands for the word $`{}_{}{}^{}010001_{}^{}`$. (f) Collecting the words that each of the flies used through all of the visual stimulus presentations, we get the total word distributions for flies 1 and 6, $`P^1(W)`$ and $`P^6(W)`$.
Figure 2. Distinguishing one fly from others based on spike trains. (a) The average rate of information gained about the identity of a fly, given the distribution of words that it used throughout the stimulus presentations, as a function of the word size used. The information rate is saturated even before we reach the maximal word length used; for more discussion of word lengths see Methods. Following Figure 1, Red marks are the average rate of information that the word distribution of fly 1 give about its identity, compared with the word distribution mixture of all of the flies. The connecting line is used for clarification only. Blue marks the results for fly 6, and the black marks the average over all 9 flies. See methods for discussion of error bars calculation. (b) Similar to the computation done for (a), we can compute the average amount of information that is gained about the identity of the fly, give its word distribution at a specific time, compared with the mixture of the word distribution of all of the 9 flies. Averaging over all times, we get the average amount of information gained about the identity of fly 1 based on its time dependent word distributions (red), fly 6 (blue), and the average over the 9 flies (black).
Figure 3. The information about the stimulus that a fly’s spike train carries is correlated with firing rate, and yet a significant part is in the temporal structure. (a) The rate at the H1 spike train provides information about the visual stimulus is shown as a function of the average spike rate, with each fly providing a single data point (Fly 1 is marked by a red point and Fly 6 by a blue one). The linear fit of the data points for the 9 flies corresponds to a universal rate of $`2.39\pm 0.24\mathrm{bits}/\mathrm{spike}`$, as noted in the text. (b) The extra amount of information that the temporal structure of the spike train of each of the flies carry about the stimulus, as a function of the average firing rate of the fly \[see Eq. (12) in Methods\]. The average amount of additional information that is carried by the temporal structure of the spike trains, over the population is $`45\pm 17\%`$.
Figure 4. What different flies mean by the same words. The word-triggered averages are shown for flies 1 (red) and 6 (blue) for 4 different words. Similar to the computation of spike triggered averages, we compute the average velocity profile of the movie presented to the flies, preceding 7-letter binary words. (a) The average stimulus waveform preceding the word $`{}_{}{}^{}0000000_{}^{}`$ for Flies 1 (red) and 6 (blue), is shown as a function of the time relative to the end of the word (shown in actual time order on the right top side of the panel). (b-d) Word triggered averages for 3 other words, reflecting that the waveforms haves similar rough structure, and that the difference between the flies is in the details.
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# Exact Eigenstates for Repulsive Bosons in Two Dimensions
## I Introduction
The recent observation of Bose-Einstein condensation in dilute gases of alkali metals has stimulated much interest in the properties of systems of interacting bosons. A question of particular interest is whether such systems will form vortices under rotation, as occurs in superfluid <sup>4</sup>He. Experimentally such vortex states have been observed, both in two component systems and in a stirred condensate of <sup>87</sup>Rb atoms . Theoretically the stability of such vortices has been considered both in the Thomas-Fermi limit of strong interactions between atoms using the Gross-Pitaevskii equation , and also in the weak interaction limit where the coherence length is much larger than the size of the atom cloud. It is the latter limit we shall focus upon in this paper.
Pursuing an analogy with the fractional quantum Hall effect , we previously introduced a model of weakly interacting bosons in a harmonic well to address the question of whether attractive bosons condense. We proved analytically that all the angular momentum in this model is carried by the center of mass for attractive bosons, whereas for repulsive bosons we numerically found that a vortex state forms in the case of one unit of angular momentum per boson. Further numerical work by one of the authors extensively investigated the properties of ground states of the repulsive model for angular momentum $`L>N`$, demonstrating that although such states are more complicated than the analytic ones known for $`LN`$, they can still be understood either within vortex or composite fermion or boson pictures. Mottelson considered the low-lying eigenstates for the case $`LN`$, whilst Bertsch and Papenbrock performed numerical computations and noted that the ground state of the repulsive model for $`LN`$ is the elementary symmetric polynomial, $`\stackrel{~}{e}_L`$, of coordinates, $`z_i=x_i+iy_i`$, relative to the center of mass, $`R=_iz_i/N`$. Finally, recent work by Kavoulakis et al and Jackson et al has considered the relationship between mean and exact numerical solutions in the limit of large $`N`$.
In this Letter we present an analytical proof that the state described above is an exact eigenstate of the model for $`LN`$. We have unfortunately been unable to show that this state is the ground state, although we have considerable numerical evidence to show that this is the case.
The model is of $`N`$ bosons in a harmonic potential in two dimensions interacting via a delta–function potential, for which the Hamiltonian is
$$H=\underset{i=1}{\overset{N}{}}\left[\frac{\mathrm{}^2}{2m}_i^2+\frac{1}{2}m\omega ^2r_i^2\right]+V\underset{i<j}{}\delta (𝐫_i𝐫_j)$$
(1)
The natural way to look at this problem is in second quantised form, so first we must solve the non-interacting problem. We do this in plane polar coordinates since we are interested in angular momentum properties. The one-particle wavefunctions and energies are then
$`\psi _{n_r,\mathrm{}}`$ $`=`$ $`N_{n_r,\mathrm{}}R(n_r,|\mathrm{}|,r^2/2)r^{\mathrm{}}e^{i\mathrm{}\theta }`$ (2)
$`E_{n_r,\mathrm{}}`$ $`=`$ $`(n_r+|\mathrm{}|+1)\mathrm{}\omega `$ (3)
where $`R(n,l,x)`$ is the confluent hypergeometric function, $`n_r`$ is the radial quantum number and $`\mathrm{}`$ is the angular momentum. The total energy and angular momentum of a system of $`N`$ non-interacting bosons in this harmonic well is thus
$$E_{tot}=\underset{i=1}{\overset{N}{}}(n_{r,i}+|\mathrm{}_i|+1)\mathrm{}\omega ,L_{tot}=\underset{i=1}{\overset{N}{}}\mathrm{}_i.$$
(4)
The ground state manifold is then obtained by putting all bosons into the lowest radial state, $`n_{r,i}=0`$, and choosing angular momentum $`\mathrm{}_i`$ all of the same sign (which we choose to be positive), such that
$$\underset{i=1}{\overset{N}{}}\mathrm{}_i=LE_{tot}=(L+N)\mathrm{}\omega .$$
(5)
It can be seen that the ground state manifold has a degeneracy which is equal to the number of $`\mathrm{𝑝𝑎𝑟𝑡𝑖𝑡𝑖𝑜𝑛𝑠}`$ of the total angular momentum $`L`$ into the $`N`$ angular momenta $`\mathrm{}_i`$ of the individual bosons.
If we assume that the dimensionless interaction strength is very small, $`|\eta |1`$, where $`\eta =V/\mathrm{}\omega `$, we can treat it as a perturbation whose sole effect will be to lift the large degeneracy of the non-interacting ground state. This means we can work within the non-interacting ground state manifold, which freezes out the one-particle part of the Hamiltonian, and consider the effect of the interaction using degenerate perturbation theory. This is essentially an extension of an approach used for fractional quantum Hall effect systems to bosons. The normalised one-particle wavefunctions are given by
$$\psi _{\mathrm{}}(r,\theta )=\frac{1}{\sqrt{\mathrm{}!2\pi }}r^{\mathrm{}}e^{i\mathrm{}\theta }e^{r^2/2}=\frac{1}{\sqrt{\mathrm{}!2\pi }}z^{\mathrm{}}e^{|z|^2/2},$$
(6)
where we have moved to complex notation. The second quantized form of the interaction Hamiltonian is thus
$$\widehat{H}=\underset{m_1,m_2,n_1,n_2}{}V_{m_1m_2n_1n_2}a_{m_1}^+a_{m_2}^+a_{n_1}a_{n_2}$$
(7)
where the matrix element
$$V_{m_1m_2n_1n_2}=\frac{\eta }{4\pi }\frac{(m_1+n_1)!}{2^{m_1+m_2}}\delta _{m_1+m_2,n_1+n_2}.$$
(8)
For future convenience we will set $`\eta =4\pi `$.
To perform the degenerate perturbation theory for given $`N`$ and $`L`$, we first need the basis states which are labelled by the partitions of $`L`$ into $`N`$ pieces. Let us write a partition $`\lambda `$ in the form
$$\lambda =0^{\lambda _0}1^{\lambda _1}2^{\lambda _2}\mathrm{}=\underset{i}{}i^{\lambda _i}$$
(9)
where $`\lambda _i=N`$ and $`i\lambda _i=L`$. The corresponding basis state is then
$$|\lambda =\left[\underset{i}{}\frac{(a_i^+)^{\lambda _i}}{\sqrt{\lambda _i!}}\right],|\mathrm{\hspace{0.17em}0}$$
(10)
where $`a_{\mathrm{}}^+`$ creates a boson of angular momentum $`\mathrm{}`$. In coordinate space, this basis state takes the form
$$\left[\frac{\underset{i}{}\lambda _i!}{(2\pi )^NN!_i(i!)^{\lambda _i}}\right]^{1/2}m_\lambda (z_1,z_2\mathrm{}z_N),$$
(11)
where $`m_\lambda `$ is the monomial symmetric polynomial corresponding to the partition $`\lambda `$. The latter is the symmetric polynomial of the $`N`$ variables $`(z_1,z_2\mathrm{}z_N)`$ which has $`\lambda _i`$ $`i`$-th powers.
To consider the problem either analytically or numerically requires us to calculate the elements of the symmetric interaction matrix, $`H_{\lambda \mu }=\lambda |\widehat{V}|\mu `$. This is obviously best performed using the second quantized approach. There are two types of matrix elements, diagonal and off-diagonal, and their evaluation is different. For the diagonal elements $`H_{\lambda \lambda }`$, we must sum over every possible pair of elements in the partition, whether distinct or identical,
$$H_{\lambda \lambda }=\underset{i}{}\lambda _i(\lambda _i1)V_{\lambda _i\lambda _i\lambda _i\lambda _i}+4\underset{i<j}{}\lambda _i\lambda _jV_{\lambda _i\lambda _j\lambda _i\lambda _j}.$$
(12)
The off-diagonal elements $`H_{\lambda \mu }`$ can only be non-zero if $`\lambda `$ and $`\mu `$ differ only in the angular momenta of two particles. This can affect either 4 separate angular momentum states, as in the case where the angular momentum transfer is $`0+41+3`$; or 3 separate angular momentum states, as in the case of angular momentum transfer $`0+42+2`$. In terms of partitions, in case (i) $`\lambda _i=\mu _i`$ except at four values of $`i`$, which we call $`i_1\mathrm{}i_4`$, and $`\lambda _{i_1}=\mu _{i_1}+1`$, $`\lambda _{i_2}=\mu _{i_2}+1`$, $`\lambda _{i_3}=\mu _{i_3}1`$, $`\lambda _{i_4}=\mu _{i_4}1`$. The matrix element is then
$$H_{\lambda \mu }=H_{\mu \lambda }=4\sqrt{\lambda _{i_1}\lambda _{i_2}\mu _{i_3}\mu _{i_4}}V_{i_1i_2i_3i_4}.$$
(13)
In case (ii) $`\lambda _i=\mu _i`$ except at three values of $`i`$ which we call $`i_1,i_2,i_3`$, and $`\lambda _{i_1}=\mu _{i_1}+2`$, $`\lambda _{i_2}=\mu _{i_2}1`$, $`\lambda _{i_3}=\mu _{i_3}1`$. The matrix element is then
$$H_{\lambda \mu }=H_{\mu \lambda }=2\sqrt{\lambda _{i_1}(\lambda _{i_1}1)\mu _{i_2}\mu _{i_3}}V_{i_1i_1i_2i_3}.$$
(14)
These formulas allow one to write down the interaction matrix $`H_{\lambda \mu }`$ which can then be diagonalised to give the energy eigenvalues and eigenstates.
## II The Subspace Property
If we look at the original Hamiltonian, we find that the interaction term depends only upon relative coordinates. To see this, change variables to the center of mass variable, $`R=_iz_i/N`$, and $`N1`$ relative coordinates such as $`\stackrel{~}{z}_i=z_iR`$, where $`i=1\mathrm{}N1`$: the interaction is then a function only of the relative coordinates $`\stackrel{~}{z}_i`$. It follows that if $`\psi (z_1\mathrm{}z_N)`$ is an eigenfunction of $`H`$ with a certain energy $`E`$ and angular momentum $`L`$, then $`R\psi (z_1\mathrm{}z_N)`$ is an eigenfunction of $`H`$ with energy $`E`$ and angular momentum $`L+1`$. This means that all eigenfunctions of $`H`$ at a given $`L`$ are reproduced at all higher $`L`$, and thus the number of new states at any $`L`$ is just $`n_Ln_{L1}`$, the difference between the number of partitions of $`L`$ and $`L1`$ respectively .
This subspace property makes it natural to think in terms of a second type of symmetric polynomial, the elementary symmetric polynomials. These are defined by, for $`N`$ variables,
$$e_L=\underset{i_1<i_2<\mathrm{}<i_L}{}z_{i_1}z_{i_2}\mathrm{}z_{i_L}$$
(15)
where $`LN`$. For a general partition $`\lambda `$ we define
$$e_\lambda =\underset{i}{}e_i^{\lambda _i}.$$
(16)
The set of new eigenstates at any total angular momentum $`L`$ is seen to be spanned by the $`\stackrel{~}{e}_\lambda `$, the elementary symmetric polynomials of the relative coordinates $`\stackrel{~}{z}_i`$, where we now include $`\stackrel{~}{z}_N=\stackrel{~}{z}_1\stackrel{~}{z}_2\mathrm{}\stackrel{~}{z}_{N1}`$. Since $`\stackrel{~}{e}_1=0`$, only partitions with $`\lambda _1=0`$ can be formed (i.e. partitions with no part equal to $`1`$), which gives exactly the correct number of states.
## III Proof that the $`\stackrel{~}{}\text{e}_L`$ are Exact Eigenstates
In this section we prove that the states $`\stackrel{~}{e}_L`$ for $`LN`$ ($`L1`$) are eigenstates of $`\widehat{H}`$ with eigenvalue $`N(N1L/2)`$. The proof is essentially a brute force method: we operate the Hamiltonian on the state $`\stackrel{~}{e}_L`$, and show that the result is the above eigenvalue times $`\stackrel{~}{e}_L`$. The derivation naturally falls into five steps, which we detail below. (1) Writing $`\stackrel{~}{}e_L`$ in Terms of $`𝐞_L`$ and $`𝐑`$ Consider $`\stackrel{~}{e}_M`$, which can be written as
$$\stackrel{~}{e}_M=\underset{i_1<i_2<\mathrm{}<i_M}{}(z_{i_1}R)\mathrm{}(z_{i_M}R).$$
(17)
If we expand out this product we will get the elementary symmetric polynomials of the $`z_i`$, namely the $`e_L`$, $`LM`$, multiplied by $`R^{ML}`$. To get the correct coefficients in this expansion we note that $`\stackrel{~}{e}_M`$ has $`N!/M!(NM)!`$ terms whilst $`e_L`$ has $`N!/L!(NL)!`$ terms. In the expansion of $`\stackrel{~}{e}_M`$, each product term will produce $`M!/L!(ML)!`$ terms of the type which will add up to produce $`R^{ML}e_L`$, so that the coefficient of $`e_L`$ in the expansion is $`(NL)!/(NM)!(ML)!`$ It follows that
$`\stackrel{~}{e}_M={\displaystyle \underset{L=2}{\overset{M}{}}}`$ $`(`$ $`1)^{ML}{\displaystyle \frac{NL!}{NM!ML!}}e_LR^{ML}`$ (18)
$`+`$ $`(`$ $`1)^{M1}{\displaystyle \frac{N!(M1)}{NM!M!}}R^M,`$ (19)
where we have noticed that, since $`e_1=R`$, the last two terms have the same form and should be combined. (2) Operating Hamiltonian $`\widehat{}H`$ on $`𝐞_L`$ An important feature of $`e_L`$ is that it is also the monomial function $`m_\lambda `$ corresponding to $`\lambda =0^{NL}1^L`$. The normalised version of $`e_L`$ can thus be written as $`|e_L|\mathrm{\hspace{0.17em}0}^{NL}1^L`$, and it is clear that the Hamiltonian $`\widehat{H}`$ can only connect this state to itself and $`|A|\mathrm{\hspace{0.17em}0}^{NL+1}1^{L2}2^1`$, where we have labelled this state as $`|A`$ for convenience in what follows. The two matrix elements can be calculated using from the formulas derived in the previous sections for diagonal and off-diagonal elements respectively. The diagonal element is given by
$$e_L|\widehat{H}|e_L==N^2N\frac{1}{2}L(L1).$$
(20)
The off-diagonal element is found by using the rule for the case where only three separate angular momentum states change (here $`1+10+2`$) to give
$$A|\widehat{H}|e_L=\frac{1}{2}\sqrt{2L(L1)(NL+1)}$$
(21)
The final result for the operation of the Hamiltonian on $`e_L`$ is thus
$`\widehat{H}|e_L`$ $`=`$ $`\left[N^2N{\displaystyle \frac{1}{2}}L(L1)\right]|e_L`$ (22)
$`+`$ $`{\displaystyle \frac{1}{2}}\sqrt{2L(L1)(NL+1)}|A.`$ (23)
(3) Removing the Normalisation Factors We want to get rid of the normalisation factors for the eigenstates so that the Hamiltonian will act directly on symmetric polynomials such as $`e_L`$. In the previous equation we should divide by the normalization factor for $`|\mathrm{\hspace{0.17em}0}^{NL}1^L`$ and multiply by that for $`|\mathrm{\hspace{0.17em}0}^{NL+1}1^{L2}2^1`$. Using the form for the normalisation factors given in Eq. (11), we find that the normalised version of Eq. (22) is
$$\widehat{H}|e_L)=[N^2N\frac{1}{2}L(L1)]|e_L)+\frac{NL+1}{2}|A),$$
(24)
where the $`|\lambda )`$ are the symmetric polynomials with no normalisation factor. (4) Relating $`|\mathrm{\hspace{0.17em}0}^{NL+1}1^{L2}2^1)`$ to $`𝐞_L`$ and $`𝐑e_{L1}`$ Consider the product
$$NR|e_{L1})=\left[\underset{i=1}{\overset{N}{}}z_i\right]_{i_1<\mathrm{}<i_{L1}}z_{i_1}z_{i_2}\mathrm{}z_{i_{L1}}$$
(25)
The above product can clearly only produce $`|0^{NL}1^L)`$ or $`|0^{NL+1}1^{L2}2^1)`$, depending upon whether the $`z_j`$ from the prefactor is not or is included in the set $`(z_{i_1},z_{i_2}\mathrm{}z_{i_{L1}})`$. Moreover we see that each element of $`|0^{NL+1}1^{L2}2^1)`$ can only be made in one way, so that
$$NRe_{L1}=|0^{NL+1}1^{L2}2^1)+Ce_L.$$
(26)
To find the coefficient $`C`$ we just count terms: $`NRe_{L1}`$ has $`(N+1)!/(NL+1)!(N1)!`$ terms, $`|A)`$ has $`N!/(NL+1)!(L2)!`$ terms, and $`e_L`$ has $`N!/(NL)!L!`$ terms. This leads to the result $`C=L`$, and hence
$$|A)=NRe_{L1}Le_L.$$
(27)
Combining this with Eq. (24) gives
$$\widehat{H}e_L=\left[N^2\left(1+\frac{L}{2}\right)N\right]e_L+\frac{N(NL+1)}{2}Re_{L1}.$$
(28)
(5) Operating $`\widehat{}H`$ onto $`\stackrel{~}{}e_M`$ If we now operate $`H`$ onto $`\stackrel{~}{e}_M`$ using Eq. (18) to write $`\stackrel{~}{e}_M`$ in terms of $`e_L`$ and $`R`$, and using Eq. (28) to operate $`\widehat{H}`$ onto $`e_L`$, we get after a lot of tedious algebra,
$$\widehat{H}\stackrel{~}{e}_M=\left[N^2\left(1+\frac{M}{2}\right)N\right]\stackrel{~}{e}_M.$$
(29)
Note that we have to compare the coefficients of three types of term separately: (i) $`e_M`$, (ii) $`R^{ML}e_L`$ for $`2L<M`$, and (iii) $`R^M`$. The proof can be simplified a little if we use the subspace property. Introduce a projection operator, $`\widehat{P}`$, which removes any term containing a factor of $`R`$. Operating $`\widehat{P}`$ onto Eq. (18) gives $`\widehat{P}e_L=\stackrel{~}{e}_L`$, whilst operating $`\widehat{P}`$ onto Eq. (28) gives Eq. (29).
## IV Discussion and Conclusions
We have considered $`N`$ bosons in a 2D harmonic potential interacting via repulsive delta-function potentials and with fixed total angular momentum $`LN`$. Within the “lowest Landau level” approximation, we have analytically shown that the elementary symmetric polynomial of coordinates relative to the center of mass, $`\stackrel{~}{e}_L`$ is an exact eigenstate of this Hamiltonian with eigenvalue $`N(NL/21)`$. Extensive numerical analysis shows that this state is actually the ground state. This is not surprising since in the special case $`L=N`$ this is just the one-vortex state discussed in Ref. which is expected to be the ground state by analogy to superfluid <sup>4</sup>He. One can also see that there is a sense in which $`\stackrel{~}{e}_L`$ distributes the angular momentum equally between particles subject to the subspace property. We have attempted to prove analytically that $`\stackrel{~}{e}_L`$ is the ground state, but have so far failed. The situation is much harder than in the proof we presented in Ref. to show that the eigenstate with the largest eigenvalue corresponds to all angular momentum being in the center of mass motion. The main problem is that the eigenvector for the smallest eigenvalue has components of both signs to reduce its energy, and frustration results between any set of three basis states that are connected by $`\widehat{H}`$: the difficulty is essentially that of a quantum antiferromagnet compared to a ferromagnet. For completeness we hope that the state considered in this paper will be analytically shown to be the ground state in the future, but from numerics there can be little doubt that it is.
ACKNOWLEDGEMENTS
We thank N.R. Cooper, J.M.F. Gunn and M.W. Long for useful discussions. We acknowledge financial support from the EPSRC Grants GR/M98975, GR/L28784 and the Nuffield Foundation.
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# Proposition 1.1
MODELLING OF STOCK PRICE CHANGES:
A REAL ANALYSIS APPROACH
Rimas Norvaiša<sup>1</sup><sup>1</sup>1\*”Research supported in part by an NSERC Canada Collaborative Grant at Carleton University, Ottawa, Canada, by a U.S. National Science Foundation Grant and by an NWO (Dutch Science Foundation) grant. Manuscript received: …; final version received: …
Institute of Mathematics and Informatics, Akademijos 4, LT-2600 Vilnius, Lithuania (e-mail: norvaisaktl.mii.lt)
Abstract. In this paper a real analysis approach to stock price modelling is considered. A stock price and its return are defined in a duality to each other provided there exist suitable limits along a sequence of nested partitions of a time interval, mimicking sum and product integrals. It extends the class of stochastic processes susceptible to theoretical analysis. Also, it is shown that extended classical calculus is applicable to market analysis whenever the local $`2`$–variation of sample functions of the return is zero, or is determined by jumps if the process is discontinuous. In particular, an extended Riemann-Stieltjes integral is used in that case to prove several properties of trading strategies.
Key words: continuous–time model, model testing, stock price, return, trading strategy
JEL Classification: G10, G12, C13
Mathematics Subject Classification (1991): 90A09, 60G17, 26A42
1 Introduction and discussion
In continuous–time financial mathematics the solution to the Doléans–Dade stochastic differential equation is often used as a model for stock price changes. The semimartingale driving this equation is called the return. Since many conclusions on the price behavior depend on the return, it plays an important role in mathematics of finance. On the other hand, the returns in econometrics of financial markets are sometimes modelled by stochastic processes which are not semimartingales. To provide a theoretical justification for such cases, one introduces a Doléans–Dade type equation with the stochastic integral replaced by a different integral. However, solutions to integral equations based on different integrals may differ considerably as demonstrated Wong and Zakai (1965). One may ask then whether it is possible to build up a model of stock price changes which is independent of a particular integration theory? The present paper addresses this question and provides a new insight into the relation between theoretical and applied financial mathematics.
1.1. Prices and returns. To begin with we discuss two continuous–time stochastic models for a frictionless stock market. Let $`R=\{R(t):\mathrm{\hspace{0.17em}0}tT\}`$ be a semimartingale such that $`R(0)=0`$ almost surely and let $`Q=\{Q(t):\mathrm{\hspace{0.17em}0}tT\}`$, where $`Q(t):=\mathrm{exp}\{R(t)\}`$ for $`0tT`$. The pair $`(Q,R)`$ will be called the exponential system of a stock. Then $`Q`$ is the price and $`R`$ is the return of a stock of the exponential system $`(Q,R)`$. Let $`P=\{P(t):\mathrm{\hspace{0.17em}0}tT\}`$ be a stochastic process satisfying the equation
$$P(t)=1+(SI)_0^tP(s)𝑑R(s),0tT,$$
$`1.1`$
where $`P(0):=1`$ and $`(SI)`$ denotes the stochastic integral defined by the $`L^2`$-isometry. Doléans–Dade (1970) proved that the unique solution to (1.1) is given by
$$P(t)=\mathrm{exp}\{R(t)\frac{1}{2}R^c,R^c(t)\}\underset{(0,t]}{}(1+\mathrm{\Delta }R)\mathrm{exp}\{\mathrm{\Delta }R\},0<tT,$$
and $`P(0)=1`$, where $`R^c`$ is the continuous local martingale part of the semimartingale $`R`$ and $`\mathrm{\Delta }R(s):=R(s)R(s)`$ for $`s(0,T]`$. If $`P`$ satisfies (1.1) and is bounded away from zero then, by associativity of the stochastic integral, we have
$$R(t)=(SI)_0^t\frac{dP(s)}{P(s)},0tT.$$
$`1.2`$
The pair $`(P,R)`$ satisfying (1.1) will be called the stochastic exponential system. Then $`P`$ is the price and $`R`$ is the return of a stock of the stochastic exponential system. Parts of continuous–time financial mathematics based on the exponential system and on the stochastic exponential system will be called respectively the exponential model and the stochastic exponential model. In general, the exponential system is different from the stochastic exponential system. Indeed, if $`R`$ is a standard Brownian motion $`B=\{B(t):t0\}`$, then the solution to (1.1) is the stochastic process $`P_B(t):=\mathrm{exp}\{B(t)t/2\}`$, $`0tT`$, often called the geometric Brownian motion. In both systems the prices are observable quantities meaning that they represent real data, while the returns are non–observable and depend on the models. In addition to being a semimartingale, $`R`$ may sometimes be assumed to satisfy certain probabilistic conditions about its distribution. An adequacy to real data of such assumptions on $`R`$ can be tested by using the price transformations: the log return $`R(t)=\mathrm{log}Q(t)`$, $`0tT`$, for the exponential model, and the return (1.2) for the stochastic exponential model. The log return is often used in econometric literature which means that certain hypotheses about the exponential model are tested. If one wishes to test the stochastic exponential model then the return (1.2) has to be used. However (1.2) is not defined for a single sample function, so that its statistical tractability is problematic. On the other hand, as pointed out Bühlmann, Delbaen, Embrechts and Shiryaev (1996), under probabilistic price analysis, the stochastic exponential model turns out to be more advantageous than the exponential model. Therefore, it is appealing to modify the stochastic exponential model in such a way that to make it more manageable for statistical analysis. Bühlmann et al. (1996) provide the analysis of the exponential model via its reduction to the stochastic exponential model using a suitable transformation in (1.1) instead of $`R`$.
1.2. Price changes as an evolution. In the present paper we define a price and its return in a duality without a priori integrability or probabilistic assumptions (Definition 2.9 below), and show that almost all sample functions of many typical stochastic processes including a Brownian motion satisfy the new definition (Propositions 2.10 and 2.11 below). An idea behind the definition is based on known results about a one–to–one correspondence between an evolution and its generator. A family of real numbers $`U=\{U(s,t):astb\}`$ is an evolution on $`[a,b]`$ if $`U(t,t)=1`$ and
$$U(t,r)U(r,s)=U(t,s)\text{for all}asrtb.$$
Let $`P`$ be a function representing a stock price over a period $`[0,T]`$ such that $`P(0)=1`$, and let $`U(t,s):=P(t)/P(s)`$ for $`0stT`$. Then $`U`$ so defined is a simple example of an evolution on $`[0,T]`$ defined by stock price changes. An evolution arise in describing the state of nonautonomous systems and they are generalizations of the concept of a one–parameter semigroup of bounded linear operators on a Banach space describing the state of autonomous linear systems. The classical Hille–Yosida theorem describes any strongly continuous, contractive semigroup in terms of its generator. In this way the Hille–Yosida theorem provides a one–to–one correspondence between semigroups and their generators. An important difficult question is when and in what sense will a given evolution $`U`$ have a generator? The answer depends on a behaviour of the function $`[a,b]tU(t,a)`$, in particular on its $`p`$-variation. If $`U`$ is defined by stock price changes, that is if $`U(,0)=P`$, then its generator is a return as defined in the present paper. Therefore the pair $`(P,R)`$ satisfying Definition 2.9 is called the (weak) evolutionary system.
1.3. The $`p`$-variation. For the approach advocated in the present paper, the notion of $`p`$-variation of a function plays a role comparable with a role of a martingale property in the stochastic exponential model. For a function $`f:[a,b]`$ and a real number $`0<p<\mathrm{}`$, the $`p`$-variation $`v_p(f)=v_p(f;[a,b])`$ is the least upper bound of sums $`s_p(f;\kappa ):=_{i=1}^n|f(x_i)f(x_{i1})|^p`$ over all partitions $`\kappa =\{x_i:i=0,\mathrm{},n\}`$ of $`[a,b]`$. We notice that the $`2`$-variation is not the same as the quadratic variation. For a standard Brownian motion $`B=\{B(t):t0\}`$ and any $`0<T<\mathrm{}`$, $`v_2(B;[0,T])=+\mathrm{}`$ almost surely, while $`v_p(B;[0,T])<\mathrm{}`$ for each $`p>2`$ and the quadratic variation of $`B`$ is defined in the almost sure sense for certain sequences of partitions. For any function $`f`$ on $`[a,b]`$, define the index of $`p`$-variation $`\upsilon (f)=\upsilon (f;[a,b])`$ by
$$\upsilon (f;[a,b]):=\{\begin{array}{cc}inf\{p>0:v_p(f)<\mathrm{}\}\hfill & \text{if the set is nonempty}\hfill \\ +\mathrm{}\hfill & \text{otherwise.}\hfill \end{array}$$
Therefore for a Brownian motion $`B`$, $`\upsilon (B;[0,T])=2`$ almost surely. Also for any $`0<T<\mathrm{}`$, $`\upsilon (X;[0,T])<2`$ almost surely if $`X`$ is a mean zero Gaussian stochastic process with stationary increments, continuous in quadratic mean and the incremental variance $`\{E[X(t+u)X(t)]^2\}^{1/2}`$ varies regularly as $`u0`$ with index $`\gamma >1/2`$, or if $`X`$ is a homogeneous Lévy process with the Lévy measure $`L`$ such that $`_{\{0\}}(1|x|^p)L(dx)<\mathrm{}`$ for some $`p<2`$.
1.4. Stochastic and classical calculi. In this paper it is proved that evolutionary systems $`(P,R)`$ possess a uniformity property with respect to sequences of partitions defining $`P`$ and $`R`$ provided the $`p`$-variation index $`\upsilon (R)<2`$. This fact is important when one deals with fitting a model to real data, or when one considers a relation between discrete–time and continuous–time models (cf. Theorem 3.5 below). With the help of the result of Föllmer (1981) one can show that a weak evolutionary system $`(P,R)`$ satisfies an integral equation similar to (1.1) with a different integral. If sample functions of the return $`R`$ in the evolutionary system $`(P,R)`$ have the $`p`$-variation index $`\upsilon (R)<2`$, then (1.1) and (1.2) hold path by path with the stochastic integral replaced by the Left Young integral, an extended Riemann–Stieltjes integral. Norvaiša (1999) proved that the values of the Left Young integral and the values of the corresponding stochastic integral agree almost surely under conditions ensuring the existence of both. In this sense the value $`\upsilon (R)=2`$ of the $`p`$-variation index is a borderline between an area where classical calculus applies and an area where stochastic calculus is needed essentially. We notice that the semimartingale property of the return $`R`$ in the stochastic exponential system $`(P,R)`$ makes a borderline between classical and stochastic calculi on a different level, e.g. the value $`\upsilon (R)=1`$ (the $`1`$-variation is the same as the total variation). Several examples of returns such as a hyperbolic Lévy motion, a normal inverse Gaussian Lévy process, the V.G. process, an $`\alpha `$-stable Lévy motion with $`\alpha [1,2)`$, or a fractional Brownian motion with the Hurst index $`H(1/2,1)`$, can be treated using classical calculus. The aim of the present paper is to find a connection between the two calculi for the mathematics of finance. However, more interesting is a question whether it is possible to develop a full fledged model of a financial market based on the evolutionary system. Clearly it is not possible to answer to this question at this writing. A model construction requires a more advanced development of theories of integral equations and optimal control for functions of bounded $`p`$-variation, as well as further development of concepts of market efficiency, equilibrium and risk in the new context.
1.5. Arbitrage. We finish with a discussion of arbitrage for the evolutionary system. In the continuous–time financial mathematics based on the semimartingale theory the first and second fundamental theorems deal with the key principals of the theory. These theorems relate suitable forms of an arbitrage with the existence of a (unique) martingale measure and completeness. Thus an applicability of these tools is restricted if arbitrage is possible. In particular, this concerns the contingent claim valuation theory based on the no arbitrage principle. Less formally, the no arbitrage principle is considered as a natural property of a model of an ideal financial market because “there is no such thing as a free lunch” in equilibria market. These arguments may give an impression that no approaches other than martingale based stochastic calculus can be useful for mathematical finance. However the situation is not as simple as it may look. There are examples from a game theory where such a thing as a free lunch is possible under equilibria (see p. 137 in Kac, Rota and Schwartz, 1992). On the other hand, non–equilibrium can explain stylized facts discovered through the statistical analysis of market data (see Chapter 4 of an overview of Focardi and Jonas, 1997, based on interviews with over 100 persons in industry and academia). A strong critique of a whole current financial mathematics comes from actuaries who use different principles to value contingent claims (see e.g. Clarkson, 1996, 1997). So instead of avoiding arbitrage it seems more fruitful to have a model which accommodates both, free lunch areas as well as areas without a free lunch, and leave the question of performance evaluation of such a model to econometrics.
Next we illustrate how a real analysis approach may shed new light on arbitrage. One way to define an arbitrage for evolutionary systems is to follow the pattern from the stochastic exponential model which requires first to define a self–financing strategy. As pointed out Harrison and Pliska (1981, Section 7), the restriction to predictable trading strategies as well as to gains defined using the stochastic integral needs a careful study. Clearly we cannot use these constructions in the present setting. Instead we define self–financing strategies pathwise following the logic of the present approach (Definition 3.3 below), and prove that the criteria suggested by Harrison and Pliska (1981) does apply to the new notions (Theorem 3.5 below). Then arbitrage can be defined either for a single function representing a price evolution, or for almost every sample function of a stochastic process using notation of Section 3 as follows: given a price $`P=(P_0,\mathrm{},P_\nu )`$ (of $`1+\nu `$ assets) during a time period $`[0,T]`$, a self–financing $`P`$–trading strategy $`\varphi =(\varphi _0,\mathrm{},\varphi _\nu )`$ is an arbitrage opportunity for $`P`$ at time $`T`$ if the portfolio value function $`V^{\varphi ,P}`$ is $`0`$ at $`0`$ and positive at $`T`$. Salopek (1998) proved that an arbitrage in this sense can be constructed whenever the return of an evolutionary system is continuous function of bounded $`p`$-variation for some $`1p<2`$. To give a short proof of the same statement we modify the ingenious construction of an arbitrage due to Shiryaev (1998, Example VII.2c.4). To this aim we replace a fractional Brownian motion with the Hurst index $`H(1/2,1)`$ in his construction, with a continuous function of bounded $`p`$-variation for some $`1p<2`$, and use the chain rule formula given by Theorem 2.1 below instead of Itô’s formula.
###### Proposition 1.1
Let $`f`$ be a continuous function of bounded $`p`$-variation on $`[0,T]`$ for some $`1p<2`$ such that $`f(0)=0`$ and $`f(T)0`$, and let $`r`$, $`\sigma `$ be real numbers. Then for the evolutionary system $`(P,R)`$ with $`R(t)=(rt,rt+\sigma f(t))`$, $`0tT`$, there exists an arbitrage opportunity for $`P`$ at time $`T`$.
###### Demonstration Proof
Let $`R_0(t):=rt`$ and $`R_1(t):=rt+\sigma f(t)`$ for $`0tT`$. By Proposition 2.6, $`P_0(t)=e^{rt}`$ and $`P_1(t)=e^{rt+\sigma f(t)}`$ for $`0tT`$. The vector function $`P=(P_0,P_1)`$ is the price in the sense defined in Section 3 below. Let $`\varphi _0(t):=1\mathrm{exp}\{2\sigma f(t)\}`$ and $`\varphi _1(t):=2[\mathrm{exp}\{\sigma f(t)\}1]`$ for $`0tT`$. By Proposition 3.2, $`\varphi =(\varphi _0,\varphi _1)`$ is the $`P`$–trading strategy. Next we show that $`\varphi `$ is self–financing $`P`$–trading strategy as defined in Definition 3.3. To this aim we apply the chain rule formula from Theorem 2.1 twice: first take $`h1`$, $`F(u_1,u_2)=u_1u_2^2`$, and second take $`h=\varphi _1`$, $`F(u_1,u_2)=u_1u_2`$. For each $`0<tT`$, we then have
$`V^\varphi (t)`$ $`=e^{rt}\left[e^{\sigma f(t)}1\right]^2={\displaystyle _0^t}\left[e^{\sigma f(t)}1\right]^2𝑑e^{rt}+{\displaystyle _0^t}2e^{rt}\left[e^{\sigma f(t)}1\right]𝑑e^{\sigma f(t)}`$
$`={\displaystyle _0^t}\left[1e^{2\sigma f(t)}\right]𝑑e^{rt}+{\displaystyle _0^t}2\left[e^{\sigma f(t)}1\right]𝑑e^{rt+\sigma f(t)}=V^\varphi (0)+G^\varphi (t),`$
where all integrals exist in the Riemann-Stieltjes sense by the Stieltjes integrability theorem of L.C. Young (1936, p. 264). Since $`V^\varphi (0)=0`$ and $`V^\varphi (T)>0`$, the self–financing $`P`$–trading strategy $`\varphi `$ is an arbitrage opportunity for $`P`$ at time $`T`$. ∎
The preceding fact shows the irrelevance of a “long memory” of a fractional Brownian motion with respect to an arbitrage. Sample function behavior of a stochastic process is responsible for arbitrage opportunities. A Weierstrass function is a non–probabilistic example of a function $`f`$ satisfying hypotheses of Proposition 1.1. Once we except the evolutionary system as a base for a model of a financial market then an arbitrage is a property of a local behavior of a sample function rather than a correlation property between indefinitely increasing time moments.
2 Duality between price and return
2.1. Returns in discretetime models. If time $`t`$ is discrete, say $`t=0,1,\mathrm{},T`$, there are at least two different notions of return. Let $`P=\{P(t):t=0,1,\mathrm{},T\}`$ be a price of a stock which pays no dividends. The simple net return $`R_1=\{R_1(t):t=0,1,\mathrm{},T\}`$ is defined by setting $`R_1(0):=0`$, and for each $`t=1,\mathrm{},T`$,
$$\widehat{R}_1(t):=R_1(t)R_1(t1):=\{\begin{array}{cc}[P(t)P(t1)]/P(t1),\hfill & \text{if }P(t1)>0,\hfill \\ 0,\hfill & \text{if }P(t1)=0.\hfill \end{array}$$
$`2.1`$
Notice that $`\widehat{R}_1(t)`$ depends on values $`P(t1)`$ and $`P(t)`$, so that $`\widehat{R}_1`$ is the function of a subinterval $`[t1,t]`$. Notation $`\widehat{R}_1`$ (as well as $`\widehat{R}_2`$ defined below) is natural to use in discrete–time models where time lags have fixed length. A work with continuous–time models requires to treat returns either as interval functions defined on all subintervals of $`[0,T]`$, or as point functions on $`[0,T]`$. In this paper we choose to use the form of a point function. Given $`P(0)>0`$, there is a one–to–one correspondence between a positive price $`P`$ and a simple net return $`R_1`$ having jumps bigger than minus one, as described by Pliska (1997, Section 3.2). Namely, in addition to (2.1), for each $`t=1,\mathrm{},T`$, we have
$$P(t)=P(0)+\underset{s=1}{\overset{t}{}}P(s1)\widehat{R}_1(s)\text{ and }P(t)=P(0)\underset{s=1}{\overset{t}{}}[1+\widehat{R}_1(s)].$$
$`2.2`$
This correspondence is used in security market models by specifying simple net returns rather than prices.
Another type of a return is the log return $`R_2=\{R_2(t):t=0,1,\mathrm{},T\}`$ defined by setting $`R_2(0):=0`$ and, for all $`t=1,\mathrm{},T`$,
$$\widehat{R}_2(t):=R_2(t)R_2(t1):=\{\begin{array}{cc}\mathrm{log}[P(t)/P(t1)],\hfill & \text{if }P(t1)>0,\hfill \\ 0,\hfill & \text{if }P(t1)=0.\hfill \end{array}$$
This return is often used in the econometrics literature on security markets. It is easy to see that $`R_2`$ satisfies the additivity property
$$\widehat{R}_2(t)+\widehat{R}_2(t1)+\mathrm{}+\widehat{R}_2(ts+1)=\mathrm{log}[P(t)/P(ts)]$$
$`2.3`$
for any $`s,t\{0,1,\mathrm{},T\}`$, $`s<t`$, because the right side is the log return for the time period between $`ts`$ and $`t`$. The additivity property of log returns is one reason of its popularity among econometricians. For a discussion of these and other related properties of returns, see Campbell, Lo and MacKinlay (1997, Section 1.4.1).
Sometimes statistical conclusions based on the log return $`R_2`$ are applied to the model (2.2) or even to the continuous–time stochastic exponential model. To justify this one usually argues that $`\widehat{R}_1`$ and $`\widehat{R}_2`$ are relatively close to each other when both are small. However, there are cases when the difference between $`\widehat{R}_1`$ and $`\widehat{R}_2`$ cannot be neglected (see e.g. Elton, Gruber and Kleindorfer, 1975). To test the stochastic exponential model one needs to use the return $`R`$ defined via the Itô integral (1.2). Then one has to be able to evaluate the Itô integral using finitely many values of a single sample function. Given a sequence of partitions into shrinking subintervals of $`[0,T]`$, by the dominated convergence in probability theorem, the value of $`R`$ can be approximated in probability by corresponding Riemann–Stieltjes sums. However, it is not possible to conclude the convergence with probability 1 without further restrictions. A relationship between the price and its return is suggested below for the continuous–time framework which makes this approximation possible path by path. This relationship is motivated by a duality relation between additive and multiplicative interval functions, which is also known as the evolution representation problem. Recall that the additivity property is satisfied by the log return $`R_2`$ (see (2.3)) while the simple net return $`R_1`$ lacks this property. On the other hand, the multiplicativity property
$$\widehat{P}(t)\widehat{P}(t1)\mathrm{}\widehat{P}(ts+1)=P(t)/P(ts)$$
for any $`s,t\{0,1,\mathrm{},T\}`$, $`s<t`$, is satisfied by the price ratios $`\widehat{P}(t):=P(t)/P(t1)`$, $`t=1,\mathrm{},T`$.
2.2. The chain rule formula. For a finite interval $`J`$, open or closed at either end, let $`Q(J)`$ be the set of all partitions $`\kappa =\{x_i:i=0,\mathrm{},n\}`$ of $`J`$. As before for $`f:J`$ and $`0<p<\mathrm{}`$, let $`v_p(f):=v_p(f;J):=sup\{s_p(f;\kappa ):\kappa Q(J)\}`$ be the $`p`$-variation of $`f`$, where $`s_p(f;\kappa ):=_{j=1}^n|f(x_i)f(x_{i1})|^p`$ for $`\kappa =\{x_i:i=0,\mathrm{},n\}`$. Denote by $`𝒲_p=𝒲_p(J)`$ the set of all functions $`f`$ such that $`v_p(f)<\mathrm{}`$. If $`f𝒲_p`$ for some $`p<\mathrm{}`$ then $`f`$ is regulated, that is there exist the limits $`f(x):=lim_{yx}f(y)`$ and $`f(x+):=lim_{yx}f(y)`$ when these are defined. The class of all regulated functions on $`J`$ will be denoted by $`(J)`$. Given a regulated function $`f`$ on $`[a,b]`$, define a left-continuous function $`f_{}^{(a)}`$ and a right-continuous function $`f_+^{(b)}`$ by
$$\{\begin{array}{cc}\hfill f_{}^{(a)}(x)& :=f_{}(x):=f(x)\text{for}a<xb\text{and}f_{}^{(a)}(a):=f(a)\hfill \\ \hfill f_+^{(b)}(x)& :=f_+(x):=f(x+)\text{for}ax<b\text{and}f_+^{(b)}(b):=f(b).\hfill \end{array}$$
Given a regulated function $`f`$ on $`J`$, define $`\mathrm{\Delta }^{}f`$ on $`J`$ by $`\mathrm{\Delta }^{}f(x):=f(x)f(x)`$ if $`f(x)`$ is defined and $`\mathrm{\Delta }^{}f(x):=0`$ otherwise. Similarly define $`\mathrm{\Delta }^+f`$ on $`J`$ by $`\mathrm{\Delta }^+f(x):=f(x+)f(x)`$ if $`f(x+)`$ is defined and $`\mathrm{\Delta }^+f(x):=0`$ otherwise. Since each regulated function $`f`$ has at most countably many jumps, one can define $`𝔖_p(f):=𝔖_p(f;J):=\{_J(|\mathrm{\Delta }^{}f|^p+|\mathrm{\Delta }^+f|^p)\}^{1/p}`$. The local $`p`$-variation $`v_p(f)^{}:=v_p(f;J)^{}`$ is defined by
$$v_p(f;J)^{}:=\underset{\lambda Q(J)}{inf}sup\{s_p(f;\kappa ):\lambda \kappa Q(J)\}.$$
Then we have the relation $`𝔖_p(f)^pv_p^{}(f)v_p(f)`$. Let $`𝒲_p^{}=𝒲_p^{}(J):=\{f𝒲_p:𝔖_p(f)^p=v_p^{}(f)\}`$ for $`1<p<\mathrm{}`$. For regulated functions $`h`$ and $`f`$ on $`[a,b]`$, define the Left Young integral, or the $`LY`$ integral, by
$$(LY)_a^bh𝑑f:=(RS)_a^bh_{}^{(a)}𝑑f_+^{(b)}+\left[h\mathrm{\Delta }^+f\right](a)+\underset{(a,b)}{}\mathrm{\Delta }^{}h\mathrm{\Delta }^+f$$
provided the Riemann–Stieltjes integral exists in the refinement sense and the sum converges absolutely. Additivity on adjacent intervals as well as some other properties of the $`LY`$ integral are proved in Norvaiša (1999). From the Stieltjes integrability theorem of L.C. Young (1936) it follows that $`(LY)_a^bh𝑑f`$ is defined if $`h𝒲_p`$, $`f𝒲_q`$ and $`1/p+1/q>1`$. The following theorem of Norvaiša (1999) extends this result to the case when $`1/p+1/q=1`$ under additional assumptions on $`h`$ and $`f`$. Let $`\nu `$ be a positive integer, and let $`F`$ be a real-valued function defined on an open set $`U^\nu `$ containing a $`\nu `$-dimensional cube $`[c,d]^\nu :=[c,d]\times \mathrm{}\times [c,d]`$. We write $`F\mathrm{\Lambda }_{1,\alpha }([c,d]^\nu )`$ for $`\alpha (0,1]`$ if $`F`$ is differentiable on $`U`$ with partial derivatives $`F_l^{}`$, $`l=1,\mathrm{},\nu `$, and there is a finite constant $`K_\alpha `$ such that the inequality
$$\underset{1l\nu }{\mathrm{max}}|F_l^{}(u)F_l^{}(v)|K_\alpha \underset{k=1}{\overset{\nu }{}}|u_kv_k|^\alpha $$
holds for all $`u=(u_1,\mathrm{},u_\nu ),v=(v_1,\mathrm{},v_\nu )[c,d]^\nu `$.
###### Theorem 2.1
For $`\alpha (0,1]`$, let $`f=(f_1,\mathrm{},f_\nu ):[a,b](c,d)^\nu `$ be a vector function with coordinate functions $`f_l𝒲_{1+\alpha }^{}([a,b])`$ for $`l=1,\mathrm{},\nu `$, let $`F\mathrm{\Lambda }_{1,\alpha }([c,d]^\nu )`$ and let $`h`$ be a regulated function on $`[a,b]`$. Then the equality
$$(LY)_a^bhd(Ff)=\underset{l=1}{\overset{\nu }{}}(LY)_a^bh(F_l^{}f)𝑑f_l$$
$$+\underset{(a,b]}{}h_{}\left[\mathrm{\Delta }^{}(Ff)\underset{l=1}{\overset{\nu }{}}(F_l^{}f)_{}\mathrm{\Delta }^{}f_l\right]+\underset{[a,b)}{}h\left[\mathrm{\Delta }^+(Ff)\underset{l=1}{\overset{\nu }{}}(F_l^{}f)\mathrm{\Delta }^+f_l\right]$$
holds meaning that all $`\nu +1`$ integrals exist provided any $`d`$ integrals exist, and the two sums converge absolutely.
We refer to the preceding statement as the chain rule formula. Its proof is given by Norvaiša (1999).
2.3. Duality relation. Turning to a continuous–time model, consider an interval $`[0,T]`$, $`0<T<\mathrm{}`$. Roughly speaking, to extend (2.1) and (2.2) to functions defined on $`[0,T]`$, we pass to a limit along a nested sequence $`\lambda `$ of partitions of $`[0,T]`$. The first limit $`_\lambda (f)`$ if exists is an extension of the (sum) integral, and for $`f𝒲_2^{}`$, its values coincide with values of the $`LY`$ integral. The second limit $`_\lambda (g)`$ if exists is an extension of the product integral, and for sample functions $`g`$ of a semimartingale, its values coincide with values of the solution to the Doléans-Dade equation (1.1). First we prove the existence of $`_\lambda (f)`$ and $`_\lambda (g)`$ for functions $`f`$ and $`g`$ from the class $`𝒲_2^{}`$. Then the duality relations (2.13) are derived for such functions. Finally, a duality relation is proved for functions having defined the quadratic variation.
###### Definition Definition 2.2
Let $`𝔔([0,T])`$ be the set of all nested sequences $`\lambda =\{\lambda (m):m1\}`$ of partitions $`\lambda (m)=\{0=t_0^m<\mathrm{}<t_{n(m)}^m=T\}`$ of $`[0,T]`$ such that $`_m\lambda (m)`$ is dense in $`[0,T]`$. Let $`I_T`$ be either $`[0,T]`$ or $`[0,T)`$, and let $`f`$ be a real–valued function on $`I_T`$. Given $`\lambda 𝔔([0,T])`$, we say that $`_\lambda =_\lambda (f)`$ is defined on $`I_T`$ if the limit
$$_\lambda (f)(t):=\underset{m\mathrm{}}{lim}\underset{i=1}{\overset{n(m)}{}}[f(t_i^mt)f(t_{i1}^mt)]/f(t_{i1}^mt)$$
$`2.4`$
exists for each $`tI_T`$. Given a nonempty subset $`𝔔𝔔([0,T])`$, if $`_\lambda (f)`$ is defined for and does not depend on each $`\lambda 𝔔`$, then set $`_𝔔=_𝔔(f)`$ to be equal to any $`_\lambda (f)`$, $`\lambda 𝔔`$.
For a regulated function $`f`$, a typical example of $`𝔔𝔔([0,T])`$ in the preceding definition is the set $`𝔔(f)`$ defined by
$$𝔔(f):=\{\begin{array}{cc}𝔔([0,T]),\hfill & \text{if }f𝒟(I_T)\text{,}\hfill \\ \{\lambda 𝔔([0,T]):_m\lambda (m)\mathrm{\Delta }_f(I_T)\},\hfill & \text{if }f(I_T)𝒟(I_T)\text{,}\hfill \end{array}$$
$`2.5`$
where $`f𝒟(I_T)`$ if, at each point of $`(0,T)`$, $`f`$ is either right-continuous or left-continuous and $`\mathrm{\Delta }_f(I_T):=\{x(0,T):\mathrm{\Delta }^{}f(x)0`$ or $`\mathrm{\Delta }^+f(x)0\}`$. Next, under stated conditions we show that $`_{𝔔(f)}`$ is defined and has values of the indefinite $`LY`$ integral.
###### Proposition 2.3
Let $`f𝒲_2^{}(I_T)`$ and let $`inf\{f(t):tI_T\}\delta `$ for some $`\delta >0`$. Then $`_{𝔔(f)}(f)`$ is defined on $`I_T`$. Moreover, for each $`tI_T`$, $`f^1`$ is $`LY`$ integrable with respect to $`f`$ on $`[0,t]`$ and the relation
$$_{𝔔(f)}(f)(t)=(LY)_0^t\frac{df}{f}=\mathrm{log}\frac{f(t)}{f(0)}\underset{(0,t]}{}\left[\mathrm{log}\frac{f}{f_{}}\frac{\mathrm{\Delta }^{}f}{f_{}}\right]\underset{[0,t)}{}\left[\mathrm{log}\frac{f_+}{f}\frac{\mathrm{\Delta }^+f}{f}\right]$$
$`2.6`$
holds, where the two sums converge absolutely.
For the proof we need an auxiliary statement, where $`Q(S):=\{\kappa Q(J):\kappa S\}`$ for any subset $`SJ`$.
###### Lemma 2.4
Let $`f𝒲_p^{}([a,b])`$ for some $`1<p<\mathrm{}`$, and let $`S`$ be a dense subset of $`[a,b]`$ containing all discontinuity points of $`f`$. For each $`ϵ>0`$, there exists $`\lambda Q(S)`$ such that $`_{j=1}^kv_p(f;(z_{j1},z_j))<ϵ`$ whenever $`\lambda \{z_j:j=0,\mathrm{},k\}Q([a,b])`$.
###### Demonstration Proof
Let $`S[a,b]`$ be as in the statement. Then $`v_p^{}(f;[a,b])`$ is equal to the greatest lower bound of sums $`_{j=1}^kv_p(f;[z_{j1},z_j])`$ over $`\{z_j:j=0,\mathrm{},k\}Q(S)`$. Let $`ϵ>0`$. Since $`𝔖_p(f)<\mathrm{}`$, there exists a finite set $`\mu [a,b]`$ such that $`_\nu (|\mathrm{\Delta }^{}f|^p+|\mathrm{\Delta }^+f|^p)>𝔖_p(f)^pϵ/2`$ for each $`\nu \mu `$. Then one can choose $`\lambda Q(S)`$ such that $`\lambda \mu `$ and $`_{j=1}^kv_p(f;[z_{j1},z_j])<v_p^{}(f)+ϵ/2`$ whenever $`\lambda \{z_j:j=0,\mathrm{},k\}Q([a,b])`$. For each small enough $`\delta >0`$ and each $`j=1,\mathrm{},k`$, we have $`v_p(f;[z_{j1},z_j])v_p(f;[z_{j1},z_{j1}+\delta ])+v_p(f;[z_{j1}+\delta ,z_j\delta ])+v_p(f;[z_j\delta ,z_j])`$. Letting $`\delta 0`$ and using Lemma 2.19 of Dudley and Norvaiša (1999), we get
$`{\displaystyle \underset{j=1}{\overset{k}{}}}v_p(f;(z_{j1},z_j))`$ $`{\displaystyle \underset{j=1}{\overset{k}{}}}v_p(f;[z_{j1},z_j]){\displaystyle \underset{j=1}{\overset{k}{}}}\left[|\mathrm{\Delta }^{}f(z_j)|^p+|\mathrm{\Delta }^+f(z_{j1})|^p\right]`$
$`<V_p^{}(f)^p+ϵ/2𝔖_p(f)^p+ϵ/2=ϵ.`$
The proof of Lemma 2.4 is complete. ∎
###### Demonstration Proof of Proposition 2.3
The existence of the integral $`(LY)_0^tf^1𝑑f`$ and the second equality in (2.6) for each $`tI_T`$ follow from Theorem 2.1. To see if it’s true take $`F(u):=\mathrm{log}u`$ for $`u[\delta ,f_{\mathrm{}}]`$, $`\nu =1`$, $`\alpha =1`$, $`h1`$, and notice that $`\mathrm{\Delta }^{}(Ff)=\mathrm{log}(f/f_{})`$, $`\mathrm{\Delta }^+(Ff)=\mathrm{log}(f_+/f)`$. To prove that $`_{𝔔(f)}(f)`$ is defined on $`I_T`$ and that the first equality in (2.6) holds, for each $`u(0,T]I_T`$ and $`v[0,T)`$, let
$$\varphi ^{}(u):=\mathrm{log}\frac{f(u)}{f(u)}\frac{\mathrm{\Delta }^{}f(u)}{f(u)}\text{and}\varphi ^+(v):=\mathrm{log}\frac{f(v+)}{f(v)}\frac{\mathrm{\Delta }^+f(v)}{f(v)}.$$
To begin with the second case in (2.5) consider $`\{\lambda (m):m1\}𝔔([0,T])`$ such that $`_m\lambda (m)`$ contains all discontinuity points of $`f`$. For each $`m1`$, let
$$\psi _m(i):=\mathrm{log}\frac{f(t_i^m)}{f(t_{i1}^m)}\left[\frac{f(t_i^m)}{f(t_{i1}^m)}1\right]\text{for }i=1,\mathrm{},n(m)\text{.}$$
Let $`S:=_m\lambda (m)`$ and $`ϵ>0`$. By Lemma 2.4, and because $`f`$ is regulated and the two sums in (2.6) converge absolutely, one can choose $`\kappa :=\{z_j:j=0,\mathrm{},k\}Q(S)`$ such that
$$\underset{j=1}{\overset{k}{}}v_2(f;(z_{j1},z_j))<ϵ,\underset{1jk}{\mathrm{max}}Osc(f;(z_{j1},z_j))<\frac{\delta }{2}\text{and}\underset{\nu }{}\left(|\varphi ^{}|+|\varphi ^+|\right)<ϵ$$
for any $`\nu I_T\kappa `$. For each partition $`\lambda (m)=\{t_i^m:i=0,\mathrm{},n(m)\}`$ containing $`\kappa `$ and for each $`j\{0,\mathrm{},k\}`$, let $`i(j)\{0,\mathrm{},n(m)\}`$ be an index such that $`z_j=t_{i(j)}^m`$. Since
$$\underset{m\mathrm{}}{lim}\psi _m(i(j))=\varphi ^{}(z_j)\text{and}\underset{m\mathrm{}}{lim}\psi _m(i(j1)+1)=\varphi ^+(z_{j1})$$
$`2.7`$
for $`j=1,\mathrm{},k`$, one can choose an integer $`M1`$ such that $`\lambda (M)\kappa `$, there are at least two elements of $`\lambda (M)`$ in each interval $`(t_{i(j1)}^m,t_{i(j)}^m)`$, $`j=1,\mathrm{},k`$, and
$$\left|\underset{j=1}{\overset{k}{}}\left[\psi _m(i(j))\varphi ^{}(z_j)\right]+\underset{j=0}{\overset{k1}{}}\left[\psi _m(i(j)+1)\varphi ^+(z_j)\right]\right|<ϵ$$
for $`mM`$. Let $`tI_T`$. First suppose $`tS`$, so that $`t=t_{l(m)}^m`$ for some $`l(m)\{1,\mathrm{},n(m)\}`$ and for all $`m`$ larger than some $`N(t)`$. For each $`mMN(t)`$, let $`l:=\mathrm{max}\{jk:z_jt\}`$ and $`J:=\{0,\mathrm{},l(m)\}\{i(j),i(j1)+1:j=1,\mathrm{},l\}`$. By the Taylor series expansion with remainder, we have $`|\mathrm{log}(1+u)u|2u^2`$ for each $`|u|1/2`$. Then, for all $`mMN(t)`$, we get
$$\left|\mathrm{log}\frac{f(t)}{f(0)}\underset{(0,t]}{}\varphi ^{}\underset{[0,t)}{}\varphi ^+\underset{i=1}{\overset{n(m)}{}}\left[\frac{f(t_i^mt)}{f(t_{i1}^mt)}1\right]\right|$$
$$<2ϵ+2\underset{iJ}{}\left|\frac{f(t_i^m)}{f(t_{i1}^m)}1\right|^2<2ϵ+\frac{2}{\delta ^2}\underset{j=1}{\overset{k}{}}v_2(f;(z_{j1},z_j))<2ϵ(1+\delta ^2).$$
$`2.8`$
If $`tI_TS`$ then we have in addition the term
$$\left|\mathrm{log}\frac{f(t)}{f(t_{l(m)}^m)}\left[\frac{f(t)}{f(t_{l(m)}^m)}1\right]\underset{(t_{l(m)}^m,t]}{}\varphi ^{}\underset{[t_{l(m)}^m,t)}{}\varphi ^+\right|,$$
where $`l(m):=\mathrm{max}\{in(m):t_i^m<t\}`$. This term tends to zero as $`m\mathrm{}`$ because $`f`$ is continuous at $`t`$ in this case. Since $`ϵ`$ in (2.8) is arbitrary, $`_{𝔔(f)}(f)`$ is defined on $`I_T`$ and the first equality in (2.6) holds for the second case in (2.5). The proof when $`f𝒟(I_T)`$ is the same except that we use Lemma 2.4 with $`S=I_T`$, and choose $`\{t_i(j)^m:j=0,\mathrm{},k\}`$ so that $`z_j(t_{i(j)1}^m,t_{i(j)}^m]`$ if $`f`$ is right–continuous at $`z_j`$ and $`z_j[t_{i(j)}^m,t_{i(j)+1}^m)`$ if $`f`$ is left–continuous at $`z_j`$. The proof of Proposition 2.3 is complete.∎
Using notation as in Definition 2.2, we have:
###### Definition Definition 2.5
Let $`g`$ be a real–valued function on $`I_T`$. Given $`\lambda 𝔔([0,T])`$, we say that $`_\lambda =_\lambda (g)`$ is defined on $`I_T`$ if the limit
$$_\lambda (g)(t):=\underset{m\mathrm{}}{lim}\underset{i=1}{\overset{n(m)}{}}[1+g(t_i^mt)g(t_{i1}^mt)]$$
$`2.9`$
exists for each $`tI_T`$. Given a nonempty subset $`𝔔𝔔([0,T])`$, if $`_\lambda (g)`$ is defined for and does not depend on each $`\lambda 𝔔`$, then we define $`_𝔔=_𝔔(g)`$ to be equal to any $`_\lambda (g)`$, $`\lambda 𝔔`$.
Next, under stated conditions we show that (2.9) is defined and has values of the product integral with respect to $`g`$ over $`[0,t]`$, $`\text{}_0^t(1+dg)`$, defined as the limit of the product from $`i=1`$ to $`n`$ of $`1+g(t_i)g(t_{i1})`$, if it exists, under refinements of partitions $`\{t_i:i=0,\mathrm{},n\}`$ of $`[0,t]`$. The set $`𝔔(g)`$ in the following statement is defined by (2.5).
###### Proposition 2.6
Let $`g𝒲_2^{}(I_T)`$ and let $`(\mathrm{\Delta }^{}g)(\mathrm{\Delta }^+g)>1`$ on $`I_T`$. Then $`_{𝔔(g)}(g)`$ is defined on $`I_T`$. Moreover, for each $`tI_T`$, the product integral $`\text{}_0^t(1+dg)`$ exists, is positive and the relation
$$_{𝔔(g)}(g)(t)=\underset{0}{\overset{t}{\text{}}}(1+dg)=e^{g(t)g(0)}\underset{[0,t]}{}\left[(1+\mathrm{\Delta }^{}g)(1+\mathrm{\Delta }^+g)\right]e^{\mathrm{\Delta }^{}g\mathrm{\Delta }^+g}$$
$`2.10`$
holds, where the product converges absolutely.
###### Demonstration Proof
The product integral $`\text{}_0^t(1+dg)`$ exists, and the second equality in (2.10) holds for each $`tI_T`$ by Theorem 4.4 of Dudley and Norvaiša (1999). Since all jumps of $`g`$ are bigger than $`1`$, the positivity of the product integral follows from its definition. To prove that $`_{𝔔(g)}(g)`$ is defined on $`I_T`$ and the first equality in (2.10) holds, let $`tI_T`$ and $`\{\lambda (m):m1\}𝔔(g)`$, where $`\lambda (m)=\{t_i^m:i=0,\mathrm{},n(m)\}`$. One can assume that $`t=t_{l(m)}^m`$ for some $`1l(m)n(m)`$. Otherwise we include $`t`$ into $`\lambda (m)`$ and change indices. For a finite set $`\mu I_T`$, let
$$A(g;\mu ):=\underset{z\mu }{}[(1+\mathrm{\Delta }^{}g(z))(1+\mathrm{\Delta }^+g(z)]e^{\mathrm{\Delta }^{}g(z)\mathrm{\Delta }^+g(z)}.$$
To begin with the second case in (2.5) consider $`\{\lambda (m):m1\}𝔔([0,T])`$ such that $`_m\lambda (m)`$ contains all discontinuity points of $`g`$. Let $`S:=_m\overline{\lambda }(m)`$, where $`\overline{\lambda }(m)=\{0=t_0^m<\mathrm{}<t_{l(m)}^m=t\}`$, and $`ϵ(0,2A)`$, where $`A(g)`$ denotes the product $`_{[0,t]}`$ in (2.10). Because $`g`$ is regulated and by Lemma 2.4, one can choose $`\kappa =\{z_j:j=0,\mathrm{},k\}Q(S)`$ such that $`Osc(g(z_{j1},z_j))<1/2`$ for $`j=1,\mathrm{},k`$,
$$\underset{j=1}{\overset{k}{}}v_2(g;(z_{j1},z_j))<\frac{ϵ}{8eA}\text{and}\left|A(g;\mu )A(g)\right|<\frac{ϵ}{4}$$
for all $`\mu \kappa `$. For each $`\overline{\lambda }(m)\kappa `$ and for each $`j\{0,\mathrm{},k\}`$, let $`i(j)\{0,\mathrm{},l(m)\}`$ be such that $`z_j=t_{i(j)}^m`$. Let $`\mathrm{\Delta }_i^mg:=g(t_i^m)g(t_{i1}^m)`$ for $`i=1,\mathrm{},l(m)`$, $`m1`$, and let
$$U(g;\kappa ):=\underset{j=1}{\overset{k}{}}(1+\mathrm{\Delta }_{i(j)}^mg)(1+\mathrm{\Delta }_{i(j1)+1}^mg)\mathrm{exp}\left\{\mathrm{\Delta }_{i(j)}^mg\mathrm{\Delta }_{i(j1)+1}^mg\right\}.$$
Then letting $`J:=\{0,\mathrm{},l(m)\}\{i(j),i(j1)+1:j=1,\mathrm{},k\}`$, we have
$$\underset{i=1}{\overset{l(m)}{}}(1+\mathrm{\Delta }_i^mg)=e^{g(t)g(0)}U(g;\kappa )\mathrm{exp}\left\{\underset{iJ}{}\left[\mathrm{log}(1+\mathrm{\Delta }_i^mg)\mathrm{\Delta }_i^mg\right]\right\},$$
$`2.11`$
for all $`\overline{\lambda }(m)\kappa `$. Since $`|\mathrm{\Delta }_i^mg|1/2`$ for $`iJ`$, by the Taylor series expansion with remainder, we get
$$\left|\underset{iJ}{}\left[\mathrm{log}(1+\mathrm{\Delta }_i^mg)\mathrm{\Delta }_i^mg\right]\right|=\underset{iJ}{}\theta (\mathrm{\Delta }_i^mg)(\mathrm{\Delta }_i^mg)^22\underset{j=1}{\overset{k}{}}v_2(g;(z_{j1},z_j))<\frac{ϵ}{4eA},$$
where $`\theta (u)[2/9,2]`$ for $`|u|1/2`$. Let $`M1`$ be an integer such that $`\overline{\lambda }(M)\kappa `$, there are at least two elements of $`\overline{\lambda }(M)`$ in each interval $`(t_{i(j1)}^m,t_{i(j)}^m)`$, $`j=1,\mathrm{},k`$, and $`|1U(g;\kappa )/A(g;\kappa )|<ϵ/(8eA)`$ for all $`mM`$. Using the inequality $`|ue^v1|e|v|+2e|u1|`$ for $`|v|1/2`$ and $`|1u|1/4`$ one can show that (2.11) differs from the right side of (2.10) by $`ϵ\mathrm{exp}\{g(t)g(0)\}`$ for all $`mM`$. Since $`ϵ`$ is arbitrary, $`_{𝔔(g)}(g)`$ is defined on $`I_T`$ and the first relation in (2.10) holds when $`𝔔(g)`$ is defined by the second case in (2.5). The proof for the first case of (2.5) is similar and therefore is omitted. The proof of Proposition 2.6 is complete. ∎
To show a duality between $`_𝔔(g)`$ and $`_𝔔(f)`$ for $`g,f𝒲_2^{}`$ first we prove it between the indefinite product integral $`𝒫(g)`$ and the indefinite $`LY`$ integral $`𝒮(f)`$ defined by
$$𝒫(g)(t):=\underset{0}{\overset{t}{\text{}}}(1+dg)\text{and}𝒮(f)(t):=(LY)_0^t\frac{df}{f}$$
$`2.12`$
for $`tI_T`$ whenever the integrals exist. The following theorem was proved by Dudley and Norvaiša (1999, Theorem 6.8 and 6.10) for functions $`f,g`$ with values in a Banach algebra under the stronger assumption: $`f,g𝒲_p`$ for some $`p(0,2)`$. To extend this result to real–valued functions from the class $`𝒲_2^{}`$ we use the chain rule formula.
###### Theorem 2.7
I. Let $`g𝒲_2^{}(I_T)`$ and $`(\mathrm{\Delta }^{}g)(\mathrm{\Delta }^+g)>1`$ on $`I_T`$. Then the indefinite product integral $`𝒫(g)`$ is defined and $`𝒫(g)𝒲_2^{}(I_T)`$. Moreover, the indefinite $`LY`$ integral in (2.12) is defined for $`f=𝒫(g)`$ and $`𝒮(𝒫(g))(t)=g(t)g(0)`$ for $`tI_T`$.
II. Let $`f𝒲_2^{}(I_T)`$ and $`inf\{f(t):tI_T\}\delta `$ for some $`\delta >0`$. Then the indefinite $`LY`$ integral $`𝒮(f)`$ is defined and $`𝒮(f)𝒲_2^{}(I_T)`$. Moreover, the product integral in (2.12) is defined for $`g=𝒮(f)`$ and $`𝒫(𝒮(f))(t)=f(t)/f(0)`$ for $`tI_T`$.
###### Demonstration Proof
I. The indefinite product integral $`𝒫(g)`$ is defined on $`I_T`$ by Theorem 4.4 of Dudley and Norvaiša (1999). It is easy to prove that $`𝒫(g)/𝒫(g)_{}=1+\mathrm{\Delta }^{}g>0`$, $`𝒫(g)_+/𝒫(g)=1+\mathrm{\Delta }^+g>0`$ on $`I_T`$, and $`𝒫(g)𝒲_2^{}(I_T)`$. Let $`tI_T`$. By Theorem 2.1, the $`LY`$ integral $`𝒮(𝒫(g))(t)`$ exists and
$$𝒮(𝒫(g))(t)=g(t)g(0)+\mathrm{log}\left[\underset{(0,t]}{}(1+\mathrm{\Delta }^{}g)e^{\mathrm{\Delta }^{}g}\right]+\mathrm{log}\left[\underset{[0,t)}{}(1+\mathrm{\Delta }^+g)e^{\mathrm{\Delta }^+g}\right]$$
$$\underset{(0,t]}{}\left[\mathrm{log}(1+\mathrm{\Delta }^{}g)\mathrm{\Delta }^{}g\right]\underset{[0,t)}{}\left[\mathrm{log}(1+\mathrm{\Delta }^+g)\mathrm{\Delta }^+g\right]=g(t)g(0).$$
The last equality follows by taking the limit of $`\mathrm{log}[_\mu \mathrm{\Phi }]=_\mu [\mathrm{log}\mathrm{\Phi }]`$ along a nested sequence of finite sets $`\mu `$ of jump points of $`\mathrm{\Phi }=(1+\mathrm{\Delta }g)\mathrm{exp}(\mathrm{\Delta }g)`$.
II. By Theorem 2.1, $`𝒮(f)`$ is defined on $`I_T`$ and its value is given by the right side of (2.6). Since $`\mathrm{log}f𝒲_2^{}`$ and the two sums in (2.6) converge absolutely, it follows that $`𝒮(f)𝒲_2^{}`$. Let $`tI_T`$. By Theorem 4.4 of Dudley and Norvaiša (1999) and by Proposition 6 of Norvaiša (1999), the product integral $`𝒫(𝒮(f))(t)`$ exists and has the representation
$$𝒫(𝒮(f))(t)=\frac{f(t)}{f(0)}\mathrm{exp}\{\underset{(0,t]}{}[\mathrm{log}\left(\frac{f}{f_{}}\right)\frac{\mathrm{\Delta }^{}f}{f_{}}]\underset{[0,t)}{}[\mathrm{log}\left(\frac{f_+}{f}\right)\frac{\mathrm{\Delta }^+f}{f}]\}\times $$
$$\times \underset{(0,t]}{}\left[(1+\frac{\mathrm{\Delta }^{}f}{f_{}})\mathrm{exp}(\frac{\mathrm{\Delta }^{}f}{f_{}})\right]\underset{[0,t)}{}\left[(1+\frac{\mathrm{\Delta }^+f}{f})\mathrm{exp}(\frac{\mathrm{\Delta }^+f}{f})\right]=\frac{f(t)}{f(0)},$$
where the last equality follows using the limiting argument as in the part I. The proof of Theorem 2.7 is complete. ∎
By Propositions 2.3, 2.6 and Theorem 2.7, for $`𝔔=𝔔([0,T])`$ and for each $`tI_T`$, it follows that
$$_𝔔(_𝔔(f))(t)=f(t)/f(0)\text{and}_𝔔(_𝔔(g))(t)=g(t)g(0)$$
$`2.13`$
whenever $`f,g𝒲_2^{}(I_T)`$ are either right– or left–continuous at each point, $`f`$ is bounded away from zero and all jumps of $`g`$ are bigger than $`1`$. Next we partly extend the duality between $`_\lambda `$ and $`_\lambda `$ for $`\lambda 𝔔([0,T])`$ and for certain functions outside of the class $`𝒲_2^{}`$.
###### Proposition 2.8
Let $`g𝒲_p([0,T])`$, $`1p<3`$, be continuous and let $`\lambda =\{\lambda (m):m1\}`$ be a sequence of partitions $`\lambda (m)=\{0=t_0^m<\mathrm{}<t_{n(m)}^m=T\}`$ such that the mesh $`|\lambda (m)|0`$ with $`m\mathrm{}`$.
(I) For each $`t[0,T]`$, the limit
$$b_\lambda (g)(t):=\underset{m\mathrm{}}{lim}\underset{i=1}{\overset{n(m)}{}}[g(t_i^mt)g(t_{i1}^mt)]^2$$
$`2.14`$
exists if and only if (2.9) so does, and then
$$_\lambda (g)(t)=\mathrm{exp}\{g(t)g(0)2^1b_\lambda (g)(t)\}.$$
$`2.15`$
(II) Suppose $`b_\lambda (g)`$ from statement (I) is defined and continuous on $`[0,T]`$. Then, for each $`t[0,T]`$, the limit (2.4) exists for $`f=_\lambda (g)`$, and satisfies the relation
$$_\lambda (_\lambda (g))(t)=g(t)g(0).$$
$`2.16`$
###### Demonstration Proof
To prove statement (I), let $`t(0,T]`$. For each $`m1`$, let $`l(m)\{1,\mathrm{},n(m)\}`$ be an integer such that $`t(t_{l(m)1}^m,t_{l(m)}^m]`$ and let $`u_i^m:=g(t_i^mt)g(t_{i1}^m)`$ for $`i=1,\mathrm{},l(m)`$. Since $`g`$ is continuous, there exists an integer $`M`$ such that $`\mathrm{max}_i|u_i^m|1/2`$ for $`mM`$. By the Taylor series expansion with remainder, we have $`\mathrm{log}(1+u)=uu^2/2+3\theta u^3`$ for $`|u|1/2`$, where $`|\theta |=|\theta (u)|1`$. Then we get the bound
$$\left|\mathrm{log}\left(\underset{i=1}{\overset{l(m)}{}}(1+u_i^m)\right)[g(t)g(0)\frac{1}{2}s_2(g;\lambda (m))]\right|$$
$$\underset{i=1}{\overset{l(m)}{}}\left|\mathrm{log}(1+u_i^m)u_i^m+\frac{1}{2}(u_i^m)^2\right|3\underset{i=1}{\overset{l(m)}{}}|u_i^m|^33v_p(g)\underset{i}{\mathrm{max}}|u_i^m|^{3p}$$
$`2.17`$
for all $`mM`$. This yields statement (I) because $`g`$ is continuous and $`g𝒲_p([0,T])`$.
To prove statement (II), suppose that $`b:=b_\lambda (g)`$ is defined and continuous on $`[0,T]`$. Given $`t(0,T]`$, for each $`m1`$, let $`l(m)`$ be as before, $`P:=_\lambda (g)`$ and let $`v_i^m:=[P(t_i^mt)P(t_{i1}^m)]/P(t_{i1}^m)`$ for $`i=1,\mathrm{},l(m)`$. Since $`g`$ and $`b`$ are continuous, there exists an integer $`M`$ such that $`\mathrm{max}_i|v_i^m|1/2`$ for $`mM`$. As in (2.17), since $`|e^x1||x|e^{|x|}`$ for $`x`$, we get
$$\left|\mathrm{log}P(t)\underset{i=1}{\overset{l(m)}{}}v_i^m+\frac{1}{2}\underset{i=1}{\overset{l(m)}{}}[v_i^m]^2\right|3\underset{i=1}{\overset{l(m)}{}}|v_i^m|^3$$
$$Cv_p(g;[0,T])\underset{i}{\mathrm{max}}|g(t_i^mt)g(t_{i1}^m)|^{3p}+Cb(T)\underset{i}{\mathrm{max}}|b(t_i^mt)b(t_{i1}^m)|^2$$
for some constant $`C`$ and all $`mM`$. Since the right side tends to zero with $`m\mathrm{}`$, the limit (2.4) exists for $`f=P`$ because
$$\underset{m\mathrm{}}{lim}\underset{i=1}{\overset{l(m)}{}}[v_i^m]^2=\underset{m\mathrm{}}{lim}\underset{i=1}{\overset{l(m)}{}}\left[g(t_i^mt)g(t_{i1}^m)\frac{1}{2}[b(t_i^mt)b(t_{i1}^m)]\right]^2=b(t).$$
It then follows that
$$_\lambda \left(_\lambda (g)\right)(t)=\underset{m\mathrm{}}{lim}\underset{i=1}{\overset{l(m)}{}}v_i^m=\mathrm{log}P(t)+\frac{1}{2}\underset{m\mathrm{}}{lim}\underset{i=1}{\overset{l(m)}{}}[v_i^m]^2=g(t)g(0).$$
The proof of Proposition 2.8 is complete. ∎
2.4. Price and return. We define a price and its return in a duality under minimal restrictions on stochastic processes. To begin with we define a random moment $`\tau _P`$ which can be interpreted as the time of the crash of a stock. In the case when a stock price $`P`$ is the solution to the Doléans–Dade equation (1.1), $`\tau _P`$ is the first moment $`t`$ when $`P(t)0`$. Given a stochastic process $`X=\{X(t):t[0,T]\}`$, let
$$\tau _P=\tau _P(X):=\{\begin{array}{cc}t,\hfill & \text{if }inf_{s[0,t)}X(s)>0\text{ and }X(t)0\text{ for }t(0,T]\text{,}\hfill \\ t+,\hfill & \text{if }inf_{s[0,t]}X(s)>0\text{ and }X^{}(t+)0\text{ for }t(0,T)\text{,}\hfill \\ T+,\hfill & \text{if }inf_{s[0,T]}X(s)>0\text{,}\hfill \end{array}$$
$`2.18`$
where $`f^{}(t+):=limsup_{st}f(s)`$. Let $`K`$ be the set consisting of points $`0`$, $`0+`$, $`t`$, $`t`$, $`t+`$ for $`t(0,T]`$ with the natural linear ordering: $`s+<t<t<t+`$ if $`s<t`$, and endowed with the interval topology. Then $`\tau _P`$ is the random variable with values in $`K`$. If $`X`$ is a price, then the event $`\{\tau _P=t\}`$ can be interpreted as the crash right before the time $`t`$, while the event $`\{\tau _P=t+\}`$ can be interpreted as the crash right after the time $`t`$. Also, let $`[0,t+):=[0,t]`$.
Recalling notation $`𝔔([0,T])`$, $`_\lambda `$ and $`_\lambda `$ from Definitions 2.2 and 2.5, we have:
###### Definition Definition 2.9
Let $`R=\{R(t):t[0,T]\}`$ and $`P=\{P(t):t[0,T]\}`$ be stochastic processes on a complete probability space $`(\mathrm{\Omega },,\mathrm{Pr})`$ such that $`R(0)=0`$ and $`P(0)=1`$ almost surely.
If the pair $`(P,R)`$ is the weak evolutionary system then we call $`P`$ the price and $`R`$ the return.
We show that if $`R`$ is a Brownian motion and $`P`$ is a geometric Brownian motion, then the pair $`(P,R)`$ is the weak evolutionary system but not the evolutionary system (Proposition 2.11 and Remark 2.12). However, according to the following statement, if almost all sample functions of $`R`$ are in $`𝒲_2^{}([0,T])𝒟([0,T])`$ then $`(P,R)`$ is the evolutionary system. Recall that $`f𝒟([0,T])`$ if, at each point of $`(0,T)`$, $`f`$ is is either right–continuous or left–continuous.
Next we define a random moment $`\tau _R`$ for a return. Given a stochastic process $`Y=\{Y(t):t[0,T]\}`$ with almost all sample functions in $`𝒟([0,T])`$, let
$$\tau _R=\tau _R(Y):=\{\begin{array}{cc}t,\hfill & \text{if }inf_{s[0,t)}\mathrm{\Delta }Y(s)>1\text{ and }\mathrm{\Delta }^{}Y(t)1\text{ for }t(0,T]\text{,}\hfill \\ t+,\hfill & \text{if }inf_{s[0,t)}\mathrm{\Delta }Y(s)>1\text{ and }\mathrm{\Delta }^+Y(s)1\text{ for }t(0,T)\text{,}\hfill \\ T+,\hfill & \text{if }inf_{s[0,T]}\mathrm{\Delta }Y(s)>1\text{.}\hfill \end{array}$$
Here $`\mathrm{\Delta }Y(s):=0`$ everywhere except at jump points $`s`$ of $`Y`$ where either $`\mathrm{\Delta }Y(s):=\mathrm{\Delta }^+Y(s)`$ if it is non-zero, or $`\mathrm{\Delta }Y(s):=\mathrm{\Delta }^{}Y(s)`$ if it is non-zero.
###### Proposition 2.10
Let $`R=\{R(t):t[0,T]\}`$ be a stochastic process on a complete probability space $`(\mathrm{\Omega },,\mathrm{Pr})`$ such that $`R(0)=0`$ and $`R𝒲_2^{}([0,T])𝒟([0,T])`$ with probability $`1`$. Then the indefinite product integral $`𝒫(R)`$ is defined with respect to almost every sample function of $`R`$, $`\tau _P(𝒫(R))=\tau _R(R)`$ almost surely and the pair $`(𝒫(R),R)`$ is the evolutionary system on $`[0,\tau _R)`$.
###### Demonstration Proof
Let $`N`$ be such that $`\mathrm{Pr}(N)=0`$ and $`R(,\omega )𝒲_2^{}([0,T])𝒟([0,T])`$ for all $`\omega \mathrm{\Omega }N`$. Let $`\lambda =\{\lambda (m):m1\}𝔔([0,T])`$. For each $`t<\tau (\omega ):=\tau _R(R(\omega ))`$, let $`P(t,\omega ):=_\lambda (R(,\omega ))(t)`$ if $`\omega \mathrm{\Omega }N`$ and $`P(t,\omega ):=0`$ if $`\omega N`$. By Proposition 2.6, $`P(t,\omega )=\text{}_0^t(1+dR(,\omega ))`$ for $`t<\tau (\omega )`$ and $`\omega \mathrm{\Omega }N`$. For all $`\omega \mathrm{\Omega }`$ and $`t\tau (\omega )`$, let $`P(t,\omega ):=\text{}_0^{\tau (\omega )}(1+dR(,\omega ))`$. Then $`P=\{P(t):t[0,T]\}`$ is a stochastic process by construction. Fix $`\omega \mathrm{\Omega }N`$ and let $`\tau :=\tau (\omega )`$, $`P(t):=P(t,\omega )`$, $`R(t):=R(t,\omega )`$. Then $`P(\tau )`$ exists and $`inf\{P(t):t[0,\tau )\}\delta `$ for some $`\delta >0`$ by Proposition 4.30 of Dudley and Norvaiša (1999). By Lemmas 5.1 and 5.2 of Dudley and Norvaiša (1999), $`P`$ and $`R`$ have the same jump points on $`[0,\tau )`$ and thus $`P𝒟([0,\tau ))`$. Moreover, $`P𝒲_2^{}([0,\tau ))`$ by Theorem 2.7. By Proposition 2.3, $`\stackrel{~}{R}(t):=_\lambda (P)(t)`$ is defined for each $`t<\tau `$ and has the same value with the indefinite $`LY`$ integral $`𝒮(P)(t)`$. By Theorem 2.7 again, $`\stackrel{~}{R}(t)=(LY)_0^tP^1𝑑P=R(t)`$ for each $`t<\tau `$. Since the null set $`N`$ does not depend on $`\lambda `$, the pair $`(P,R)`$ is evolutionary system on $`[0,\tau _P)`$. It is clear that $`\tau _P(P(\omega ))=\tau _R(R(\omega ))`$ for all $`\omega \mathrm{\Omega }N`$. The proof of Proposition 2.10 is complete. ∎
For example, a fractional Brownian motion $`B_H`$ with $`H(1/2,1)`$ and a symmetric $`\alpha `$-stable Lévy motion $`X_\alpha `$ with $`\alpha (0,2)`$ are the returns of the evolutionary systems $`(P_H,B_H)`$ and $`(P_\alpha ,X_\alpha )`$, respectively, where
$$P_H(t):=\mathrm{exp}\{B_H(t)\}\text{and}P_\alpha (t):=\mathrm{exp}\{X_\alpha (t)\}\underset{(0,t]}{}(1+\mathrm{\Delta }^{}X_\alpha )\mathrm{exp}\{\mathrm{\Delta }^{}X_\alpha \}$$
for $`t[0,T]`$. The price $`P_\alpha `$ is positive until the first moment $`t`$ when $`\mathrm{\Delta }^{}X_\alpha (t)1`$.
###### Proposition 2.11
Let $`B=\{B(t):t[0,T]\}`$ be a standard Brownian motion and let $`P_B:=\{\mathrm{exp}\{B(t)t/2\}:t[0,T]\}`$. Then the pair $`(P_B,B)`$ is the weak evolutionary system on $`[0,T]`$.
###### Remark Remark 2.12
The proof of the above proposition rely on Théorème 5 of Lévy (1940, p. 510): for a standard Brownian motion $`B`$ and for a sequence $`\lambda =\{\lambda (m):m1\}𝔔([0,1])`$, the limit $`lim_m\mathrm{}s_2(B;\lambda (m))=1`$ exists with probability 1. However, the exceptional null set $`N(\lambda )`$ of this implication depends on $`\lambda `$ and $`\{N(\lambda ):\lambda 𝔔([0,1])\}=\mathrm{\Omega }`$. Moreover, for almost all $`\omega \mathrm{\Omega }`$ there exist $`\lambda 𝔔([0,1])`$ such that $`lim_m\mathrm{}s_2(B(,\omega );\lambda (m))=\mathrm{}`$, and hence, $`_\lambda (B(,\omega ))=0`$. The proofs of these properties are given by Freedman (1983, p. 48) because $`_m\lambda (m)`$ is everywhere dense in $`[0,1]`$ if and only if the mesh $`|\lambda (m)|0`$ with $`m\mathrm{}`$.
###### Demonstration Proof
The claim will follow from Proposition 2.8 once we show that, given $`\lambda =\{\lambda (m):m1\}𝔔([0,T])`$, the limit (2.14) with $`g=B`$ exists with probability 1 for each $`t[0,T]`$. For each $`t(0,T]`$ and $`m1`$, let $`b_m(t,\omega ):=_{i=1}^{n(m)}\left[B(t_i^mt,\omega )B(t_{i1}^mt,\omega )\right]^2`$. Let $`N=N(\lambda )`$ be a null set such that, for each $`\omega \mathrm{\Omega }N`$, $`B(,\omega )`$ is continuous function of bounded $`p`$-variation for some $`2<p<3`$, and the limit $`lim_mb_m(t,\omega )=t`$ exists for each $`t`$ in the countable set $`S:=_m\lambda (m)`$. Let $`t(0,T)S`$. Since $`b_m(t,\omega )t=[B(t,\omega )B(t_{i1}^m,\omega )]^2+b_m(t_{i1}^m,\omega )t_{i1}^m+[t_{i1}^mt]`$ for $`t_{i1}^m<t<t_i^m`$, it follows that the limit as $`m\mathrm{}`$ of $`b_m(t,\omega )`$ is $`t`$ for each $`t[0,T]`$. Therefore an appeal to Proposition 2.8 completes the proof. ∎
3 Pathwise trading strategies
A trading strategy is a collection of instructions for buying and selling a stock, depending on its price fluctuations. A mathematical notion of a trading strategy should be defined so that one can calculate the portfolio value and portfolio gain for each single trajectory of a stock price. In the stochastic exponential model, the portfolio gain is the stochastic integral of a trading strategy with respect to a price. Its value can be approximated by portfolio gains based on simple trading strategies in probability. In this section we take a pathwise approach to trading strategies.
Consider a frictionless stock market with $`\nu +1`$ non-dividend-paying stocks and open for trade during a time period $`[0,T]`$. A vector function $`P=(P_0,\mathrm{},P_\nu )`$ defined on $`[0,T]`$ will be called the price during the time period $`[0,T]`$ if, for each $`k=0,\mathrm{},\nu `$, $`inf\{P_k(t):t[0,T]\}>0`$. The value $`P_k(t)`$ refers to the price of the $`k`$th stock at time $`t[0,T]`$ for $`k=0,\mathrm{},\nu `$. For example, the price may be a vector of sample functions of the price stochastic processes defined in the preceding section. A possible dependence of the price on a randomness and the duality relation between the price and its return are disregarded in this section.
As before, $`𝔔([0,T])`$ denotes the set of all nested sequences $`\lambda =\{\lambda (m):m1\}`$ of partitions of $`[0,T]`$ such that $`_m\lambda (m)`$ is dense in $`[0,T]`$. A sequence $`\kappa =\{\kappa (m):m1\}𝔔([0,T])`$ is a refinement of a sequence $`\lambda =\{\lambda (m):m1\}𝔔([0,T])`$ if each $`\kappa (m)`$ is a refinement of $`\lambda (m)`$.
###### Definition Definition 3.1
Let $`P=(P_0,\mathrm{},P_\nu )`$ be a price during a time period $`[0,T]`$. Given $`\lambda 𝔔([0,T])`$, a vector function $`\varphi =(\varphi _0,\mathrm{},\varphi _\nu )`$ defined on $`[0,T]`$ will be called the $`(\lambda ,P)`$–trading strategy during the time period $`[0,T]`$ if, for each $`k=0,\mathrm{},\nu `$ and $`t[0,T]`$, there exists the limit
$$(LCS)_0^t\varphi _kd_\lambda P_k:=\underset{m\mathrm{}}{lim}\underset{i=1}{\overset{n(m)}{}}\varphi _k(t_{i1}^mt)[P_k(t_i^mt)P_k(t_{i1}^mt)],$$
$`3.1`$
where $`\lambda =\{\lambda (m):m1\}`$ and $`\lambda (m)=\{t_i^m:i=0,\mathrm{},n(m)\}`$. If there exists $`\lambda _0𝔔([0,T])`$ such that for each refinement $`\lambda 𝔔([0,T])`$ of $`\lambda _0`$, $`\varphi `$ is the $`(\lambda ,P)`$–trading strategy on $`[0,T]`$ and (3.1) does not depend on $`\lambda `$, then we call $`\varphi `$ the $`P`$–trading strategy during the time period $`[0,T]`$ and replace $`d_\lambda `$ with $`d`$ in the left side of (3.1).
Föllmer (1981) proved that (3.1) exists whenever $`\varphi _k=fP_k`$ for some $`fC^1`$ and the quadratic variation is defined for $`P_k`$ along the sequence $`\lambda 𝔔([0,T])`$. The above notion of $`(\lambda ,P)`$–trading strategy is similar to the convergence of trading strategies introduced by Bick and Willinger (1994, p. 356). These authors derived the Black and Scholes formula without probabilistic arguments using Föllmer’s variant of Itô’s formula. Remark 2.12 ensure that $`\varphi `$ may be $`(\lambda ,P)`$–trading strategy without being a $`P`$-trading strategy. It is clear that each $`\varphi _k`$ is Left Cauchy–Stieltjes integrable with respect to $`P_k`$ if $`\varphi `$ is $`P`$–trading strategy. As the rest of this section show we could use the $`LY`$ integral instead of the $`LCS`$ integral in the definition of $`P`$–trading strategies. However in econometric analysis it seems easier to handle with the latter integral because to evaluate the $`LY`$ integral we would need to know jumps of the price. The following statement provides sufficient conditions for a vector function to be the $`P`$–trading strategy.
###### Proposition 3.2
Let $`P=(P_0,\mathrm{},P_\nu )`$ be a price during a time period $`[0,T]`$. A vector function $`\varphi =(\varphi _0,\mathrm{},\varphi _\nu )`$ on $`[0,T]`$ is the $`P`$–trading strategy, each $`\varphi _k`$ is $`LY`$ integrable with respect to $`P_k`$ on $`[0,T]`$ and
$$(LCS)_0^t\varphi _k𝑑P_k=(LY)_0^t\varphi _k𝑑P_k$$
$`3.2`$
for each $`t[0,T]`$ and each $`k=0,\mathrm{},\nu `$ in either of the following two cases:
###### Demonstration Proof
In case (2) the conclusion follows from Theorem 2 and Corollary 3 of Norvaiša (1999). To prove the conclusion in case (1), for notation simplicity we suppress the index $`k=0,\mathrm{},\nu `$ for $`f_k`$ and $`P_k`$. By Theorem 2.1, $`fP`$ is $`LY`$ integrable with respect to $`P`$ on $`[0,T]`$ and
$$(LY)_0^t(fP)𝑑P=FP(t)FP(0)\underset{(0,t]}{}\varphi ^{}\underset{[0,t)}{}\varphi ^+,$$
where $`\varphi ^{}:=\mathrm{\Delta }^{}(FP)(fP)_{}\mathrm{\Delta }^{}P`$, $`\varphi ^+:=\mathrm{\Delta }^+(FP)(fP)\mathrm{\Delta }^+P`$ and $`F(u):=_0^uf(x)𝑑x`$ for $`u0`$. Fix $`t(0,T]`$ and let $`ϵ>0`$. By Lemma 2.4, there exists $`\lambda =\{s_j:j=0,\mathrm{},m\}Q([0,t])`$ such that $`_{j=1}^mv_2(P;(s_{j1},s_j))<ϵ`$, and for each refinement $`\{t_i:i=0,\mathrm{},n\}`$ of $`\lambda `$,
$$\left|\underset{(0,t]}{}\varphi ^{}+\underset{[0,t)}{}\varphi ^+\underset{i=1}{\overset{n}{}}\left[\varphi ^{}(t_i)+\varphi ^+(t_{i1})\right]\right|<ϵ.$$
Then choose $`\{u_{j1},v_j:j=1,\mathrm{},m\}[0,t]`$ such that $`s_{j1}<u_{j1}<v_j<s_j`$ for $`j=1,\mathrm{},m`$ and
$$\underset{j=1}{\overset{m}{}}\left[Osc(P;(z_{j1},u_{j1}])+Osc(P;[v_j,z_j))\right]<ϵ.$$
Then by the mean value theorem and using the Lipschitz condition with the constant $`K`$, we have
$$\left|S_{LCS}(fP,P;\kappa )\left[FP(t)FP(0)\underset{(0,t]}{}\varphi ^{}\underset{[0,t)}{}\varphi ^+\right]\right|$$
$$<ϵ+\underset{i=1}{\overset{n}{}}\left|fP(t_{i1})[P(t_i)P(t_{i1})][FP(t_i)FP(t_{i1})][\varphi ^{}(t_i)\varphi ^+(t_{i1})]\right|$$
$$=ϵ+\underset{i=1}{\overset{n}{}}|[f(\theta _i)f(P(t_i))][P(t_i)P(t_{i1}+)]+[fP(t_i)fP(t_{i1})]$$
$$\times [P(t_i)P(t_{i1}+)]|<ϵ+Kϵ+2\underset{t}{sup}|P(t)|Kϵ+Kϵ,$$
where $`\theta _i[P(t_{i1}+)P(t_i),P(t_{i1}+)P(t_i)]`$. Since $`ϵ`$ is arbitrary the proof of Proposition 3.2 is complete. ∎
Having defined the $`P`$–trading strategies via the Left Cauchy–Stieltjes integral, the following definition of self–financing strategy corresponds naturally to its counterpart in the stochastic exponent model.
###### Definition Definition 3.3
Let $`P=(P_0,\mathrm{},P_\nu )`$ be a price during a time period $`[0,T]`$ and let $`\varphi =(\varphi _0,\mathrm{},\varphi _\nu )`$ be the $`P`$–trading strategy.
For the sake of illustration, next we partially extend Proposition 3.24 of Harrison and Pliska (1981) to the present setting. Let us denote the discounted price $`(1,P_1/P_0,\mathrm{},P_\nu /P_0)`$ by $`\overline{P}`$.
###### Proposition 3.4
Let $`0<p<2`$ and let $`P=(P_0,\mathrm{},P_\nu )`$ be a price during a time period $`[0,T]`$ such that $`P_k𝒲_p([0,T])`$ for $`k=0,\mathrm{},\nu `$. Suppose that $`\varphi =(\varphi _0,\mathrm{},\varphi _\nu )`$ is a vector function on $`[0,T]`$ such that $`\varphi _k𝒲_p([0,T])`$ for $`k=0,\mathrm{},\nu `$. Then $`\varphi `$ is self–financing $`P`$–trading strategy if and only if it is self–financing $`\overline{P}`$–trading strategy.
###### Demonstration Proof
Let $`\beta :=1/P_0`$. We have $`\beta P_k𝒲_p([0,T])`$ for each $`k=1,\mathrm{},\nu `$. Thus, by case (2) of Proposition 3.2 with $`q=p`$, $`\varphi `$ is the $`P`$–trading strategy if and only if $`\varphi `$ is the $`\overline{P}`$–trading strategy
To prove the “only if” part of the statement suppose $`\varphi `$ is self–financing $`P`$-trading strategy. Let $`V:=V^{\varphi ,P}`$ and $`\overline{V}:=V^{\varphi ,\overline{P}}=\beta V`$. By Theorem 2.1 with $`F(u)=u_1u_2`$ for $`u=(u_1,u_2)`$, $`h=\varphi _k`$, and $`\alpha =1`$, for each $`t[0,T]`$, we have
$$(LY)_0^t\varphi _kd(\beta P_k)=(LY)_0^t\varphi _k\beta 𝑑P_k+(LY)_0^t\varphi _kP_k𝑑\beta +\underset{(0,t]}{}(\varphi _k)_{}\mathrm{\Delta }^{}\beta \mathrm{\Delta }^{}P_k$$
$`3.3`$
for $`k=0,\mathrm{},\nu `$, where the left side of (3.3) is equal to 0 when $`k=0`$. Since $`P`$–trading strategy $`\varphi `$ is self–financing, by the substitution rule for the Left Young integral (Theorem 9 of Norvaiša, 1999), for each $`t[0,T]`$, we have
$$(LY)_0^t\beta 𝑑V=\underset{k=0}{\overset{\nu }{}}(LY)_0^t\beta d\left((LY)_0^{}\varphi _k𝑑P_k\right)=\underset{k=0}{\overset{\nu }{}}(LY)_0^t\beta \varphi _k𝑑P_k.$$
$`3.4`$
Using Theorem 2.1 again except that now $`h=1`$, and Proposition 7 of Norvaiša (1999) about jumps of the indefinite $`LY`$ integral, for each $`t[0,T]`$, we get
$`\overline{V}(t)\overline{V}(0)`$ $`=(LY){\displaystyle _0^t}\beta 𝑑V+(LY){\displaystyle _0^t}V𝑑\beta +{\displaystyle \underset{(0,t]}{}}\mathrm{\Delta }^{}\beta \mathrm{\Delta }^{}V`$
by (3.4) $`={\displaystyle \underset{k=0}{\overset{\nu }{}}}\left\{(LY){\displaystyle _0^t}\varphi _k\beta 𝑑P_k+(LY){\displaystyle _0^t}\varphi _kP_k𝑑\beta +{\displaystyle \underset{(0,t]}{}}(\varphi _k)_{}\mathrm{\Delta }^{}\beta \mathrm{\Delta }^{}P_k\right\}`$
by (3.3) $`={\displaystyle \underset{k=1}{\overset{\nu }{}}}(LY){\displaystyle _0^t}\varphi _kd(\beta P_k)=G^{\varphi ,\overline{P}}(t).`$
Thus $`\varphi `$ is self–financing $`\overline{P}`$–trading strategy. The proof of the converse implication is similar and therefore is omitted. ∎
We finish with the main argument in favor of the pathwise approach to trading strategies. In their discussion of the notion of trading strategy, Harrison and Pliska (1981, Section 7) made several suggestions. For example, it would be desirable to show that a claim is attainable if and only if it is the limit (in some appropriate sense) of claims generated by simple self–financing strategies. Duffie and Protter (1992) and Eberlein (1992) proved that the portfolio gain processes are approximable by their discrete counterparts under certain conditions. Next we show a kind of approximation of a “contingent claim” by simple self–financing strategies in the present context. A trading strategy $`\varphi =(\varphi _0,\mathrm{},\varphi _\nu )`$ is simple if each $`\varphi _k`$ is a step function on $`[0,T]`$. The idea of the following statement originated from Harrison, Pitbladdo and Schaefer (1984, Proposition 9), where this claim is proved for price processes with continuous sample functions of bounded variation.
###### Theorem 3.5
Let $`P`$ be a price during a time period $`[0,T]`$ and let $`\varphi `$ be a vector function on $`[0,T]`$, both satisfying either of the two conditions of Proposition 3.2. If $`inf\{V^\varphi (t):t[0,T]\}>0`$ then there exists a sequence $`\{\varphi ^N:N1\}`$ of simple self–financing $`P`$–trading strategies such that $`V^{\varphi ^N}(0)=V^\varphi (0)`$ and $`lim_N\mathrm{}V^{\varphi ^N}(T)=V^\varphi (T)`$.
###### Demonstration Proof
We start with the construction of the sequence of simple self–financing $`P`$–trading strategies based on a given nested sequence $`\{\lambda ^N:N1\}`$ of partitions $`\lambda ^N=\{0=t_0<\mathrm{}<t_n=T\}`$. Given an integer $`N1`$, we define $`\varphi ^N=(\varphi _0^N,\mathrm{},\varphi _\nu ^N)`$ recursively with constant values on each interval of the partition $`\lambda ^N`$. For each $`k=0,\mathrm{},\nu `$, let $`\varphi _k^N:=\varphi _k(0)`$ on $`[0,t_1)`$. Suppose that all $`\varphi _k^N`$ are defined on $`[0,t_i)`$ for some $`1in`$. Let $`V^N(t_i):=_{k=0}^\nu \varphi _k^N(t_{i1})P_k(t_i)`$. Then for each $`k=0,\mathrm{},\nu `$, let $`\varphi _k^N`$ be equal to $`\varphi _k(t_i)V^N(t_i)/V^\varphi (t_i)`$ either on $`[t_i,t_{i+1})`$ if $`i<n`$, or on $`\{T\}`$ if $`i=n`$. It is clear that each $`\varphi ^N`$ is a simple $`P`$–trading strategy. Moreover, the portfolio value function $`V^{\varphi ^N}`$ has values $`V^{\varphi ^N}(0)=V^\varphi (0)`$ and
$$V^{\varphi ^N}(t_i)=\underset{k=0}{\overset{\nu }{}}\varphi _k^N(t_i)P_k(t_i)=\frac{V^N(t_i)}{V^\varphi (t_i)}\underset{k=0}{\overset{\nu }{}}\varphi _k(t_i)P_k(t_i)=V^N(t_i)$$
$`3.5`$
for each $`i=1,\mathrm{},n`$. Next we show that each $`P`$–trading strategy $`\varphi ^N`$ is self–financing. Let $`u=t_{i1}`$ and $`v(t_{i1},t_i]`$ for some $`i=1,\mathrm{},n`$. Since $`\varphi _k^N`$ is constant on $`[u,v)`$ we get
$$(LY)_u^v\varphi _k^N𝑑P_k=(RS)_u^v(\varphi _k^N)_{}^{(u)}d(P_k)_+^{(v)}+\varphi _k^N(u)\mathrm{\Delta }^+P_k(u)$$
$$=\varphi _k^N(u)[P_k(v)P_k(u+)]+\varphi _k^N(u)\mathrm{\Delta }^+P_k(u)=\varphi _k^N(u)[P_k(v)P_k(u)]$$
for each $`k=0,\mathrm{},\nu `$. Given $`t(0,T]`$, let $`l:=\mathrm{max}\{in:t_it\}`$. Then using the additivity of the $`LY`$ integral over adjacent intervals (Theorem 4 of Norvaiša, 1999) and changing the order of summation over $`k`$ and $`i`$, we get
$`G^{\varphi ^N}(t)`$ $`={\displaystyle \underset{k=0}{\overset{\nu }{}}}\left\{(LY){\displaystyle _{t_l}^t}\varphi _k^N𝑑P_k+{\displaystyle \underset{i=1}{\overset{l}{}}}(LY){\displaystyle _{t_{i1}}^{t_i}}\varphi _k^N𝑑P_k\right\}`$
$`={\displaystyle \underset{k=0}{\overset{\nu }{}}}\varphi _k^N(t_l)[P_k(t)P_k(t_l)]+{\displaystyle \underset{i=1}{\overset{l}{}}}\left[V^N(t_i)V^{\varphi ^N}(t_{i1})\right]`$
by (3.5) $`={\displaystyle \underset{k=0}{\overset{\nu }{}}}\left[\varphi _k^N(t_{i(t)})P_k(t)\varphi _k^N(0)P_k(0)\right]=V^{\varphi ^N}(t)V^{\varphi ^N}(0).`$
Thus the $`P`$–trading strategy $`\varphi ^N`$ is self–financing for each $`N1`$. By (3.5) and by the additivity of the $`LY`$ integral over adjacent intervals again, we get
$$\mathrm{\Delta }^N(t_{i1},t_i):=V^\varphi (t_i)V^\varphi (t_{i1})\frac{V^{\varphi ^N}(t_i)}{V^{\varphi ^N}(t_{i1})}=V^\varphi (t_i)\underset{k=0}{\overset{\nu }{}}\varphi _k(t_{i1})P_k(t_i)$$
$$=\underset{k=0}{\overset{\nu }{}}\{(LY)_{t_{i1}}^{t_i}\varphi _kdP_k\varphi _k(t_{i1})[P_k(t_i)P_k(t_{i1})]\}=:\underset{k=0}{\overset{\nu }{}}\mathrm{\Delta }_k^N(t_{i1},t_i)$$
$`3.6`$
for each $`i=1,\mathrm{},n`$. Suppose that one can choose a sequence $`\{\lambda ^N:N1\}`$ such that
$$\underset{1in}{\mathrm{max}}|\mathrm{\Delta }^N(t_{i1},t_i)|\underset{i=1}{\overset{n}{}}|\mathrm{\Delta }^N(t_{i1},t_i)|ϵ_N$$
$`3.7`$
for some $`ϵ_N0`$. Then it follows that
$$\left|\frac{V^{\varphi ^N}(t_i)}{V^\varphi (t_i)}\frac{V^\varphi (t_{i1})}{V^{\varphi ^N}(t_{i1})}1\right|=\left|\frac{\mathrm{\Delta }^N(t_{i1},t_i)}{V^\varphi (t_i)}\right|ϵ_N/\delta $$
for each $`t_{i1},t_i\lambda ^N`$. We conclude then recursively that $`V^N(t_i)0`$ for $`i=1,\mathrm{},n`$ whenever $`2ϵ_N\delta `$. By the mean value theorem, $`|\mathrm{log}(1+u)|2|u|`$ for each $`|u|1/2`$. Thus, for all $`N`$ such that $`2ϵ_N\delta `$, using the telescoping sum representation, we get
$`\left|\mathrm{log}{\displaystyle \frac{V^{\varphi ^N}(T)}{V^\varphi (T)}}\right|`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}\left|\mathrm{log}{\displaystyle \frac{V^{\varphi ^N}(t_i)}{V^\varphi (t_i)}}\mathrm{log}{\displaystyle \frac{V^{\varphi ^N}(t_{i1})}{V^\varphi (t_{i1})}}\right|2{\displaystyle \underset{i=1}{\overset{n}{}}}\left|{\displaystyle \frac{V^{\varphi ^N}(t_i)}{V^\varphi (t_i)}}{\displaystyle \frac{V^\varphi (t_{i1})}{V^{\varphi ^N}(t_{i1})}}1\right|`$
$`{\displaystyle \frac{2}{\delta }}{\displaystyle \underset{i=1}{\overset{n}{}}}|\mathrm{\Delta }^N(t_{i1},t_i)|2ϵ_N/\delta ,`$
where the last inequality follows from the second inequality in (3.7). Since $`V^{\varphi ^N}(0)=V^\varphi (0)`$ by construction, this yields the conclusion of the theorem.
It remains to find $`\{\lambda ^N:N1\}`$ such that (3.7) holds for some $`ϵ_N0`$. Suppose that condition (1) of Proposition 3.2 holds. By Theorem 2.1 and by the mean value theorem, each term $`\mathrm{\Delta }_k^N(u,v)`$ in (3.6) with $`u=t_{i1}`$, $`v=t_i`$, $`i=1,\mathrm{},n`$ and $`k=0,\mathrm{},\nu `$ is equal to
$$[f_k(\theta _k)f_k(P_k(v))][P_k(v)P_k(u+)]+[f_k(P_k(v))f_k(P_k(u))][P_k(v)P(u+)]$$
$$\underset{(u,v)}{}\left[\mathrm{\Delta }^{}(F_kP_k)(f_kP_k)_{}\mathrm{\Delta }^{}P_k\right]\underset{(u,v)}{}\left[\mathrm{\Delta }^+(F_kP_k)(f_kP_k)\mathrm{\Delta }^+P_k\right],$$
where $`\theta _k(P_k(u+)P_k(v),P_k(u+)P_k(v)]`$ and $`F_k(u)=_0^uf_k(x)𝑑x`$ for $`u0`$. Then, given $`ϵ_N0`$, one can find $`\lambda ^N`$ such that (3.7) holds just as in the proof of Proposition 2.3. Finally, suppose that condition (2) of Proposition 3.2 holds. Choose $`p^{}>p`$ and $`q^{}>q`$ so that $`1/p^{}+1/q^{}>1`$. By the Love–Young inequality (p. 256 in Young, 1936), each term $`\mathrm{\Delta }_k^N(u,v)`$ in (3.6) with $`u=t_{i1}`$, $`v=t_i`$, $`i=1,\mathrm{},n`$ and $`k=0,\mathrm{},\nu `$ can be bounded as follows:
$`|\mathrm{\Delta }_k^N(u,v)|`$ $`=\left|(RS){\displaystyle _u^v}[(\varphi _k)_{}^{(u)}(\varphi _k)_{}^{(u)}(u)]d(P_k)_+^{(v)}+{\displaystyle \underset{(u,v)}{}}\mathrm{\Delta }^{}\varphi _k\mathrm{\Delta }^+P_k\right|`$
$`KV_p^{}(P_k;[u+,v])V_q^{}(\varphi _k;[u,v])+{\displaystyle \underset{(u,v)}{}}|\mathrm{\Delta }^{}\varphi _k\mathrm{\Delta }^+P_k|`$
for some finite constant $`K`$ depending on $`p^{}`$ and $`q^{}`$ only. Again, given $`ϵ_N0`$, one can find $`\lambda ^N`$ such that (3.7) holds as in the proof of Proposition 2.3. The proof of Theorem 3.5 is complete. ∎
4 Implications and conclusions
The results of the present paper provide an alternative construction of a stock price model, and show that many concrete financial models can be treated using classical calculus. By its definition, the evolutionary system is the continuous–time model obtained as the limit of the discrete–time model (2.2) along a sequence of partitions of a time interval into shrinking subintervals. The evolutionary system separates analytical and probabilistic aspects of analysis casting new light on important problems of stock price modelling. The arbitrage construction in Subsection 1.5 illustrates implications of this separation.
New definition of the return $`R`$ in the evolutionary system $`(P,R)`$ makes easier to use it in a statistical analysis as compared with the definition (1.2). Statistical analysis of analytical properties of functions developed in relation to natural sciences could be applied in econometric analysis of the evolutionary system. For example, an interesting task is to distinguish the hypotheses that the $`p`$-variation index $`\upsilon (R)<2`$ against the hypotheses that $`\upsilon (R)2`$. This is important because the value $`\upsilon (R)=2`$ separates a fundamentally different behaviour of $`R`$. Also, testing hypotheses $`\upsilon (R)<2`$ and $`R`$ is continuous could be used to test market efficiency related to arbitrage. Naturally that there are no ready to use statistical tests for estimating the $`p`$-variation index. In this case one needs to extract from data an information about a local behavior of a sample function rather than an information about tail distribution, or correlation estimates. The first step in this direction has been taken up by Norvaiša and Salopek (1999). These authors suggest a statistic based on old results of G. Baxter and E.G. Gladyshev concerning quadratic variation for Gaussian processes. Also, they compare the results of data analysis using the new definition of the return and the log return.
Let dim$`{}_{HB}{}^{}(G)`$ be the Hausdorff–Besicovitch dimension of a set $`G`$. Then for a large class of stochastic processes, the relation dim$`{}_{HB}{}^{}(`$graph $`X)=21/(1\upsilon (X))`$ holds for almost all sample functions of $`X`$. This fact can be used to construct new statistics for estimating the $`p`$-variation index $`\upsilon (X)`$ because statistical analysis of fractal dimensions is already highly developed in various natural sciences. The real analysis approach to modelling of stock price changes provides a new meaning to stylized facts discovered in Econophysics (see e.g. Bouchaud and Potters, 1999), and opens a way for exploring new tools for investigating financial markets.
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# Cotangent and tangent modules on quantum orbits
## 1 Introduction
Since the creation of super-theory it became clear that numerous aspects of commutative algebra and usual geometry could be generalized to the super-case. In particular, for any two $`Z_2`$-graded one-sided (say, left) $`A`$-modules $`M_1`$ and $`M_2`$ over a super-commutative algebra $`A`$ their tensor product $`M_1_AM_2`$ is well-defined.
In the latter 80’s it was recognized that many properties of (super-)commutative algebra and geometry could be further generalized onto objects related to a Yang-Baxter (YB) operator, i.e., a solution of the quantum YB equation
$$(S\mathrm{id})(\mathrm{id}S)(S\mathrm{id})=(\mathrm{id}S)(S\mathrm{id})(\mathrm{id}S),$$
$`S`$ being an operator acting on $`V^2`$ where $`V`$ is a vector space. If $`S`$ is an involutary YB operator $`(S^2=\mathrm{id})`$, the notion of an $`S`$-commutative algebra can be introduced in a natural way. If $`A`$ is such an algebra, the product $`M_1_AM_2`$ of two one-sided $`A`$-modules can be introduced by means of the operator $`S`$ (under some natural conditions on it). Also, the operator $`S`$ plays the crucial role in a twisted or quantum (i.e. related to an operator $`S`$) version of differential calculus. It can be used for ordering ”quantum functions” and differentials mixed in virtue of the Leibniz rule.
Unfortunately, a straightforward application of this method to algebras related to a non-involutary YB operator $`S`$ leads to non-flat deformations<sup>1</sup><sup>1</sup>1Let us recall that a deformation $`V_{\mathrm{}}`$ of a vector space $`V`$ where $`\mathrm{}`$ is a formal parameter is called flat if $`V_{\mathrm{}}/\mathrm{}V_{\mathrm{}}=V`$ and $`V_{\mathrm{}}`$ is isomorphic to $`Vk[[\mathrm{}]]`$ as $`k[[\mathrm{}]]`$-module (the tensor product is completed in the $`\mathrm{}`$-adic topology)..
As was shown in \[Ar\], \[AAM\] the differential calculus on the quantum algebras
$$k_q(G),G=SL(n),SO(n),Sp(n)$$
initiated in \[W1\] and \[W2\] and making use of the Leibniz rule gives rise to non-flat deformations of the classical differential algebras. However, this differential calculus plays the central role in all known attempts to introduce a quantum version of gauge theory related to Drinfeld-Jimbo quantum groups (cf., i.e. \[BM\], \[HM\] and the references therein). From our viewpoint, non-flat deformations are somewhat pointless since in this case classical objects are not limits of their quantum counterparts.
Quantum orbits pertain to another type of algebras related to quantum groups. The simplest example of such an orbit is the quantum sphere $`k(S_q^2)`$ introduced in \[P1\]. A version of differential calculus on it was suggested in \[P2\]. However, as follows from our results, on the quantum sphere no flatly deformed $`U_q(su(2))`$-covariant differential calculus exists which makes use of the Leibniz rule. An explanation of this phenomenon, given in the present paper, consists in the following. It is possible to realize a flat deformation of the cotangent module<sup>2</sup><sup>2</sup>2We use this term for the module corresponding to the tangent vector bundle in the framework of the Serre-Swan approach, cf. \[Se\], \[Sw\]. In a similar way we consider other classical and quantum modules. (All modules are assumed to be finitely generated.) considered as a one-sided $`k(S^2)`$-module. However, the flatness of deformation breaks down when one tries to deform the tangent vector bundle considered as a two-sided $`k(S^2)`$-module.
Thus, if we want deformed modules to be flatly deformed objects, the use of one-sided modules on quantum varieties <sup>3</sup><sup>3</sup>3By abusing the language we call ”quantum varieties” the corresponding ”coordinate rings”. are only relevant in general.
Note, that some one-sided modules which are q-analogues of line bundles on ”quantum generic orbits” were constructed in \[GS\] in the spirit of the Serre-Swan approach. Here we introduce the cotangent and tangent modules on the same orbits. Also, we discuss the problem of an operator meaning of the tangent modules, i.e. that of representing them by ”braided vector fields” (note that in the case of the quantum sphere this problem was solved in \[A\]).
In the sequel we prefer to use the quantum group (QG) $`U_q(sl(n))`$ instead of this $`U_q(su(n))`$ since we disregard any involution in quantum algebras in question. The basic field $`k`$ is assumed to be $`\mathrm{I}\mathrm{R}`$ or $`\mathrm{I}\mathrm{C}`$. Throughout the whole of the paper the parameter $`q`$ is assumed to be generic.
Acknowledgment One of the authors (D.G.) is supported by the grant PICS-608.
## 2 Cotangent $`k(S_q^2)`$-module
First, let us introduce the quantum sphere in a form appropriate for our goals. Let $`V`$ be a three dimensional vector space being a left $`U_q(sl(2))`$-module. The action of the quantum group $`U_q(sl(2))`$ can be extended to any tensor power of $`V`$ via the coproduct. Let us decompose the space $`V^2`$ into a direct sum of irreducible $`U_q(sl(2))`$-modules:
$$V^2=V_0V_1V_2$$
(2.1)
where $`V_i`$ stands for the spin $`i`$ $`U_q(sl(2))`$-module. Let $`v_0`$ be a generator of the one-dimensional component $`V_0`$.
Let us introduce quantum sphere by
$$k(S_q^2)=T(V)/\{V_1,v_0c\},ck,c0$$
(2.2)
where $`T(V)`$ stands for the free tensor algebra of the space $`V`$ and $`\{I\}`$ stands for the ideal generated by a family $`IT(V)`$. This algebra is a particular (in some sense q-commutative) case of the Podles’ sphere introduced in \[P1\]. (We disregard involution operators in this algebra. So, in fact we do not distinguish quantum sphere and quantum hyperboloid.)
Remark, that in definition (2.2) of the quantum sphere we do not use any coordinate form of the quantum sphere. In a similar non-coordinate way we introduce the cotangent module on it.
Let $`V^{}`$ be the space isomorphic to $`V`$ as $`U_q(sl(2))`$-modules but spanned by the differentials $`dx,xV`$. Let $`V^{}k(S_q^2)`$ be the free finitely generated right $`k(S_q^2)`$-module and its submodule $`M_r`$ generated by $`(V^{}V)_0`$. Thus, we have
$$M_r=\mathrm{Im}\mu ^{23}((V^{}V)_0k(S_q^2)).$$
Hereafter $`(V^{}V)_i`$ stands for the spin $`i`$ component in the product $`V^{}V`$, $`\mu `$ stands for the product in the algebra $`k(S_q^2)`$, and $`\mu ^{ij}`$ is the operator $`\mu `$ applied to the i-th and j-th factors. We call a right cotangent module on the quantum sphere the following quotient
$$T_r^{}(S_q^2)=(V^{}k(S_q^2))/M_r.$$
We define a left cotangent module $`T_l^{}(S_q^2)`$ on the quantum sphere in the same way as the quotient of the left $`k(S_q^2)`$-module $`k(S_q^2)V^{}`$ over its submodule
$$M_l=\mu ^{12}(k(S_q^2)(VV^{})_0).$$
(Hereafter we omit the symbol $`\mathrm{Im}`$.)
These modules are quantum analogues (respectively, right and left ones) of the cotangent vector bundle over quantum sphere (or quantum hyperboloid) realized in the spirit of the Serre-Swan approach.
As was shown in \[AG\], the left cotangent module $`T_l^{}(S_q^2)`$ is a flat deformation of its classical counterpart (in fact, the proof consists in showing that the classical and quantum objects are built from the same, respectively, $`U(sl(2))`$\- and $`U_q(sl(2))`$-irreducible components). By the same reason the right cotangent module $`T_r^{}(S_q^2)`$ is a flatly deformed object.
Now, define a two-sided cotangent module (in the sequel called cotangent bimodule) on the quantum sphere. Let us set
$$\overline{T^{}(S_q^2)}=(k(S_q^2)V^{}k(S_q^2))/(M_lk(S_q^2)+k(S_q^2)M_r).$$
This $`k(S_q^2)`$-bimodule is much bigger than one-sided one even in the classical case $`(q=1)`$ because so far we do not have any rule for transposing the ”quantum functions” and differentials (i.e., elements of $`V^{}`$). In what follows we omit the subscription $`q`$ if $`q=1`$.
First, let us consider the classical case in details. In order to reduce this bimodule to the seize of the one-sided one we should define a commutation rule between elements of the algebra and those of $`V^{}`$. In the classical case it is always done by the flip. Namely, we set
$$T^{}(S^2)=\overline{T^{}(S^2)}/\{avva\}ak(S^2),vV^{}.$$
(2.3)
It is not difficult to see that there exists a one-to-one correspondence between the one-sided (say, right) $`k(S^2)`$-module $`T_r^{}(S^2)`$ and the bimodule $`T^{}(S^2)`$. Indeed, modulo the denominator of (2.3) any element of the $`k(S^2)`$-bimodule $`k(S^2)V^{}k(S^2)`$ can be reduced to an element of the right module $`V^{}k(S^2)`$. Thus, we have a map
$$\rho :k(S^2)V^{}k(S^2)V^{}k(S^2).$$
The following inclusion is clear
$$\rho (M_lk(S^2)+k(S^2)M_r)M_r.$$
This implies that the map $`\rho `$ sends the two-sided module $`T^{}(S^2)`$ into $`T_r^{}(S^2)`$. Moreover, it is isomorphism of linear spaces.
An analogous construction for algebras related to an involutary YB operator $`S`$ can be introduced in a similar way. In this case the denominator of formula (2.3) should be replaced by $`\{avS(av)\}`$. However, if we want to realize a similar approach for the algebra $`k(S_q^2)`$ (or for other algebras related to non-involutary YB operators $`S`$) it is not clear what should be a proper analogue of the denominator in (2.3).
Let $`\overline{S}:VV^{}V^{}V`$ be any $`U_q(sl(2))`$-covariant invertible operator (called in the sequel a transposition). Let us replace the denominator in (2.3) by $`\{av\overline{S}(av)\}`$ (we also assume that $`S^2`$ in the numerator is replaced by $`S_q^2`$). The problem consists in finding all transpositions $`\overline{S}`$ such that the corresponding quotient denoted $`T^{}(S_q^2)`$ would be a flat deformation of the quotient (2.3).
It is evident that in order to give rise to a flatly deformed object a transposition $`\overline{S}`$ should preserve the ideal of the formula (2.2) and take the submodule $`M_l`$ into $`M_r`$. Otherwise, by passing to the right $`k(S_q^2)`$-module $`T_r^{}(S^2)`$ we would get some supplementary relations in it what would lead to a collapse of the final object. Essentially, this means that the map $`\rho `$ takes the $`k(S_q^2)`$-bimodule $`\overline{T^{}(S_q^2)}`$ onto some proper quotient of the right $`k(S_q^2)`$-module $`T_r^{}(S^2)`$.
As we will see in the next section the only transpositions preserving the ideal in (2.2) are $`\pm S^{\pm 1}`$ where $`S`$ is the YB operator coming from $`U_q(sl(2))`$. Then we will show that even these operators do not take the submodule $`M_l`$ into $`M_r`$. So, any transposition $`\overline{S}`$ leads to the collapse mentioned above.
## 3 Non-existence of a flat deformation
Now we represent the quantum sphere in a more explicit (coordinate) form. Let us fix the base $`(u,v,w)`$ in the space $`V`$ with the following action of the QG $`U_q(sl(2))`$
$$\begin{array}{ccccccccccc}X.u& =& 0,\hfill & & X.v& =& (q+q^1)u,\hfill & & X.w& =& v,\hfill \\ Y.u& =& v,\hfill & & Y.v& =& (q+q^1)w,\hfill & & Y.w& =& 0,\hfill \\ H.u& =& 2u,\hfill & & H.v& =& 0,\hfill & & H.w& =& 2w.\hfill \end{array}$$
Hereafter $`X,H,Y`$ are the standard generators of the QG $`U_q(sl(2))`$ (cf. \[CP\]).
Note, that the QG $`U_q(sl(2))`$ acts on the space $`V^{}`$ in the same way (we should only replace the generators $`(u,v,w)`$ by $`(du,dv,dw)`$ in the formulae above).
Thus, the spaces $`V_i,i=0,1,2`$ being irreducible $`U_q(sl(2))`$-modules are as follows
$$V_0=\mathrm{span}(v_0),v_0=(q^3+q)uw+v^2+(q+q^1)wu,$$
$$V_1=\mathrm{span}(q^2uvvu,(q^3+q)(uwwu)+(1q^2)v^2,wvq^2vw),$$
$$V_2=\mathrm{span}(u^2,uv+q^2vu,uwqv^2+q^4wu,vw+q^2wv,w^2)$$
(the sign $``$ is systematically omitted).
It is well known that the YB operator $`S`$ being restricted onto each component becomes scalar. Namely,
$$S|_{V_0}=q^4\mathrm{id},S|_{V_1}=q^2\mathrm{id},S|_{V_2}=q^2\mathrm{id}$$
(cf. f.e. \[G\]). This implies that being applied to the product $`VV^{}`$ the operator $`S`$ acts as follows
$$S(udu)=\alpha duu$$
(3.4)
$$S(q^2udvvdu)=\beta (q^2duvdvu)$$
(3.5)
$$S((q^3+q)udw+vdv+(q+q^1)wdu)=\gamma ((q^3+q)duw+dvv+(q+q^1)dwu)$$
(3.6)
with $`\alpha =q^2,\beta =q^2,\gamma =q^4`$ and similarly for other elements of each component.
Now consider an arbitrary invertible transposition $`\overline{S}:VV^{}V^{}V`$ being a $`U_q(sl(2))`$-morphism (we do not require it to be a YB operator). It is given by the same formulae (3.4)-(3.6) (and all their descendants) but with arbitrary non-trivial $`\alpha ,\beta `$ and $`\gamma `$. By applying such an operator many times we can transform any element from $`V^kV^{}`$ into that from $`V^{}V^k`$. Let us note that the operator $`\overline{S}=S`$ has the following remarkable property (this property is also valid for the operator $`S^1`$).
###### Proposition 1
We have
$$S\mu ^{12}=\mu ^{23}S^{12}S^{23}\mathrm{and}S\mu ^{23}=\mu ^{12}S^{23}S^{12}$$
.
Proof The defining relations of the algebra $`k(S_q^2)`$ are coordinated with the action of the QG $`U_q(sl(2))`$ in the following sense
$$X.\mu (ab)=\mu \mathrm{\Delta }(X).(ab)XU_q(sl(2)),a,bk(S_q^2).$$
Let $``$ be the universal quantum R-matrix corresponding to the QG $`U_q(sl(2))`$. It satisfies the relations
$$\mathrm{\Delta }^{12}=^{13}^{23}\text{and}\mathrm{\Delta }^{23}=^{13}^{12}$$
(this means that the QG in question is quasitriangular).
Thus, for any $`a,b,ck(S_q^2)`$ we have (hereafter $`\sigma `$ is the usual flip)
$`S\mu ^{12}(abc)`$ $`=`$ $`\sigma \mu ^{12}(abc)`$
$`=`$ $`\sigma (_1(\mu ^{12}(ab))_2c)`$
$`=`$ $`\sigma (\mu ^{12}\mathrm{\Delta }(_1)(ab)_2c)=\sigma \mu ^{12}(\mathrm{\Delta }^{12})(abc).`$
Here we use the Sweedler’s notation and apply the components of $``$ to the elements $`a,b,c`$ w.r.t. the action $`U_q(sl(2))`$ on the algebra $`k(S_q^2)`$. Moreover, we use the relation $`S=\sigma `$.
The following chain of identities completes the proof of the first relation of proposition (the second one can be verified in a similar way)
$`\sigma \mu ^{12}\mathrm{\Delta }^{12}`$ $`=`$ $`\mu ^{23}\sigma ^{12}\sigma ^{23}\mathrm{\Delta }^{12}`$
$`=`$ $`\mu ^{23}\sigma ^{12}\sigma ^{23}^{13}^{23}=\mu ^{23}\sigma ^{12}\sigma ^{23}^{13}\sigma ^{23}\sigma ^{23}^{23}`$
$`=`$ $`\mu ^{23}\sigma ^{12}\sigma ^{23}^{13}\sigma ^{23}S^{23}=\mu ^{23}\sigma ^{12}^{12}S^{23}`$
$`=`$ $`\mu ^{23}S^{12}S^{23}.`$
It is evident that the above proposition is still valid if we replace the operator $`S`$ by $`S^1`$ in the formulae above. However, if we replace the operator $`S`$ by any other transposition $`\overline{S}`$ the identities from this proposition become broken in virtue of the following.
###### Proposition 2
The only operators $`\overline{S}`$ such that
$$\overline{S}(v_0v^{})=v^{}v_0,v^{}V^{}\mathrm{and}\overline{S}(V_1V^{})V^{}V_1$$
are $`\pm S`$ and $`\pm S^1`$.
Proof In the sequel we represent an arbitrary transposition as
$$\overline{S}=xP_0+yP_1+zP_2,x,y,zk$$
where the operators
$$P_i:VV^{}(V^{}V)_i,i\{0,1,2\}$$
(3.7)
become the projectors $`V^2V_i`$ if we identify $`V`$ and $`V^{}`$. For $`\overline{S}=S`$ we have $`x=\gamma =q^4,`$ $`y=\beta =q^2,`$ $`z=\alpha =q^2.`$ In what follows we need the images of some elements under the action of the operators $`P_i`$:
$`P_0(vdv)`$ $`=`$ $`\alpha _1d\overline{v}_0,P_0(udw)=\alpha _1^{}d\overline{v}_0,P_2(udu)=duu,`$
$`P_2(udv)`$ $`=`$ $`\beta ^{}(duv+q^2dvu),P_2(vdu)=\overline{\beta }^{}(duv+q^2dvu),`$
$`P_2(vdv)`$ $`=`$ $`\gamma _1(duwqdvv+q^4dwu),`$
$`P_1(vdv)`$ $`=`$ $`\beta _1[(q^3+q)(duwdwu)+(1q^2)dvv],`$
$`P_1(vdu)`$ $`=`$ $`\overline{\alpha }^{}(q^2duvdvu),P_1(udv)=\alpha ^{}(q^2duvdvu),`$
$`dv_0`$ $`=`$ $`(q^3+q)udw+vdv+(q+q^1)wdu,d\overline{v}_0=(q^3+q)duw+`$
$`dvv+(q+q^1)dwu,`$
$`P_1^{12}(ud\overline{v}_0)`$ $`=`$ $`\alpha ^{}(q^2duvdvu)v+\beta _1^{}(q+q^1)[(q^3+q)(duwdwu)+`$
$`(1q^2)dvv]u,`$
$`P_2^{12}(ud\overline{v}_0)`$ $`=`$ $`(q^3+q)duuw+\beta ^{}(duv+q^2dvu)v+\gamma _1^{}(q+q^1)(duw`$
$`qdvv+q^4dwu)u,`$
$`P_0^{12}(ud\overline{v}_0)`$ $`=`$ $`(q+q^1)\alpha _1^{}d\overline{v}_0u,`$
where
$`\beta ^{}`$ $`=`$ $`(1+q^4)^1,\alpha ^{}=q^2\beta ^{},\overline{\alpha }^{}=\beta ^{},\overline{\beta }^{}=q^2\beta ^{},\gamma _1^{\prime \prime }=q^4\gamma _1^{},`$
$`\beta _2`$ $`=`$ $`\beta ^{},\alpha _2=\alpha ^{},\alpha _2^{}=\beta ^{},\beta _2^{}=q^2\alpha _2^{},\alpha _1^{\prime \prime }=q^2\alpha _1^{},`$
$`\alpha _1`$ $`=`$ $`q^2(1+q^2+q^4)^1,\beta _1=(1q^2)\beta ^{},\beta _1^{\prime \prime }=\beta _1^{},`$
$`\gamma _1`$ $`=`$ $`q(1+q^2)^2(1+q^2+q^4)^1\beta ^{},\gamma _1^{}=(1+q^2+q^4)^1\beta ^{},`$
$`\alpha _1^{}`$ $`=`$ $`q^3(1+q^2)^1(1+q^2+q^4)^1,\beta _1^{}=q(1+q^2)^1\beta ^{}.`$
On applying the transposition $`\overline{S}`$ to the element $`(q^2uvvu)du`$ we get
$`\overline{S}((q^2uvvu)du)`$ $`=`$ $`\overline{S}^{12}\overline{S}^{23}((q^2uvvu)du)`$
$`=`$ $`q^2\overline{S}^{12}\overline{S}^{23}(uvdu)\overline{S}^{12}\overline{S}^{23}(vudu).`$
By a straightforward but tedious computations with the use of the formulae above we get the following result for the coefficient at the element $`dvuu`$ in the image above
$$q^4\beta ^2(\alpha ^2+(q^4+q^4)\alpha \beta +\beta ^2).$$
(3.8)
The condition $`\overline{S}(V_1V^{})(V^{}V_1)`$ implies
$$\alpha ^2+(q^4+q^4)\alpha \beta +\beta ^2=0.$$
(3.9)
This equation has two solutions (in the projective sense)
$$\alpha =q^4\beta \mathrm{and}\alpha =q^4\beta .$$
(3.10)
Let us remark that the first (resp., second) solution is satisfied by the operator $`cS`$ (resp. $`cS^1`$) with an arbitrary $`c0.`$ Then, the operator $`\overline{S}`$ in general can be represented as follows
$$\overline{S}=cS^{\pm 1}+\delta P_0,\delta k.$$
(3.11)
Now it remains to show that $`\delta =0`$ and $`c=\pm 1`$. Let us do it for $`S`$ (the $`S^1`$ case is analogous). From the above form of $`\overline{S}`$ we have
$$\overline{S}^{12}\overline{S}^{23}=c^2S^{12}S^{23}+c\delta (S^{12}P_0^{23}+P_0^{12}S^{23})+\delta ^2P_0^{12}P_0^{23}.$$
(3.12)
Let us consider the image of the element $`(q^2uvvu)dv`$ w.r.t. the action of the transposition $`\overline{S}`$. In this image we are interested in terms containing $`duv_0`$ or $`dwuu`$. Let us denote $`I_1`$ (resp. $`I_2`$) the coefficient at $`duv_0`$ (resp. $`dwuu`$) in $`\overline{S}^{12}\overline{S}^{23}((q^2uvvu)dv).`$ A straightforward computation shows that
$$\begin{array}{c}I_1=(q^3+q)\alpha _1^{\prime \prime }I_0+(q+q^1)\alpha _1\alpha _1^{}[q^6(q+q^1)\beta ^{}+q^2\gamma _1^{}]c\delta ,\\ I_2=(q+q^1)I_0+(q^3+q)\alpha _1[\beta ^{}+q^6\gamma _1^{}]c\delta \mathrm{where}\\ I_0=(q^3+q)\alpha _1\alpha _1^{}\delta ^2+\alpha _1[2q^3(1+q^2)\alpha _1^{}+(\beta _1q^8(1+q^2)\gamma _1^{})1]c\delta .\end{array}$$
(3.13)
Since the element $`duv_0`$ is that of the highest weight and the element $`(q^2uvvu)dv`$ is not, the coefficient $`I_1`$ is trivial. Moreover, the condition $`\overline{S}(V_1V^{})(V^{}V_1)`$ implies that the coefficient $`I_2`$ is trivial. These two relations imply $`c\delta =0.`$
Since $`c0`$ (unless the operator $`\overline{S}`$ is not invertible) we have $`\delta =0`$.
Thus, if we admit the first condition of (3.10) the only transposition $`\overline{S}=cS`$ could preserve the defining ideal of $`k(S_q^2)`$. But in fact only factors $`c=\pm 1`$ are compatible with the centrality of the element $`v_0`$. This completes the proof.
Now we pass to showing that even the operator $`S`$ does not preserve the flatness of the deformation (for the operators $`S`$ and $`\pm S^1`$ the proof is analogous).
###### Proposition 3
The image of $`M_l`$ w.r.t. the YB operator $`S`$ does not belong to $`M_r.`$
Proof By definition of the submodules $`M_l`$ and $`M_r`$, they consist respectively of the following elements
$$\mu ^{12}(fdv_0)\mathrm{and}\mu ^{23}(d\overline{v}_0f),fk(S_q^2).$$
(3.14)
Let us show that there exists an element $`fk(S_q^2)`$ such that $`S(\mu ^{12}(fdv_0))`$ $`M_r`$. Let $`f=u`$. We have
$`S(\mu ^{12}(udv_0))`$ $`=`$ $`\mu ^{23}S^{12}S^{23}(udv_0)`$
$`=`$ $`\mu ^{23}(S^{23})^1S^{23}S^{12}S^{23}(udv_0)`$
$`=`$ $`\mu ^{23}(S^{23})^1(d\overline{v}_0u).`$
By using
$$(S^{23})^1(d\overline{v}_0.u)=\gamma ^1P_0^{23}(d\overline{v}_0u)+\beta ^1P_1^{23}(d\overline{v}_0u)+\alpha ^1P_2^{23}(d\overline{v}_0u),$$
(3.15)
we get
$$\mu ^{23}(S^{23})^1(d\overline{v}_0u)=\gamma ^1\mu ^{23}(P_0^{23}(d\overline{v}_0u))+\alpha ^1\mu ^{23}(P_2^{23}(d\overline{v}_0u)).$$
(3.16)
We state that there is no element $`gk(S_q^2)`$ such that $`\mu ^{23}(d\overline{v}_0g)`$ would be equal to the r.h.s. of (3.16). Indeed, it could be only an element of the form $`g=\nu u,\nu k`$. However, since $`\alpha \gamma `$ and the both components in (3.16) are non-trivial we conclude that no appropriate factor $`\nu `$ exists. This completes the proof.
###### Remark 1
The statement of this proposition can be generalized to other quantum algebras related to non-involutary YB operators. The crucial property of the operator $`S`$ used in the proof is the following one: this operator has more than one distinct eigenvalues and the corresponding components do not vanish in the algebra $`k(S_q^2)`$.
However, for certain algebras and certain their two-sided modules the flatness of deformation is valid. Let us consider for instance, the quantum cone (it corresponds to the case $`c=0`$). The module $`T^{}(S_q^2)`$ defined as above is a flat deformation of its classical counterpart. This follows from the fact that in the corresponding quantum coordinate ring defined by $`V_0=0,V_1=0`$ the only component $`V_2V^2`$ survives. This prevents us from the effect used in the proof above.
For the same reason, in quantum geometry dealing with non-involutary YB operators $`S`$ it is not convenient to use this operator (either any other transposition) in order to define a product $`M_1_AM_2`$ of two one-sided $`A`$-modules assuming $`A`$ to be a quantum algebra looking like that $`k(S_q^2)`$. However, the product of two modules can be apparently defined as the quantum deformation of the product of their classical counterparts. So, the notation $`M_1_AM_2`$ must be regarded in this restricted sense without any transposition of the elements of $`A`$ and those of $`M_1`$ (or $`M_2`$).
## 4 Generic quantum orbits and modules on them
The main purpose of this section is to generalize the construction of the cotangent module on quantum sphere to some other quantum orbits. Also, we define the tangent modules on these orbits and discuss the problem of equipping the tangent module with an action on the quantum coordinate ring in question. All constructions are done in the framework of one-sided modules over algebras in question. This allows us to hope that these modules are flatly deformed objects.
First of all, describe quantum orbits in question. Let us begin with evoking their quasiclassical counterparts (i.e. the corresponding Poisson structures).
As was shown in \[DGS\], on any orbit (of a semisimple element) $`𝒪𝐠^{}`$ where $`𝐠`$ is a simple Lie algebra there exists a family of Poisson-Lie structures (for the compact form of the Lie algebra this family is labeled by the elements of $`H^2(𝒪)`$). Morover, in this family there exists a bracket which is compatible with the Kirillov-Kostant-Souriau one. A quantization of this particular Poisson bracket can be realized (at least in the $`𝐠=sl(n)`$ case) in terms of the so-called reflection equation (RE) algebra. The resulting algebra can be described as an appropriate quotient of the RE algebra. Thus, we get an explicit realization of such an algebra in the spirit of algebraic geometry by means of a system of braided algebraic equations. (As for other Poisson-Lie structures they can be quantized by means of formal series in the spirit of deformation quantization. Their description in terms of so-called Hopf-Galois extension is also known, cf. \[Sh\], \[HM\].)
Let $`𝐠=sl(n)`$ and $`𝐠_q`$ be the same as vector space but equipped with a $`U_q(sl(n))`$-action which is a deformation of the adjoint action of $`𝐠`$ onto itself. Let us extend this action to the space $`𝐠_q^2`$ by means of the coproduct in $`U_q(sl(n))`$ and decompose it into a direct sum
$$𝐠_q^2=I_+I_{}$$
(4.17)
of two $`U_q(sl(n))`$-invariant subspaces $`I_+`$ and $`I_{}`$ so that the corresponding algebras
$$_\pm =_\pm (𝐠_q)=T(𝐠_q)/\{I_{}\}$$
would be flat deformations of the symmetric $`_+(𝐠)`$ and skewsymmetric $`_{}(𝐠)`$ algebras respectively. Since the space $`𝐠_q^2`$ is not multiplicity free (the component isomorphic to $`𝐠_q`$ itself comes twice in the decomposition of $`𝐠_q^2`$ into a direct sum of irreducible $`U_q(sl(n))`$-modules) it is not evident that decomposition (4.17) exists.
Nevertheless, it does exist and can be constructed by means of the RE algebra mentioned above and of a $`U_q(sl(n))`$-covariant pairing. Let us recall that by the RE algebra one means the algebra generated by $`n^2`$ elements $`l_i^j,\mathrm{\hspace{0.17em}\hspace{0.17em}1}i,jn`$ subject to the relations
$$SL_1SL_1L_1SL_1S=0$$
(4.18)
where $`L=(l_i^j)`$ is the matrix with the entries $`l_i^j`$ and $`L_1=L\mathrm{id}`$.
It is known that this algebra has the center generated by the elements $`C_q^p=\mathrm{tr}_qL^p,p=1,\mathrm{},n`$ where $`\mathrm{tr}_q`$ is the quantum analogue of the usual trace. Then the space $`\mathrm{span}(l_i^j)`$ is a direct sum of a one-dimensional $`U_q(sl(n))`$-module generated by $`\mathrm{tr}_qL`$ and a $`n^21`$-dimensional one which we identify with $`𝐠_q`$. Then the space $`I_{}`$ can be treated as the l.h.s. of (4.18) modulo the elements of the form
$$l\mathrm{tr}_qL,\mathrm{tr}_qLl,l\mathrm{span}(l_i^j).$$
(4.19)
In virtue of \[L\] the algebra $`_+`$ is a flat deformation of its classical counterpart. (Let us note that the RE algebra and adjacent objects are also well defined for non-quasiclassical YB operators, cf. \[GPS\]).
We introduce the space $`I_+`$ as that orthogonal to $`I_{}`$ w.r.t. the pairing
$$(,):𝐠_q^2𝐠_q^2k,(,)=<,><,>^{23}$$
where $`<,>:𝐠_q^2k`$ is a (unique up to a factor) $`U_q(sl(n))`$-covariant pairing. Then following \[D\] we can state that the algebra $`_{}`$ is a flatly deformed object as well.
Now, let us introduce ”the generic quantum orbits” by the following system of equations
$$C_q^pc^p=0,c^pk,p=1,\mathrm{},n.$$
(4.20)
The constant $`c^1`$ is equal to 0 while the other constants are assumed to be generic. Let us note $`k(M_q)`$ the quotient of the RE algebra over the ideal $`\{J\}`$ generated by the l.h.s. elements of (4.20). This ”quantum coordinate ring” is a flat deformation of coordinate ring of a generic orbit in $`sl(n)^{}`$.
As for q-analogues of other orbits of semisimple elements in $`sl(n)^{}`$ the reader is referred to \[DGK\] where the case of the ”$`\mathrm{I}\mathrm{CI}\mathrm{P}^n`$ type orbits” was studied.
Now let us introduce quantum analogues of the cotangent module and its exterior powers on the orbits in question. (In the sequel all the modules are left.)
Consider the elements $`dC_q^p`$ looking like $`dv_0`$ of the previous section. This means that the differential $`d`$ is applied only to the last factor of the element $`C_q^p`$. Let us multiply the elements $`dC_q^p`$ by those of $`k(M_q)`$ from the left and the elements of $`_{}`$ (in the sense of the algebra $`_{}`$) from the right.
Now, consider the quotient of the left $`k(M_q)`$-module $`k(M_q)_{}^l`$ over its submodule formed by the elements
$$f_pdC_q^pg_p,f_pk(M_q),g_p_{}^{l1}$$
(hereafter $`_{}^l`$ is the degree $`l`$ homogeneous component of the algebra $`_{}`$). Conjecturally, this quotient is a flat deformation of the space of degree $`l`$ differential forms. A proof of this conjecture for the quantum sphere is given in \[AG\]. Also suggested in that paper was a $`U_q(sl(n))`$-covariant de Rham type complex which was a deformation of its classical counterpart without making use of the Leibniz rule.
However, the orbits in question are not multiplicity free any more and the problem of constructing a $`U_q(sl(n))`$-covariant complex which would be a flat deformation of the usual de Rham one becomes more delicate. Nevertheless, if we are only interested in q-analogues of 2-forms being generators of the cohomology ring on quantum orbits in question we can explicitly construct them in the following way (by analogy with the classical case).
It suffices to apply the q-cobracket to the last factor of each element $`C_q^p`$ and treat its image as an element of $`_{}^2`$ (i.e. by realizing it as a sum of the summands $`dx_idx_j`$). By definition, the q-cobracket is the inverse (in some natural sense) of the q-Lie bracket whose construction was featured in \[LS\].
Let us call the above quotient module corresponding to the case $`l=1`$ as the cotangent module and denote it as $`T^{}(M_q)`$. Emphasize once more that this module is introduced as a one-sided module (namely, the left one but in the same way we can introduce the right one). Moreover, it is introduced explicitly by a system of equations. In a similar way we can realize the other modules defined above.
Now let us pass to defining the tangent module $`T(M_q)`$ on the orbits in question. In the classical case the tangent module $`T(M)`$ on a given regular affine variety $`M`$ has an operator realization by vector fields, i.e., there exists a map
$$T(M)k(M)k(M)$$
(4.21)
which commutes with the $`k(M)`$-module structure product
$$k(M)T(M)T(M).$$
Moreover, if $`M`$ is an orbit in $`𝐠^{}`$ there exists an embedding
$$𝐠T(M)$$
such that map (4.21) realizes a representation of Lie algebra $`𝐠`$ by vector fields in the coordinate ring $`k(M)`$.
As for the tangent modules on quantum orbits $`k(M_q)`$ we define them by the same system as the cotangent ones. This is motivated by the fact that in the classical case the tangent and cotangent modules on orbits in question are isomorphic. So, conjecturally the deformation of the tangent module is flat.
However, in the quantum case there exists the problem of an operator meaning of the tangent module. For the case of the quantum sphere this problem was solved in \[A\]. Namely, it was shown that for the tangent quantum module $`T(S_q^2)`$ there exists a map
$$T(S_q^2)k(S_q^2)k(S_q^2)$$
(4.22)
commuting with the module structure product
$$k(S_q^2)T(S_q^2)T(S_q^2).$$
Also in \[A\] an embedding was constructed of the form
$$sl(2)_qT(S_q^2)$$
such that map (4.22) realized a representation of the q-Lie algebra $`sl(2)_q`$ (this means that the relations between the generators of $`sl(2)_q`$ in its enveloping algebra are preserved under map (4.22)).
We call braided vector fields the elements of $`T(S_q^2)`$ realized as operators on $`k(S_q^2)`$ via the map (4.22). However the problem of a similar treatement of the tangent modules on the quantum orbits in question is still open.
Let us complete the paper with the following remark. There exists a lot of articles devoted to different aspects of ”braided geometry”. However, often they do not consider the problem of flatness of quantum deformation. Nevertheless, a flat deformation is rather subtle phenomenon. Even if the flatness is fulfilled for a deformation of complexes related to a vector space, in general it disappears if one tries to restrict the differential algebras to a ”quantum variety”. We are sure that the approach making use of one-sided modules on quantum varieties developed here (as well as in \[AG\], \[GS\]) is more adequate for the needs of ”braided geometry” on quantum varieties since conjecturally it allows us to preserve the flatness of deformation (at the expense of the Leibniz rule).
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# Photoproduction of top quarks in peripheral heavy ion collisions
## Abstract
In relativistic heavy ion collisions, top quarks can be produced by photon-gluon fusion when a photon from the Weizsäcker-Williams virtual photon field of one nucleus interacts with a gluon in the other nucleus. Photoproduction with heavy ions at the Large Hadron Collider (LHC) will be the first accessible non-hadronic top production channel. We calculate the $`t\overline{t}`$ photoproduction cross sections, pair mass and top quark rapidity distributions in peripheral lead-lead and oxygen-oxygen collisions. The cross sections are sensitive to the top quark charge and the large-$`Q^2`$ gluon distribution in the nucleus. We find a cross section of 15 pb in oxygen-oxygen collisions, leading to 210 pairs in a one month ($`10^6`$ s) LHC run. In $`pA`$ collisions, the rate is higher, 1100 pairs per month for $`p`$O. A comparison of the $`AA`$ and $`pA`$ data might allow for a study of gluon shadowing at high $`Q^2`$.
preprint: LBNL-45743
Due to their large charges, relativistic heavy ions carry strong electromagnetic fields which may be treated as virtual photon fields. In a relativistic ion collider, these fields interact with target nuclei in the opposing beam, resulting in high luminosities for photonuclear interactions. Because of the large Lorentz boosts, high photon-nucleon center of mass energies are reached. Previous studies have considered photoproduction of charm and bottom quarks as well as nuclear breakup and vector meson production . These calculations all considered peripheral collisions, with impact parameter, $`b`$, greater than twice the nuclear radius $`R_A`$ so that the two nuclei do not interact hadronically. Peripheral collisions will be studied experimentally at RHIC and the LHC, now under construction at CERN.
Here, we consider the production of top quarks via photon-gluon fusion, paralleling previous calculations of photoproduction in heavy ion collisions . We calculate top production from the heaviest, Pb, and lightest, O, ions planned for the LHC. In these collisions, a $`t\overline{t}`$ is produced in the reaction $`\gamma (k)+A(P)t(p_1)+\overline{t}(p_2)+X`$ where $`k`$ is the four momentum of the photon emitted from the virtual photon field of the projectile nucleus, $`P`$ is the four momentum of the interacting nucleon in target nucleus $`A`$, and $`p_1`$ and $`p_2`$ are the four momenta of the produced $`t`$ and $`\overline{t}`$ quarks. Note that a phton from the target can also interact with a nucleon in the projectile. We work in the center of mass (lab) frame. The photons are almost real, with $`|q^2|<(\mathrm{}c/R_A)^2`$. The slight virtuality is neglected.
On the parton level, the photon-gluon fusion reaction is $`\gamma (k)+g(x_2P)t(p_1)+\overline{t}(p_2)`$ where $`x_2`$ is the fraction of the target momentum carried by the gluon. To lowest order (LO), the $`t\overline{t}`$ production cross section is
$$s^2\frac{d^2\sigma }{dt_1du_1}=\pi \alpha _s(Q^2)\alpha e_t^2\left(\frac{t_1}{u_1}+\frac{u_1}{t_1}+\frac{4m^2s}{t_1u_1}\left[1\frac{m^2s}{t_1u_1}\right]\right)\delta (s+t_1+u_1)$$
(1)
where the partonic invariants $`s`$, $`t_1`$, and $`u_1`$ are defined as $`s=(k+x_2P)^2`$, $`t_1=(kp_1)^2m^2=(x_2Pp_2)^2m^2`$, and $`u_1=(x_2Pp_1)^2m^2=(kp_2)^2m^2`$ with top quark mass $`m=175`$ GeV. In this case, $`s=4k\gamma x_2m_p`$ where $`\gamma `$ is the Lorentz boost of a single beam and $`m_p`$ is the proton mass. Here, $`e_t=2/3`$ is the expected top quark charge, $`\alpha =e^2/\mathrm{}c`$ is the electromagnetic coupling constant, and $`\alpha _s(Q^2)0.11`$ is the one-loop strong coupling constant evaluated at scale $`Q^2=m^2+p_T^2`$ where $`p_T`$ is the transverse momentum of the produced top quark. The large mass of the top quark prevents toponium production, allowing the $`t`$ and $`\overline{t}`$ to be treated as free quarks, even near threshold .
The hadronic top production cross section is obtained by integrating Eq. (1) over $`x_2`$ and $`k`$, since virtual photons are emitted by the nucleus in a continua of four-momenta,
$$S^2\frac{d^2\sigma _{\gamma At\overline{t}X}}{dT_1dU_1}=2_{k_{\mathrm{min}}}^{\mathrm{}}𝑑k\frac{dN}{dk}_{x_{2_{\mathrm{min}}}}^1\frac{dx_2}{x_2}g(x_2,Q^2)s^2\frac{d^2\sigma }{dt_1du_1},$$
(2)
where $`dN/dk`$ is the photon flux. The factor of two in Eq. (2) arises because both nuclei emit photons and thus serve as targets. The incoherence of top production guarantees that there is no interference between the two production sources. The hadronic invariants can be defined for a given photon four momuntum $`k`$. If the top quark is detected, the invariants are $`S=(k+P)^2`$, $`T_1=(Pp_1)^2m^2`$, and $`U_1=(kp_1)^2m^2`$ . The partonic and hadronic invariants are related by $`s=x_2S`$, $`t_1=U_1`$, and $`u_1=x_2T_1`$. Four-momentum conservation at leading order gives $`x_{2_{\mathrm{min}}}=U_1/(S+T_1)`$. In addition to the total cross sections, we also present the top quark rapidity distribution and the $`t\overline{t}`$ invariant mass distribution. We define the $`t`$ rapidity as $`y`$ and the $`\overline{t}`$ rapidity as $`y_2`$. The top quark rapidity is related to the invariant $`T_1`$ by $`T_1=\sqrt{S}Qe^y`$. The invariant mass of the pair can be determined if both the $`t`$ and $`\overline{t}`$ are detected. The invariant mass is $`s=M^2=2Q^2(1+\mathrm{cosh}(yy_2))`$. The minimum photon momentum necessary to produce a $`t\overline{t}`$ pair is $`k_{\mathrm{min}}=M^2/(4\gamma x_2m_p)`$.
We use the MRST LO gluon distribution $`xg(x,Q^2)`$ which is considerably softer than the older, flat parameterizations such as $`xg(x)(1x)^n`$ used in earlier photoproduction predictions. Calculations with $`n=5`$ result in cross sections $`23`$ times higher than those with the MRST LO gluon distribution. Shadowing, the modification of the gluon distribution in nuclei, is included via a parameterization based on a fit to data. There are two caveats regarding shadowing. First, the large $`Q^2`$ of top production is outside the upper limit of the parameterization, $`Q_{\mathrm{max}}^2=10^4`$ GeV<sup>2</sup>. However, the shadowing parameterization continues to evolve beyond this point with changes on the percent level between $`Q^2=10^4`$ GeV<sup>2</sup> and $`6.4\times 10^5`$ GeV<sup>2</sup>. Secondly, the probable impact parameter dependence of the shadowing has been neglected.
The photon flux is given by the Weizsäcker-Williams formulae. The flux is a function of the distance from the nucleus, $`r`$,
$$\frac{d^3N}{dkd^2r}=\frac{Z^2\alpha w^2}{\pi ^2kr^2}\left[K_1^2(w)+\frac{1}{\gamma ^2}K_0^2(w)\right],$$
(3)
where $`w=kr/\gamma `$ and $`K_0(w)`$ and $`K_1(w)`$ are modified Bessel functions. The photon flux is cut off at an energy determined by the size of the nucleus. In the rest frame of the target nucleus, the cutoff is boosted to $`(2\gamma ^21)\mathrm{}c/R_A`$, or 500 TeV for lead and 1800 TeV for oxygen. The $`t\overline{t}`$ production cross section, $`\gamma pt\overline{t}`$ , for a photon with the cutoff energy of oxygen is a factor of 7.2 times larger than the cross section at the cutoff energy of lead. Thus the difference in the energy cutoffs of lead and oxygen is significant.
The total photon flux striking the target nucleus is the integral of Eq. (3) over the transverse area of the target at all impact parameters subject to the constraint that the two nuclei do not interact hadronically. The numerical result agrees to within 15% of the analytical result,
$$\frac{dN}{dk}=\frac{2Z^2\alpha }{\pi k}\left[w_RK_0(w_R)K_1(w_R)\frac{w_R^2}{2}\left(K_1^2(w_R)K_0^2(w_R)\right)\right],$$
(4)
the integral of Eq. (3) for $`r>2R_A2.4A^{1/3}`$ where $`w_R=2kR_A/\gamma `$. The 15% difference can serve as a conservative estimate of the uncertainty on the photon flux; this can be checked using other photoproduction and two-photon interactions.
Although our $`t\overline{t}`$ calculation is at leading order, the large mass of the top quark ensures faster convergence of the perturbative expansion than for the lighter charm and bottom quarks where the next-to-leading order corrections lead to factors of $`23`$ enhancements over the LO cross section. The next-to-leading order $`t\overline{t}`$ cross sections are only $`40`$% larger than the LO cross sections in $`pp`$ interactions . We assume the NLO enhancement to be similar for photoproduction.
The expected center of mass energies and average luminosities for O+O and Pb+Pb collisions at the LHC are given in Table I. The numbers are taken from the latest LHC machine studies on LHC ion operation . The calculation of the average luminosities assumes two experiments, e.g. ALICE and CMS, and a bunch spacing of 125 ns.
The total cross sections for $`t\overline{t}`$ pair production with lead and oxygen beams are 550 pb and 15 pb. As Table I shows, a $`10^6`$ s (one month) heavy ion run at design luminosity will produce 0.2 and 210 pairs respectively. The most significant factor in the rate difference comes from the luminosity, $`3\times 10^4`$ times higher for oxygen.
Shadowing plays little role in the total cross sections since at $`y=0`$ and $`p_T=0`$, $`x_20.06`$ for lead and 0.05 for oxygen. These values are in a region where shadowing effects are small. In addition, at the large $`Q^2`$ required for top production, shadowing effects are further reduced by $`Q^2`$ evolution . The uncertainty in the large-$`Q^2`$ shadowing is important since no data are available. However, because of the large luminosity gain, lighter mass ions are better for studies of nuclear gluon distributions, reducing any possible effect still further.
The top quark rapidity distributions, obtained from Eq. (2), are shown in Fig. 1. The calculations assume that the photon is in the field of the nucleus coming from positive rapidity so that $`y<0`$ corresponds to $`k<\gamma x_2m_p`$ in the center of mass (lab) frame. If the photon is emitted by the target instead of the projectile, the resulting top quark rapidity distribution would then be the mirror image of the distribution in Fig. 1 around $`y=0`$. (Note that this mirror distribution is equivalent to detecting the $`\overline{t}`$ instead of the $`t`$. Thus the $`t`$ and $`\overline{t}`$ distributions are not symmetric around $`y=0`$.) The total top quark rapidity distribution is the sum of the curve in Fig. 1 with its mirror image when both nuclei emit photons. Roughly half the production is within $`|y|<1`$, in the central ALICE acceptance and almost all $`t\overline{t}`$ production falls within the CMS acceptance, $`|y|<2.4`$. The $`t\overline{t}`$ pair invariant mass distributions, $`d\sigma /dM`$, are shown in Fig. 2. The larger photon energy of oxygen results in a broader pair mass distribution.
The only previous study of $`t\overline{t}`$ production via photon-gluon fusion in heavy ion collisions used a very different photon flux. The flux was integrated over all $`r`$, including the nuclear interior, modelling the nucleus as a homogenously charged sphere. This approach is flawed because the photon flux inside the nucleus is poorly defined. Ref. shows that the photon flux varies by orders of magnitude inside the nucleus, depending on the chosen nuclear form factor. More importantly, the calculation is incorrect because the photon flux inside the nucleus is much higher than that outside the nucleus. The authors apply a correction factor of $`1/2`$ but this factor is inappropriately large. Their photon flux (and corresponding cross section) for Pb+Pb collisions is a factor of 10 greater than ours.
Top quark pairs can be observed via their decays, predominantly $`t\overline{t}W^+bW^{}\overline{b}`$, where the $`W`$ decays to $`\mathrm{}\nu `$ or $`q\overline{q}^{}`$. The major background to these channels is likely to be hadroproduction of top in grazing peripheral collisons with slightly smaller impact parameter. The two reactions can be separated based on the presence of rapidity gaps in the collision, the breakup of the colliding nuclei, and a small multiplicity difference. Because the photon is colorless, even if the nucleus breaks up, photoproduction events should also have a rapidity gap between the photon-emitting nucleus and the $`t\overline{t}`$ system. The average multiplicity in $`pp`$ collisions at the LHC is expected to be about 45, slightly higher than the roughly 35 particles expected in a $`\gamma p`$ collision at typical top-production energies. The photoproduction multiplicity should be further reduced since much of the photon energy is needed to produce the $`t\overline{t}`$ pair, leaving less energy for produce additional particle production.
For lead, nuclear breakup can be a complication. A lead nucleus may be excited into a giant dipole resonance with probability $`35\%(b/2R_A)^2`$. When the resonance decays, the nucleus will emit one or more neutrons. To a good approximation, the excitation and the photon-gluon fusion occur independently. Thus the two processes factorize and can be studied separately . For oxygen, the excitation probability is very small, so one of the interacting ions will almost always remain intact.
The photoproduction cross sections are large compared with the corresponding $`ppt\overline{t}X`$ cross sections. At lowest order, $`\sigma (ppt\overline{t}X)`$ = 41 pb at $`\sqrt{s}=5.5`$ TeV, 7% of the lead photoproduction cross section, while $`\sigma (ppt\overline{t}X)`$ = 80 pb at $`\sqrt{s}=7`$ TeV, 5.3 times the corresponding oxygen photoproduction cross section. Thus, only a moderate hadronic rejection factor is required.
Other backgrounds should be small. Hadronic single or double diffractive production without accompanying colored interactions, $`AAAA\overline{t}tX`$ occurs only in a narrow range of impact parameters. Diffractive production is also suppressed by the $`1/M^4`$ final-state mass dependence. Any single diffractive top production will be at larger rapidities than photoproduction which is more central. We have calculated $`\sigma (\gamma \gamma t\overline{t})`$ and found it to be negligible. Backgrounds from other photoproduction channels should be significantly smaller than at hadron colliders because $`\sigma (\gamma pQ\overline{Q}X)/\sigma (\gamma pX)`$ is much larger than $`\sigma (ppQ\overline{Q}X)/\sigma (ppX)`$ at comparable energies.
We believe that these criteria should allow for at least statistical separation of events containing photoproduction of top in oxygen on oxygen collisions.
The top photoproduction cross sections are also measurable in $`pA`$ collisions, as may be possible at the LHC. Because the proton and the ion must have the same magnetic rigidity and $`Z=A`$ for the proton, these collisions are at somewhat higher per nucleon energies than the corresponding $`AA`$ collisions. The cost is that the center of mass is no longer at rest in the lab. Our $`pA`$ results are calculated in the equal speed system so that the $`\gamma `$ is that of the equal speed system. In this case, the photon flux is calculated using the analytical expression in Eq. (4) with $`w_R=(r_p+R_A)k/\gamma `$ and $`r_p=0.6`$ fm is the proton radius.
The expected nucleon-nucleon center-of-mass energies and luminosities for $`pA`$ collisions are shown in Table II. No official LHC $`pA`$ luminosities are yet available . The values in Table II are based on estimates in Ref. , which were obtained assuming a 125 ns bunch spacing and one experiment.
The $`pA`$ collisions would allow a measurement of the gluon structure function in free protons. A comparison between photoproduction in $`pA`$ and $`AA`$ collisions at the same energy would provide a straightforward measurement of nuclear gluon shadowing at $`Q^2`$ values far above those currently available. The rates are much higher at the maximum $`pA`$ energies, however, since the cross sections per nucleon are larger than in $`AA`$ collisions. The center-of-mass energy is 40-60% higher, leading to a larger boost in the target rest frame, increasing the photon cutoff energy to 1200 TeV for lead, more than a factor of two greater than in Pb+Pb, and 3600 TeV for oxygen. In addition, the average photon flux is higher in $`pA`$ because single protons can approach the photon-emitting nucleus more closely than a proton in the center of another nucleus. The estimated luminosities for $`pA`$ interactions lead to 40 pairs in $`p`$Pb and 1100 pairs in $`p`$O over a $`10^6`$ s run, as shown in Table II.
Although the top quarks are produced inside the target nucleus, because of their very high boost with respect to the target, they decay well outside the nucleus. For a standard model top width of $`\mathrm{\Gamma }_t=`$1.5 GeV the top quarks typically travel over 100 fm in the nuclear rest frame before decaying. Thus the $`t\overline{t}`$ pair acts as a dipole with separation $`\mathrm{}c/m10^3`$ fm, resulting in a small interaction cross section. Therefore, aside from debris from the target nucleon, the target nucleus will be relatively undisturbed, leaving the $`t\overline{t}`$ pair with relatively few accompanying particles.
In conclusion, top quark pairs will be produced via photon-gluon fusion in heavy ion collisions at the LHC. Heavy ion photoproduction will be the first accessible non-hadronic production channel for $`t\overline{t}`$ pairs. A $`10^6`$ s O+O run will produce 210 $`t\overline{t}`$ pairs while $`pA`$ runs will result in higher rates, up to 1100 pairs in $`p`$O collisions. These events could at least be statistically identifiable on the basis of accompanying rapidity gaps, the presence of an intact nucleus, and a slightly smaller multiplicity. The data should allow for measurements of the top charge and mass. In addition, the nuclear gluon distribution may be measurable at large $`Q^2`$.
This work was supported in part by the Division of Nuclear Physics of the Office of High Energy and Nuclear Physics of the U. S. Department of Energy under Contract No. DE-AC-03-76SF00098.
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# 1 Introduction
## 1 Introduction
Several integrable quantum spin chain models within the range of the algebraic Bethe ansatz method have a distinguished basis of states, which minimalizes quantum effects. That is, the quasiparticle creation and annihilation operators in this basis have an appearance devoid of polarization clouds. For the $`XXX`$ and $`XXZ`$ models with underlying group $`sl(2)`$ the bases in question were found by Maillet and Sanchez de Santos through the construction of a generalized Drinfeld twist .
The ensuing representation of the quantum monodromy matrices coincides, as noted by Terras , for the case of the rational $`XXX`$ model with the representation provided by Sklyanin’s functional Bethe ansatz method . An obvious generalization of Sklyanin’s method (substituting polynomials in the spectral parameter by polynomials in the exponential of the spectral parameter) leads us to the conclusion that an analogous coincidence holds true for the trigonometric $`XXZ`$ model.
The purpose of the present communication is to report on the generalization of the above results to the elliptic $`sl(2)`$-$`XYZ`$ model. For this sake we will make use of Baxter’s map of the $`XYZ`$ model onto an ice type model (which is akin to a cyclic solid on solid model ). This brings us formally near to the $`XXX`$ and $`XXZ`$ models and allows us to use the technique developed in for the handling of the analogous problem in the $`sl(n)`$-$`XXX`$ model. The plan of the paper is as follows : Section 2 provides a short survey of the $`XYZ`$ model and its reformulation as an ice-type model. Section 3 deals with the factorizing twists and the computation of the operator valued entries of the monodromy matrix. Section 4 contains the conclusions. The appendix is devoted to proving some relation needed in Sect.3.
## 2 $`XYZ`$ model and its relation to ice-type models
In the framework of the Algebraic Bethe Ansatz the $`XYZ`$-model is determined by the elliptic solution of the Yang-Baxter equation
$`R_{12}(\lambda _1\lambda _2)R_{13}(\lambda _1\lambda _3)R_{23}(\lambda _2\lambda _3)=R_{23}(\lambda _2\lambda _3)R_{13}(\lambda _1\lambda _3)R_{12}(\lambda _1\lambda _2)`$ (1)
with
$`R^{xyz}(\lambda \mu )=\left(\begin{array}{cccc}a& 0& 0& d\\ 0& b& c& 0\\ 0& c& b& 0\\ d& 0& 0& a\end{array}\right)`$ (6)
where
$`a(\lambda \mu )`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Theta }(2\eta )\mathrm{\Theta }(\lambda \mu )}{\mathrm{\Theta }(0)\mathrm{\Theta }(\lambda \mu +2\eta )}}`$
$`b(\lambda \mu )`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Theta }(2\eta )H(\lambda \mu )}{\mathrm{\Theta }(0)H(\lambda \mu +2\eta )}}`$
$`c(\lambda \mu )`$ $`=`$ $`{\displaystyle \frac{H(2\eta )\mathrm{\Theta }(\lambda \mu )}{\mathrm{\Theta }(0)H(\lambda \mu +2\eta )}}`$
$`d(\lambda \mu )`$ $`=`$ $`{\displaystyle \frac{H(2\eta )H(\lambda \mu )}{\mathrm{\Theta }(0)\mathrm{\Theta }(\lambda \mu +2\eta )}}`$ (7)
with the notation $`H(u)=\vartheta _1(\frac{u}{2K},q),\mathrm{\Theta }(u)=\vartheta _4(\frac{u}{2K},q)`$ and $`\vartheta _4=_m(1)^nq^{n^2}e^{2\pi inz},\vartheta _1(z,q)=iq^{\frac{1}{4}}e^{i\pi z}\theta _4(z+\tau /2,q)`$ are the standard theta-functions of a single complex variable.<sup>1</sup><sup>1</sup>1For a concise introduction into the theory of theta and elliptic functions we refer the reader to the appendix of , whose conventions we shall use throughout our work.
The somewhat different parametrization as compared to is due to the normalization in order to achieve unitarity for the R-matrix.
The monodromy matrix $`T(\lambda ,\{\lambda _i\})`$ (generalized to the inhomogeneous chain , ) is given as the ordered product of Lax operators $`L_i(\lambda \lambda _i)=R_{0i}(\lambda \lambda _i)`$
$`T(\lambda ,\left\{\lambda _i\right\})=L_N(\lambda z_N)\mathrm{}L_2(\lambda z_2)L_1(\lambda z_1)=\left(\begin{array}{cc}A(\lambda ,\left\{\lambda _i\right\})& B(\lambda ,\left\{\lambda _i\right\})\\ C(\lambda ,\left\{\lambda _i\right\})& D(\lambda ,\left\{\lambda _i\right\})\end{array}\right).`$ (10)
The presence of the Boltzmann weight $`d`$ in Eq. (6) reflects the non-conservation of spin, which is responsible for the absence of a local vacuum for the Lax operator associated with the above R-matrix.
To circumvent the problems arising from the eight vertex nature, we use the vertex–face map established by Baxter to obtain a $`XXZ`$ type (six vertex) R-matrix by exploiting the relation
$`R^{xyz}(\lambda \mu )\varphi _{l,l^{}}z_{m^{},l^{}}={\displaystyle \underset{m}{}}w(m,m^{}|l,l^{})\varphi _{m,m^{}}z_{m,l}`$ (11)
valid for all integers $`l,l^{},m,m^{}`$ such that $`|ll^{}|=|mm^{}|=1`$ and the summation on the r.h.s. is over integers s.t. $`|mm^{}|=1,|ml|=1`$.
The two dimensional vectors $`\varphi ,z`$ are given by
$`\varphi _{l,l+1}`$ $`=`$ $`X(s_l+\mu );z_{l+1,l}=X(s_l+\lambda )`$
$`\varphi _{l+1,l}`$ $`=`$ $`X(t_{l+1}\mu );z_{l1,l}=X(t_l\lambda )`$ (12)
with $`X(u)=\left(\genfrac{}{}{0pt}{}{H(u)}{\mathrm{\Theta }(u)}\right)`$ and the abbrevation $`s_l=s+2\eta l`$, where $`s,t`$ are arbitrary complex parameters.
Relation (11) provides a change of basis in which the Lax operator has a local vacuum independent of the spectral parameter $`\lambda `$. There exist a whole family of bases labeled by the integer $`l`$ and the parameters $`s,t`$, which achieve this goal.
The six weights $`w(m,m^{}|l,l^{})`$ are represented graphically in Figure 1.
FIG. 1 The six resulting Boltzmann weights.
The Boltzmann weights are subject to the star-triangle relation
$`{\displaystyle \underset{m}{}}w(l_2,l_3|l_1,m)w(l_3,l_4|m,l_5)w(m,l_5|l_1,l_6)={\displaystyle \underset{m}{}}w(l_3,l_4|l_2,m)w(l_2,m|l_1,l_6)w(m,l_4|l_6,l_5).`$ (13)
They are parametrized by $`\left(h(u)=\mathrm{\Theta }(0)H(u)\mathrm{\Theta }(u);\omega _l=(\frac{s+t}{2}+2\eta lK)\right)`$:
$`a_l`$ $`=`$ $`a_l^{}=1;`$
$`b_l`$ $`=`$ $`{\displaystyle \frac{h(\lambda )h(\omega _{l1})}{h(\lambda +2\eta )h(\omega _l)}};b_l^{}={\displaystyle \frac{h(\lambda )h(\omega _{l+1})}{h(\lambda +2\eta )h(\omega _l)}}`$
$`c_l`$ $`=`$ $`{\displaystyle \frac{h(2\eta )h(\omega _l\lambda )}{h(\lambda +2\eta )h(\omega _l)}};c_l^{}={\displaystyle \frac{h(2\eta )h(\omega _l+\lambda )}{h(\lambda +2\eta )h(\omega _l)}}`$ (14)
These weights can be arranged into a matrix
$`R_{12}(l)=\left(\begin{array}{cccc}a_l& 0& 0& 0\\ 0& b_l& c_l& 0\\ 0& c_l^{}& b_l^{}& 0\\ 0& 0& 0& a_l^{}\end{array}\right)`$ (19)
which fulfills the modified Yang-Baxter equation ,
$`R_{12}(l\sigma _3)R_{13}(l)R_{23}(l\sigma _1)=R_{23}(l)R_{13}(l\sigma _2)R_{12}(l).`$ (20)
The monodromy matrix related to this modified Yang-baxter equation is <sup>2</sup><sup>2</sup>2We use the convention $`X_{0,1\mathrm{}N}`$ to denote an operator $`X`$ represented in space $`0`$ with entries in the tensorproduct of spaces $`1,\mathrm{},N`$.
$`T_{0,1\mathrm{}N}(l)=R_{0N}(l\sigma _1\mathrm{}\sigma _{N1})\mathrm{}R_{02}(l\sigma _1)R_{01}(l)`$ (21)
where $`0`$ denotes the horizontal auxiliary space (with the asssociated spectral parameter $`\lambda _0`$), $`1,\mathrm{},N`$ label the vertical quantum spaces which span the physical Hilbertspace $`_N`$ (with associated local inhomogeneities $`\left\{\lambda _i\right\}`$), and $`\sigma _i`$ equals $`\pm 1`$ depending on whether the arrow in the $`i`$-th space is up or down (right/left for the horizontal space). It also sets our convention to associate the integer in the right lower corner of the graphical representation with the operator (cf. Figure 2).
FIG. 2 Elements of the monodromy matrix
From (20) follows the equation for the monodromy matrices
$`R_{00^{}}(l\sigma _1\mathrm{}\sigma _N)T_0(l)T_0^{}(l\sigma _0)=T_0^{}(l)T_0(l\sigma _0^{})R_{00^{}}(l)`$ (22)
and can be represented graphically as in Figure 3.
FIG. 3 The Yang–Baxter equation for the monodromy matrix
The unitarity relation $`R_{21}R_{12}=1\mathrm{I}`$ can be checked directly by using the following relation
$`h(u+x)h(ux)h(v+y)h(vy)h(u+y)h(uy)h(v+x)h(vx)`$
$`=h(xy)h(uy)h(v+x)h(v+x)`$ (23)
which is a consequence of the equality $`H(u,k)\mathrm{\Theta }(u,k)=H((1+k)u,2\sqrt{k}/(1+k))`$ with $`k`$ the nome of the elliptic function (Eq. 15.10.20 of ), and the analogue of (23) for $`H(u)`$ (Eq. 15.3.10 of ).
The representation of eigenvalues of the transfer matrix (10) through these of (21) is explained in reference (One has however to keep in mind that our model (21) differs from that of by an additional change of basis in the quantum space).
We will concentrate in what follows on the computation of a factorizing $`F`$-matrix for the monodromy matrix (21) for an arbitrary fixed value of $`l`$.
## 3 The $`F`$ basis
The factorizing $`F`$-matrix for two sites defined by the relation $`F_{21}R_{12}=F_{12}`$ is
$`F_{12}=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& c^{}& b^{}& 0\\ 0& 0& 0& 1\end{array}\right)_{[12]}.`$ (28)
The proof of the factorization property amounts to checking the same relations as those in the proof of the unitarity of the R-matrix above.
The factorizing $`F`$-matrix for $`N`$ sites ($`N`$ quantum spaces) turns out to be given by formally the same expression as found in for the $`XXX`$ model<sup>3</sup><sup>3</sup>3We thank Frank Göhmann for communicating this representation of $`F_{1\mathrm{}N}`$, which is slightly simpler than that quoted in .
$`F_{1\mathrm{}N}(l)`$ $`=`$ $`{\displaystyle \underset{\alpha _2^N}{}}P_\alpha R_{1\mathrm{}N}^{\sigma _\alpha }(l)(z_1,\mathrm{},z_N)`$
$`P_\alpha `$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}P_i^{\alpha _i}`$ (29)
where $`P_i^{\alpha _i}`$ projects on the $`\alpha _i`$-th component in the $`i`$-th space and the permutation $`\sigma _\alpha `$ is uniquely determined through the conditions
$`\alpha _{\sigma _\alpha (i+1)}`$ $``$ $`\alpha _{\sigma _\alpha (i)}\text{if}\sigma _\alpha (i+1)>\sigma _\alpha (i)`$
$`\alpha _{\sigma _\alpha (i+1)}`$ $`>`$ $`\alpha _{\sigma _\alpha (i)}\text{if}\sigma _\alpha (i+1)<\sigma _\alpha (i).`$ (30)
An algorithm for finding $`\sigma _\alpha `$ for a given $`\alpha `$ is described in the appendix.
The modification of the Yang-Baxter equation (20) enforces a particular rule for the handling of the integer valued parameter $`l`$ in the formation of the intertwining matrix $`R^\sigma (l)`$ (related to the permutation $`\sigma `$), which can be read off from the modified composition law
$`R^{\sigma \sigma _i}(l)`$ $`=`$ $`R_{\sigma (i),\sigma (i+1)}(\stackrel{~}{l}_i)R^\sigma (l)`$
$`\stackrel{~}{l}_i`$ $`=`$ $`l\sigma _{\sigma (1)}\mathrm{}\sigma _{\sigma (i1)}`$ (31)
where $`\sigma _i`$ is the transposition of $`i,i+1`$, and $`\sigma `$ an arbitrary permutation.
$`R^\sigma (l)`$ has the intertwining property
$`R^\sigma (l)T_{0,1\mathrm{}N}(l)=T_{0,\sigma (1)\mathrm{}\sigma (N)}(l)R^\sigma (l\sigma _0).`$
The matrix $`F_{1\mathrm{}N}(l)`$ satisfies the factorizing equation
$`R_{1\mathrm{}N}^\sigma (l)=F_{\sigma (1\mathrm{}N)}^1(l)F_{1\mathrm{}N}(l).`$ (32)
A proof of the latter equation can be found in – the modification of the composition law going along with the parameter $`l`$ being rather immaterial. An alternative argument emphasizing the geometric meaning of (32) (loosely speaking a pure gauge field representation of $`R^\sigma `$ on the symmetric group $`𝒮_N`$ as base space) may be sketched as follows:
It is sufficient to prove Eq. (32) for arbitrary elementary transpositions $`\sigma =\sigma _i`$; $`\sigma _i:(i,i+1)(i+1,i)`$,
$`R_{1\mathrm{}N}^{\sigma _i}=F_{\sigma _i(1\mathrm{}N)}^1F_{1\mathrm{}N}.`$ (33)
A special solution $`\widehat{F}^i`$ of (33) is the two site matrix $`F_{ii+1}`$, Eq. (28), embedded into the N-fold tensor product
$`\widehat{F}^i=1\mathrm{I}\mathrm{}F_{ii+1}\mathrm{}1\mathrm{I}`$
and the most general solution of (33) is given by
$`F_{1\mathrm{}N}=X(1,\mathrm{},\{i,i+1\},\mathrm{},N)\widehat{F}^i`$ (34)
where $`X(1,\mathrm{},\{i,i+1\},\mathrm{},N)`$ denotes a non-singular matrix in the quantum space $`^N^2`$ which is symmetric with respect to an exchange of the space labels $`i`$ and $`i+1`$ and of the respective local inhomogeneity parameters. So we note that the proof of Eq. (32) is complete if we can show that
$`F_{1\mathrm{}N}\left(\widehat{F}^i\right)^1=X_{(1,\mathrm{}N)}^{(i)}`$ (35)
is symmetric in the labels $`i`$ and $`i+1`$. For this purpose it is convenient to consider the permutation group element $`\sigma _\alpha `$ of Eq. (35) as being decomposed into products of elementary transpositions, which also leads to a factorization of $`R^{\sigma _\alpha }`$ into a product of elementary (two site) matrices $`R^{\sigma _x}`$, and to subdivide the sum $`_\alpha `$ in Eq. (29) into terms with the transposition $`\sigma _i`$ and into terms without. It is furthermore possible to arrange the order of the transpositions such that $`R^{\sigma _i}`$ does occur at most once and then on the very right end of the respective product of matrices $`R^{\sigma _x}`$. Schematically we obtain
$`{\displaystyle \underset{\alpha }{}}P_\alpha R^{\sigma _\alpha }`$ $`=`$ $`{\displaystyle \underset{\alpha ^{}}{}}P_\alpha ^{}R^{\sigma _\alpha ^{}}R^{\sigma _i}+{\displaystyle \underset{\alpha ^{\prime \prime }}{}}P_{\alpha ^{\prime \prime }}R^{\sigma _{\alpha ^{\prime \prime }}}.`$ (36)
Each factor $`P_\alpha ^{}R^{\sigma _\alpha ^{}}`$ in the first sum has a counterpart $`P_{\alpha ^{\prime \prime }}R^{\sigma _{\alpha ^{\prime \prime }}}`$ in the second sum related by an interchange of the labels $`i`$ and $`i+1`$.
Keeping this in mind one verifies straightforwardly by inspection that the multiplication of Eq. (36) by
$`\widehat{F}_{\sigma _i}^1`$ $`=`$ $`\left({\displaystyle \underset{2\alpha _i>\alpha _{i+1}}{}}1\mathrm{I}_{i,i+1}P_i^{\alpha _i}P_{i+1}^{\alpha _{i+1}}+{\displaystyle \underset{2\alpha _{i+1}\alpha _i}{}}R^{\sigma _i}P_i^{\alpha _i}P_{i+1}^{\alpha _{i+1}}\right)\times diag\{1,b^1(\lambda _j,\lambda _i),b^1(\lambda _i,\lambda _j),1\}`$
from the right side produces an expression which is symmetric in $`i`$ and $`i+1`$.
The operators of the monodromy matrix (21) in the $`F`$ basis are obtained by using a recursion relation, which enables one to express the monodromy matrix for $`N`$ sites in terms of that for $`N1`$ sites. Starting from the definition $`\stackrel{~}{T}_{0,1\mathrm{}N}(l)=F_{1\mathrm{}N}(l)T_{0,1\mathrm{}N}(l)F_{1\mathrm{}N}^1(l\sigma _0)`$ we use the identity $`F_{1\mathrm{}N}(l)=F_{2\mathrm{}N}(l\frac{1+\sigma _1}{2})F_{1,2\mathrm{}N}(l)`$ with $`F_{1,2\mathrm{}N}(l)=\left(P_1^1+P_1^2T_{1,23\mathrm{}N}\right)`$ (for the proof of this identity we refer the reader to the appendix)
$`\stackrel{~}{T}_{0,1\mathrm{}N}(l)=F_{1\mathrm{}N}(l)T_{0,1\mathrm{}N}(l)F_{1\mathrm{}N}^1(l\sigma _0)`$
$`=`$ $`F_{2\mathrm{}N}(l{\displaystyle \frac{1+\sigma _1}{2}})F_{1,2\mathrm{}N}(l)T_{0,1\mathrm{}N}(l)F_{1,2\mathrm{}N}^1(l\sigma _0)F_{2\mathrm{}N}^1(l\sigma _0{\displaystyle \frac{1+\sigma _1}{2}}).`$
Splitting $`T_{0,12\mathrm{}N}(l)`$ into $`T_{0,2\mathrm{}N}(l\sigma _1)R_{01}(l)`$ we arrive at
$`\stackrel{~}{T}_{0,1\mathrm{}N}(l)`$
$`=`$ $`F_{2\mathrm{}N}(l{\displaystyle \frac{1+\sigma _1}{2}})F_{1,2\mathrm{}N}(l)F_{2\mathrm{}N}^1(l\sigma _1)F_{2\mathrm{}N}(l\sigma _1)T_{0,2\mathrm{}N}(l\sigma _1)R_{01}(l)\times `$
$`\times F_{2\mathrm{}N}^1(l\sigma _0\sigma _1)F_{2\mathrm{}N}(l\sigma _0\sigma _1)F_{1,2\mathrm{}N}^1(l\sigma _0)F_{2\mathrm{}N}^1(l\sigma _0{\displaystyle \frac{1+\sigma _1}{2}})`$
$`=`$ $`\stackrel{~}{F}_{1,2\mathrm{}N}(l)\stackrel{~}{T}_{0,2\mathrm{}N}(l\sigma _1)R_{01}(l)\stackrel{~}{F}_{1,2\mathrm{}N}^1(l\sigma _0)`$
Inserting the explicit form of
$`\stackrel{~}{F}_{1,2\mathrm{}N}(l)`$ $`=`$ $`F_{2\mathrm{}N}\left(l{\displaystyle \frac{1+\sigma _1}{2}}\right)F_{1,2\mathrm{}N}(l)F_{2\mathrm{}N}^1(l\sigma _1)=\left(\begin{array}{cc}1\mathrm{I}& 0\\ \stackrel{~}{C}_{2\mathrm{}N}^1(l)& \stackrel{~}{D}_{2\mathrm{}N}^1(l)\end{array}\right)_{[1]}`$
we obtain the final expression
$`\stackrel{~}{T}_{0,1\mathrm{}N}(l)=\left(\begin{array}{cc}1\mathrm{I}& 0\\ \stackrel{~}{C}_{2\mathrm{}N}^1(l)& \stackrel{~}{D}_{2\mathrm{}N}^1(l)\end{array}\right)_{[1]}\stackrel{~}{T}_{0,2\mathrm{}N}(l\sigma _1)R_{01}(l)\left(\begin{array}{cc}1\mathrm{I}& 0\\ \stackrel{~}{C}_{2\mathrm{}N}^1(l\sigma _0)& \stackrel{~}{D}_{2\mathrm{}N}^1(l\sigma _0)\end{array}\right)_{[1]}^1.`$ (42)
This relation can be solved recursively starting with the one site monodromy matrix which coincides with the Lax operator $`L_i=R_{0i}`$ (19).
We demonstrate the derivation and possible difficulties and their resolution in the case of the operator $`\stackrel{~}{D}_{1\mathrm{}N}(l)`$:
$`\stackrel{~}{D}_{1\mathrm{}N}(l)`$ $`=`$ $`\left(\begin{array}{cc}\stackrel{~}{D}_{2\mathrm{}N}^0(l1)b_{01}^{}(l)& 0\\ X& \stackrel{~}{D}_{2\mathrm{}N}^1(l)\stackrel{~}{D}_{2\mathrm{}N}^0(l+1)\left(\stackrel{~}{D}_{2\mathrm{}N}^1\right)^1(l+1)\end{array}\right)_{[1]}`$
where the lower left element $`X`$ is given by
$`\stackrel{~}{C}_{2\mathrm{}N}^1(l)\stackrel{~}{D}_{2\mathrm{}N}^0(l1)b_{01}^{}(l)+\stackrel{~}{D}_{2\mathrm{}N}^1(l)\stackrel{~}{C}_{2\mathrm{}N}^0(l+1)c_{01}(l)\stackrel{~}{D}_{2\mathrm{}N}^0(l)\stackrel{~}{C}_{2\mathrm{}N}^1(l)`$
This combination vanishes identically by virtue of the Yang–Baxter equation (22), or in graphical representation:
We thus end with a diagonal matrix of the form
$`\stackrel{~}{D}_{1\mathrm{}N}^0(l)`$ $`=`$ $`\left(\begin{array}{cc}\stackrel{~}{D}_{2\mathrm{}N}^0(l1)b_{01}^{}(l)& 0\\ 0& \stackrel{~}{D}_{2\mathrm{}N}^0(l)\end{array}\right)_{[1]}`$
which can now be solved iteratively
$`\left(\begin{array}{cc}\stackrel{~}{D}_{2\mathrm{}N}^0(l1)b_{01}^{}(l)& 0\\ 0& \stackrel{~}{D}_{2\mathrm{}N}^0(l)\end{array}\right)_{[1]}=\stackrel{~}{D}_{2\mathrm{}N}^0(l{\displaystyle \frac{1+\sigma _1}{2}})\left(\begin{array}{cc}b_{01}^{}(l)& 0\\ 0& 1\end{array}\right)_{[1]}`$
$`=`$ $`\stackrel{~}{D}_{3\mathrm{}N}^0(l{\displaystyle \frac{1+\sigma _1}{2}}{\displaystyle \frac{1+\sigma _2}{2}})\left(\begin{array}{cc}b_{02}^{}(l\frac{1+\sigma _1}{2})& 0\\ 0& 1\end{array}\right)_{[2]}\left(\begin{array}{cc}b_{01}^{}(l)& 0\\ 0& 1\end{array}\right)_{[1]}`$
$`\mathrm{}`$
$`=`$ $`_{i=1}^N\left(\begin{array}{cc}b_{0i}^{}(l\frac{1+\sigma _1}{2}\frac{1+\sigma _2}{2}\mathrm{}\frac{1+\sigma _{i1}}{2})& 0\\ 0& 1\end{array}\right)_{[i]}.`$
The computation of the operators $`\stackrel{~}{B}_{1\mathrm{}N}(l),\stackrel{~}{C}_{1\mathrm{}N}(l)`$ proceeds along the same lines to give
$`\stackrel{~}{B}_{1\mathrm{}N}(l)`$ $`=`$ $`\left(\begin{array}{cc}\stackrel{~}{B}_{2\mathrm{}N}^0(l1)b_{01}^{}(l)& 0\\ Y& \stackrel{~}{D}_{2\mathrm{}N}^1(l)\stackrel{~}{B}_{2\mathrm{}N}^0(l+1)\left(\stackrel{~}{D}_{2\mathrm{}N}^1\right)^1(l+1)\end{array}\right)_{[1]}`$
with
$$Y=\stackrel{~}{C}_{2\mathrm{}N}^0(l)\stackrel{~}{B}_{2\mathrm{}N}^0(l1)b_{01}^{}(l)+\stackrel{~}{D}_{2\mathrm{}N}^1(l)\stackrel{~}{A}_{2\mathrm{}N}^0(l+1)c_{01}(l)\stackrel{~}{D}_{2\mathrm{}N}^1(l)\stackrel{~}{B}_{2\mathrm{}N}^0(l+1)\left(\stackrel{~}{D}_{2\mathrm{}N}^1\right)^1(l+1)\stackrel{~}{C}_{2\mathrm{}N}^1(l+1)$$
and for $`\stackrel{~}{C}`$
$`\stackrel{~}{C}_{1\mathrm{}N}(l)`$ $`=`$ $`\left(\begin{array}{cc}Z& \stackrel{~}{D}_{2\mathrm{}N}^0(l1)\left(\stackrel{~}{D}_{2\mathrm{}N}^0\right)^1(l1)c_{01}^{}(l)\\ 0& \stackrel{~}{C}_{2\mathrm{}N}^0(l)\end{array}\right)_{[1]}`$
with
$$Z=\stackrel{~}{C}_{2\mathrm{}N}^0(l1)\stackrel{~}{D}_{2\mathrm{}N}^0(l1)\left(\stackrel{~}{D}_{2\mathrm{}N}^0\right)^1(l1)\stackrel{~}{C}_{2\mathrm{}N}^1(l1)c_{01}^{}(l).$$
The solution of the recursion relation finally yields the result for the operators of the monodromy matrix $`(`$we use a slight change in notation: $`b(\lambda )=\frac{h(\lambda )}{h(\lambda +2\eta )}`$ and denote $`k=l_{i=1}^{N1}\sigma _i`$$`)`$:
$`\stackrel{~}{D}_l(\lambda _0)`$ $`=`$ $`{\displaystyle \frac{h(\omega _{l+1})}{h(\omega _{1+\frac{k+lN}{2}})}}_{i=1}^N\left(\begin{array}{cc}b(\lambda _0\lambda _i)& 0\\ 0& 1\end{array}\right)_{[i]}`$ (53)
$`\stackrel{~}{B}_l(\lambda _0)`$ $`=`$ $`{\displaystyle \frac{h(\omega _{l+1})}{h(\omega _k)}}{\displaystyle \underset{i=1}{\overset{N}{}}}c_{k1}(\lambda _0\lambda _i)\sigma _i^{}_{ji}^N\left(\begin{array}{cc}b(\lambda _0\lambda _j)& 0\\ 0& b^1(\lambda _j\lambda _i)\end{array}\right)_{[j]}`$ (56)
$`\stackrel{~}{C}_l(\lambda _0)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}c_l^{}(\lambda _0\lambda _i)\sigma _i^+_{ji}^N\left(\begin{array}{cc}b(\lambda _0\lambda _j)b^1(\lambda _i\lambda _j)& 0\\ 0& 1\end{array}\right)_{[j]}`$ (59)
(60)
The above mentioned basis transformation (11) amounts to splitting the model into sectors with a fixed number of turned spins. To obtain the spectrum of the $`XYZ`$ model one uses the operators $`\stackrel{~}{B}_l(\lambda ),\stackrel{~}{C}_l(\lambda )`$ to construct eigenvectors in the form proposed by (there denoted by $`B_{k,l}(\lambda )`$ etc.).
The remaining operator $`\stackrel{~}{A}_l(\lambda _0)`$ is obtained from the quantum determinant $`det_qT(\lambda _0)`$ ,
$`det_qT(\lambda _0)`$ $`=`$ $`{\displaystyle \frac{h(\omega _{k+1})}{h(\omega _{l+1})}}D_{k,l}(\lambda _0)A_{k+1,l+1}(\lambda _0+2\eta ){\displaystyle \frac{h(\omega _{l+1})}{h(\omega _k)}}B_{k,l}(\lambda _0)C_{k1,l+1}(\lambda _0+2\eta )`$ (61)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{h(\lambda _0\lambda _i)}{h(\lambda _0\lambda _i+2\eta )}}`$
As the operator $`D_l`$ is diagonal it can be inverted easily, and the operator $`C_l`$ possesses a global vacuum which is annihilated by its action. This enables one to compute $`A_l`$
$`\stackrel{~}{A}_l(\lambda _0)`$ $`=`$ $`{\displaystyle \frac{h(\omega _{\frac{k+lN}{2}})}{h(\omega _k)}}\{_{i=1}^N\left(\begin{array}{cc}1& 0\\ 0& b(\lambda _i\lambda _0)^1\end{array}\right)_{[i]}`$ (62)
$`+`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{c_{k1}(\lambda _0\lambda _i)c_l^{}(\lambda _0\lambda _i)}{b(\lambda _0\lambda _i)}}\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)_{[i]}_{ji}\left(\begin{array}{cc}\frac{b(\lambda _0\lambda _j)}{b(\lambda _i\lambda _j)}& 0\\ 0& b(\lambda _j\lambda _i)^1\end{array}\right)_{[j]}`$
$`+`$ $`{\displaystyle \underset{ij}{\overset{N}{}}}{\displaystyle \frac{c_{k1}(\lambda _0\lambda _i)c_l^{}(\lambda _0\lambda _j)}{b(\lambda _j\lambda _k)}}\sigma _{}^i\sigma _+^j_{ki,j}\left(\begin{array}{cc}\frac{b(\lambda _0\lambda _k)}{b(\lambda _j\lambda _k)}& 0\\ 0& b(\lambda _k\lambda _i)^1\end{array}\right)_{[k]}\}`$
## 4 Conclusion
The form of the $`F`$-matrix, Eq. (29) and the appearance of the monodromy matrix in the basis supplied by the $`F`$-matrix, Eq’s (60), (62), are completely analogous to what has been found in and for the rational and trigonometric models. The concrete expressions for $`\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C},\stackrel{~}{D}`$ are in particular manifestly symmetric with respect to exchanges of the local inhomogeneity parameters $`\lambda _i`$. The quasiparticle operators $`\stackrel{~}{B}`$ and $`\stackrel{~}{C}`$ are free from polarization effects due to non-local exchange terms.
The argument used in \- borrowed from \- concerning the identification of operators corresponding to different entries of the monodromy matrix relied on the $`sl(n)`$ symmetry of the rational model. It is not available for the trigonometric and elliptic model. The recursive procedure followed instead in the preceding section is equally applicable to the $`XXX`$, $`XXZ`$ and $`XYZ`$ model.
It seems rather plausible in view of the formal similarities of the rational, trigonometric and elliptic models that some version of Sklyanin’s functional Bethe ansatz should also be feasible in the latter case as has already been achieved for the $`XYZ`$ Gaudin magnet in .
Acknowledgement: We thank Frank Göhmann for a discussion. H.B. and R.H.P. acknowledge support of the Alexander von Humboldt Foundation. R.F. was supported by the TMR network contract FMRX-CT96-0012 of the European Commission.
## Appendix
In this appendix we derive the identity used in the recursion relation (42) using the factorizing $`F`$-matrix for $`N`$ sites (for the sake of simplicity we omit the explicit dependence on $`l`$, which can be restored by using Eq.(31)):
$`F_{1\mathrm{}N}`$ $`=`$ $`{\displaystyle \underset{\alpha _2^N}{}}P_\alpha R_{1\mathrm{}N}^{\sigma _\alpha }(z_1,\mathrm{},z_N)`$
$`P_\alpha `$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}P_i^{\alpha _i}`$ (63)
and $`P_i^{\alpha _i}`$ projects on the $`\alpha _i`$-th component in the $`i`$-th space.
With a given $`\alpha `$ the permutation $`\sigma _\alpha 𝒮_N`$ is uniquely determined by the following algorithm:
Let $`k`$ be the number of 1’s in the sequence $`\alpha _1,\mathrm{},\alpha _N`$. For $`ik`$ we put $`\sigma _\alpha (i)=l`$, where $`l`$ is the label of the $`i`$-th 1 in the above mentioned sequence.
For $`i>k`$ we put $`\sigma _\alpha (i)=k+l^{}`$, where $`l^{}`$ is the label of the $`ik`$-th 2 in the sequence.
The permutation constructed this way is the only one satisfying the constraints
$`\alpha _{\sigma _\alpha (i+1)}`$ $``$ $`\alpha _{\sigma _\alpha (i)}\text{if}\sigma _\alpha (i+1)>\sigma _\alpha (i)`$
$`\alpha _{\sigma _\alpha (i+1)}`$ $`>`$ $`\alpha _{\sigma _\alpha (i)}\text{if}\sigma _\alpha (i+1)<\sigma _\alpha (i).`$ (64)
Now we will prove that
$$F_{12\mathrm{}N}=F_{23\mathrm{}N}\left(P_1^1+P_1^2T_{1,23\mathrm{}N}\right).$$
If $`\alpha _2^N`$ has the form $`(1,\alpha _2,\mathrm{},\alpha _N)(1,\stackrel{~}{\alpha });\stackrel{~}{\alpha }=(\alpha _2,\mathrm{},\alpha _N)_2^{N1}`$, we have $`P_\alpha =P_1^1P_{\stackrel{~}{\alpha }}`$ and $`\sigma _\alpha =\sigma _{\stackrel{~}{\alpha }}`$ (here and in what follows we identify the symmetric group $`𝒮_{N1}`$ acting on the elements $`(2,\mathrm{},N)`$ as the subgroup of $`𝒮_N`$ acting on the elements $`(1,2,\mathrm{},N)`$). Thus the part of the sum of (63) corresponding to such $`\alpha `$’s exactly gives $`P_1^1F_{2\mathrm{}N}`$.
Now consider the case when $`\alpha =(2,\alpha _2,\mathrm{},\alpha _N)(2,\stackrel{~}{\alpha })`$. It is easy to check that $`P_\alpha =P_1^2P_{\stackrel{~}{\alpha }}`$ and $`\sigma _\alpha =\sigma _{\stackrel{~}{\alpha }}\sigma _{1,\stackrel{~}{2}}\mathrm{}\sigma _{1,\stackrel{~}{k+1}}`$ where $`k`$ is the number of 1’s in $`\stackrel{~}{\alpha }`$ and $`\sigma _{i,j}`$ are elementary transpositions and $`\stackrel{~}{i}=\sigma _{\stackrel{~}{\alpha }}(i)`$. We have
$`P_1^2P_{\stackrel{~}{\alpha }}R^{\sigma _{\stackrel{~}{\alpha }}}T_{1,23\mathrm{}N}`$ $`=`$ $`P_\alpha T_{1,\sigma _{\stackrel{~}{\alpha }}(2)\mathrm{}\sigma _{\stackrel{~}{\alpha }}(N)}R^{\sigma _{\stackrel{~}{\alpha }}}`$ (65)
$`=`$ $`P_\alpha R_{1,\sigma _{\stackrel{~}{\alpha }}(k+1)}\mathrm{}R_{1,\sigma _{\stackrel{~}{\alpha }}(2)}R^{\sigma _{\stackrel{~}{\alpha }}}`$
$`=`$ $`P_\alpha R^{\sigma _{\stackrel{~}{\alpha }}\sigma _{1,\stackrel{~}{2}}\mathrm{}\sigma _{1,\stackrel{~}{k+1}}}=P_\alpha R^{\sigma _\alpha }`$
where we have taken into account that $`\stackrel{~}{\alpha }_{\stackrel{~}{k+2}}=\mathrm{}=\stackrel{~}{\alpha }_{\stackrel{~}{N}}=2`$ and that $`P_1^2P_i^2R_{1i}=P_1^2P_i^21\mathrm{I}`$.
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# Plane–fronted waves in metric–affine gravity
## I Introduction
Even though Einstein’s general relativity appears almost fully corroborated experimentally, there are several reasons to believe that the validity of such a description is limited to macroscopic structures and to the present cosmological era. The only available finite perturbative treatment of quantum gravity, namely the theory of the quantum superstring , suggests that non–Riemannian features are present on the scale of the Planck length. On the other hand, recent advances in the study of the early universe, as represented by the inflationary model, involve, in addition to the metric tensor, at the very least a scalar dilaton induced by a Weyl geometry, i.e. again an essential departure from Riemannian metricity . Even at the classical cosmological level, a dilatonic field has recently been used to describe the presence of dark matter in the universe as well as to explain certain cosmological observations which contradict the fundaments of the standard cosmological model .
There is good experimental evidence that, at the present state of the universe, the geometrical structure of spacetime corresponds to a metric–compatible geometry in which nonmetricity, but not necessarily the torsion, vanishes. Consequently, the full metric–affine geometry is irrelevant for the geometrical description of the universe today. However, during the early universe, when the energies of the cosmic matter were much higher than today, we expect scale invariance to prevail and, according to MAG, the canonical dilation (or scale) current of matter, i.e. the trace of the hypermomentum current $`\mathrm{\Delta }^\gamma _\gamma `$ becomes coupled to the Weyl covector $`Q:=\frac{1}{4}g^{\alpha \beta }Q_{\alpha \beta }`$, here $`Q_{\alpha \beta }:=Dg_{\alpha \beta }`$ is the nonmetricity of spacetime. Moreover, shear type excitations of the material multispinors (Regge trajectory type of constructs) are expected to arise, thereby liberating the (metric–compatible) Riemann–Cartan spacetime from its constraint of vanishing nonmetricity $`Q_{\alpha \beta }=0`$. It is therefore important to derive and investigate exact solutions of these theories which contain information about the new geometric objects like torsion and nonmetricity (for a survey of these theories see ).
For restricted irreducible pieces of torsion and nonmetricity, there are similarities between the Einstein–Maxwell system and the vacuum MAG field equations . This observation encourages us to find new solutions for MAG theories . However, the coupling of the post–Riemannian structures of a metric–affine spacetime to matter is still under investigation.
The search for plane–fronted wave solutions in MAG was first restricted to its Einstein–Cartan sector . Later, plane wave solutions with non–vanishing nonmetricity were found by Tucker et al. . Colliding waves with the appropriate metric and an excited post–Riemannian triplet are studied in , the corresponding generalization to the electrovac case can be found in .
In this paper we study plane–fronted gravitational and electromagnetic waves in metric–affine gravity theories with nonzero cosmological constant in their triplet ansatz sector. The plane–fronted electrovacuum–MAG waves comprize curvature, nonmetricity, torsion, and an electromagnetic field.
The plan of the paper is as follows: In Sec. II, we review the plane–fronted gravitational and electromagnetic waves in Einstein–Maxwell theory. In Sec. III, we present the plane–fronted gravitational and electromagnetic waves in MAG. In Sec. IV, we specialize to particular wave solutions and, in Sec. V, we discuss the results.
## II Plane–fronted gravitational and electromagnetic waves in Einstein–Maxwell theory
In this section we summarize the main results of Ref. : Using the null tetrad formalism, in a coordinate system $`(\rho ,\sigma ,\zeta ,\overline{\zeta })`$ (the bar denotes complex conjugation), the metric reads
$$ds^2=2\left(\vartheta ^{\widehat{0}}\vartheta ^{\widehat{1}}+\vartheta ^{\widehat{2}}\vartheta ^{\widehat{3}}\right),$$
(1)
with the coframe
$$\vartheta ^{\widehat{0}}=\frac{1}{p}d\zeta ,\vartheta ^{\widehat{1}}=\frac{1}{p}d\overline{\zeta },\vartheta ^{\widehat{2}}=d\sigma ,\vartheta ^{\widehat{3}}=\left(\frac{q}{p}\right)^2\left(sd\sigma +d\rho \right),$$
(2)
where the structural functions $`p`$, $`q`$, and $`s`$ are given as follows:
$`p(\zeta ,\overline{\zeta })`$ $`=`$ $`1+{\displaystyle \frac{\lambda _{\text{cosm}}}{6}}\zeta \overline{\zeta },`$ (3)
$`q(\sigma ,\zeta ,\overline{\zeta })`$ $`=`$ $`\left(1{\displaystyle \frac{\lambda _{\text{cosm}}}{6}}\zeta \overline{\zeta }\right)\alpha (\sigma )+\zeta \overline{\beta }(\sigma )+\overline{\zeta }\beta (\sigma ),`$ (4)
$`s(\rho ,\sigma ,\zeta ,\overline{\zeta })`$ $`=`$ $`{\displaystyle \frac{\lambda _{\text{cosm}}}{6}}\rho ^2\alpha ^2(\sigma )\rho ^2\beta (\sigma )\overline{\beta }(\sigma )+\rho _\sigma \left(\mathrm{ln}q\right)+{\displaystyle \frac{p}{2q}}H(\sigma ,\zeta ,\overline{\zeta }).`$ (5)
Here $`\alpha `$, $`\beta `$, and $`H`$ are arbitrary functions.
Let $`\stackrel{~}{R}_{\alpha \beta }`$ denote the Riemannian part of the curvature 2-form. Then we can subtract out the irreducible scalar curvature piece
$${}_{}{}^{(6)}\stackrel{~}{R}_{\alpha \beta }^{}:=\frac{1}{12}(e_\nu e_\mu \stackrel{~}{R}^{\nu \mu })\vartheta _\alpha \vartheta _\beta ,$$
(6)
see , and can define the 2-form
$$S_{\alpha \beta }:=\stackrel{~}{R}_{\alpha \beta }^{(6)}\stackrel{~}{R}_{\alpha \beta }=^{(1)}\stackrel{~}{R}_{\alpha \beta }+^{(4)}\stackrel{~}{R}_{\alpha \beta }=C_{\alpha \beta }+^{(4)}\stackrel{~}{R}_{\alpha \beta }.$$
(7)
Here $`e_\alpha `$ denotes the (vector) frame dual to the coframe $`\vartheta ^\alpha `$. If the Einstein vacuum field equations (with or without cosmological constant) are fulfilled — in this specific case $`{}_{}{}^{(4)}\stackrel{~}{R}_{\alpha \beta }^{}=0`$ —, then $`S_{\alpha \beta }`$ becomes the Weyl conformal curvature 2-form $`C_{\alpha \beta }:=^{(1)}\stackrel{~}{R}_{\alpha \beta }`$. Moreover, we will introduce the propagation 1-form $`k:=k_\mu \vartheta ^\mu `$ which inherits the properties of the geodesic, shear–free, expansion–free and twistless null vector–field $`k^\mu `$ representing the propagation vector of a plane–fronted wave.
The gravitational and null electromagnetic fields are subject to the radiation conditions
$$S_{\alpha \beta }k=0,(e^\alpha k)S_{\alpha \beta }=0,$$
(8)
and
$$Fk=0,(e^\alpha k)e_\alpha F=0.$$
(9)
In the following we will solve the Einstein–Maxwell equations (for the notion compare )
$`{\displaystyle \frac{1}{2}}\eta _{\alpha \beta \gamma }\stackrel{~}{R}^{\beta \gamma }+\lambda _{\mathrm{cosm}}\eta _\alpha `$ $`=`$ $`\kappa \mathrm{\Sigma }_\alpha ^{\mathrm{Max}},`$ (10)
$`dF`$ $`=`$ $`0,`$ (11)
$`d{}_{}{}^{}F`$ $`=`$ $`0,`$ (12)
where $`\mathrm{\Sigma }_\alpha ^{\mathrm{Max}}`$ represents the energy-momentum 3-form of the Maxwell field given by
$$\mathrm{\Sigma }_\alpha ^{\mathrm{Max}}:=\frac{1}{2}[(e_\alpha F){}_{}{}^{}F(e_\alpha ^{}F)F].$$
(13)
Writing the electromagnetic field as
$`F`$ $`=`$ $`{\displaystyle \frac{1}{2}}F_{ab}dx^adx^b`$ (14)
$`=`$ $`f(\zeta ,\sigma )d\zeta d\sigma +\overline{f}(\overline{\zeta },\sigma )d\overline{\zeta }d\sigma ,`$ (15)
with $`f(\zeta ,\sigma )`$ an arbitrary function of its arguments, one finds for the energy-momentum 3-form of the electromagnetic field as nonvanishing component
$$\mathrm{\Sigma }_{\widehat{2}}^{\mathrm{Max}}=2p^2f\overline{f}\vartheta ^{\widehat{0}}\vartheta ^{\widehat{1}}\vartheta ^{\widehat{2}},$$
(16)
in agreement with the result for $`T_{ab}`$ mentioned in Ref. Eq. (3.7).
The surfaces of constant $`\sigma `$ are the wave fronts of the electromagnetic waves. The above conditions (8)–(9) restrict the function $`\alpha (\sigma )`$ to the real domain whereas $`\beta (\sigma )`$ can be complex valued.
The function $`H`$, for a combined gravitational and electromagnetic wave, has to fulfill the equation
$$H_{,\zeta \overline{\zeta }}+\frac{\lambda _{\text{cosm}}}{3p^2}H=\frac{2\kappa p}{q}f\overline{f}.$$
(17)
In order to solve this non–homogeneous equation, one observes that a complex combination of an arbitrary holomorphic function $`\mathrm{\Phi }=\mathrm{\Phi }(\zeta ,\sigma )`$ of the form $`\mathrm{\Phi }_{,\zeta }(\lambda _{\text{cosm}}/3)(\overline{\zeta }/p)\mathrm{\Phi }`$ is the general complex solution to the corresponding homogeneous equation of (17). Thus, the real $`H_\mathrm{h}`$–solution to the homogeneous equation is given by
$$H_\mathrm{h}=\mathrm{\Phi }_{,\zeta }\frac{\lambda _{\text{cosm}}}{3}\frac{\overline{\zeta }}{p}\mathrm{\Phi }+\overline{\mathrm{\Phi }}_{,\overline{\zeta }}\frac{\lambda _{\text{cosm}}}{3}\frac{\zeta }{p}\overline{\mathrm{\Phi }}.$$
(18)
This structure sheds light on how to find the general solution of the non–homogeneous equation. Let us look for the particular solution $`H_\mathrm{p}`$ of the form
$$H_\mathrm{p}=\mu _{,\zeta }\frac{\lambda _{\text{cosm}}}{3}\frac{\overline{\zeta }}{p}\mu +\overline{\mu }_{,\overline{\zeta }}\frac{\lambda _{\text{cosm}}}{3}\frac{\zeta }{p}\overline{\mu },$$
(19)
where $`\mu =\mu (\sigma ,\zeta ,\overline{\zeta })`$, such that the function
$$H_{(1)}:=\mu _{,\zeta }\frac{\lambda _{\text{cosm}}}{3}\frac{\overline{\zeta }}{p}\mu ,$$
(20)
satisfies the equation
$$H_{(1),\zeta \overline{\zeta }}+\frac{\lambda _{\text{cosm}}}{3}\frac{H_{(1)}}{p^2}=\frac{\kappa p}{q}f\overline{f}.$$
(21)
Then it follows that $`\mu `$ itself is subject to
$$\left(\mu _{,\overline{\zeta }}\right)_{,\zeta \zeta }\frac{\lambda _{\text{cosm}}}{3}\left(\frac{\overline{\zeta }}{p}\mu _{,\overline{\zeta }}\right)_{,\zeta }=\frac{\kappa p}{q}f\overline{f},$$
(22)
with the general solution
$$\mu =\kappa ^{\overline{\zeta }}𝑑\overline{\zeta }p^2^\zeta \frac{d\zeta ^{}}{p^2}^\zeta ^{}𝑑\zeta ^{\prime \prime }\frac{p}{q}f\overline{f}.$$
(23)
For any given function $`f`$ one integrates for $`\mu `$ and, by using (19), one obtains $`H_\mathrm{p}`$. The general $`H`$ is constructed simply by adding the homogeneous solution $`H_\mathrm{h}`$ to $`H_\mathrm{p}`$,
$$H=H_\mathrm{h}+H_\mathrm{p}.$$
(24)
The general solution $`H`$ is characterized by the selfdual part of the conformal Weyl 2–form
$${}_{}{}^{+}C_{\alpha \beta }^{}:=\frac{1}{2}(C_{\alpha \beta }+i{}_{}{}^{}C_{\alpha \beta }^{}),$$
(25)
the trace–free Ricci 1–form
$$\stackrel{~}{R}_\alpha :=e_\beta \stackrel{~}{R}_\alpha ^\beta \frac{1}{4}\stackrel{~}{R}\vartheta _\alpha ,$$
(26)
the Ricci scalar
$$\stackrel{~}{R}:=e_\alpha e_\beta \stackrel{~}{R}^{\alpha \beta },$$
(27)
and the electromagnetic 2–form $`F`$. The ansatz (1)–(5) yields
$`{}_{}{}^{+}C_{\widehat{2}\widehat{0}}^{}`$ $`=`$ $`^+C_{\widehat{0}\widehat{2}}={\displaystyle \frac{1i}{4}}pq\left(H_{,\zeta \zeta }+{\displaystyle \frac{\lambda _{\mathrm{cosm}}}{3}}{\displaystyle \frac{\overline{\zeta }}{p}}H_{,\zeta }\right)\vartheta ^{\widehat{0}}\vartheta ^{\widehat{2}},`$ (28)
$`{}_{}{}^{+}C_{\widehat{2}\widehat{1}}^{}`$ $`=`$ $`^+C_{\widehat{1}\widehat{2}}={\displaystyle \frac{1+i}{4}}pq\left(H_{,\overline{\zeta }\overline{\zeta }}+{\displaystyle \frac{\lambda _{\mathrm{cosm}}}{3}}{\displaystyle \frac{\zeta }{p}}H_{,\overline{\zeta }}\right)\vartheta ^{\widehat{1}}\vartheta ^{\widehat{2}},`$ (29)
$`\stackrel{~}{R}_{\widehat{2}}`$ $`=`$ $`pq\left(H_{,\zeta \overline{\zeta }}+{\displaystyle \frac{\lambda _{\mathrm{cosm}}}{3p^2}}H\right)\vartheta ^{\widehat{2}}=2\kappa p^2f\overline{f}\vartheta ^{\widehat{2}},`$ (30)
$`\stackrel{~}{R}`$ $`=`$ $`4\lambda _{\mathrm{cosm}},`$ (31)
$`F`$ $`=`$ $`dA=d\left[\left({\displaystyle ^\zeta }f(\zeta ^{},\sigma )𝑑\zeta ^{}+{\displaystyle ^{\overline{\zeta }}}\overline{f}(\overline{\zeta }^{},\sigma )𝑑\overline{\zeta }^{}\right)\vartheta ^{\widehat{2}}\right].`$ (32)
The Weyl 2–form could be written still a bit more compactly according to
$$pq\left(H_{,\zeta \zeta }+\frac{\lambda _{\mathrm{cosm}}}{3}\frac{\overline{\zeta }}{p}H_{,\zeta }\right)=_\zeta \left[q^2_\zeta \left(\frac{p}{q}H\right)\right],$$
(33)
but the form given above is more practical if a certain function $`H`$ is explicitly given and calculations need to be done.
It is worthwhile to mention the existence of a conformally flat solution given by
$$H=\frac{1}{p}\left(u+\overline{v}\zeta +v\overline{\zeta }+w\zeta \overline{\zeta }\right),$$
(34)
where $`u`$, $`v`$, $`w`$ are arbitrary and $`u`$, $`w`$ real functions of $`\sigma `$. The subbranch of the studied metric with constant curvature arises form the above expression by setting $`w=(\lambda _{\text{cosm}}/6)u`$.
If the electromagnetic field is switched off, one arrives at the non–twisting type N solutions of Garcia et al. .
## III Plane–fronted gravitational and electromagnetic waves in MAG
In this section we generalize the type N gravitational and electromagnetic waves to the metric–affine gravity theories. We will present exact solutions of the field equations belonging to the Lagrangian
$$L=V_{\mathrm{MAG}}+V_{\mathrm{Max}},$$
(35)
where $`V_{\mathrm{Max}}=(1/2)F^{}F`$ is the Lagrangian of the Maxwell field and $`F=dA`$ is the electromagnetic field strength. The MAG Lagrangian considered here reads (a more general MAG Lagrangian can be found in ):
$`V_{\mathrm{MAG}}={\displaystyle \frac{1}{2\kappa }}`$ $`[a_0R^{\alpha \beta }\eta _{\alpha \beta }2\lambda _{\text{cosm}}\eta `$ (41)
$`+T^\alpha {}_{}{}^{}\left({\displaystyle \underset{I=1}{\overset{3}{}}}a_I^{(I)}T_\alpha \right)`$
$`+2\left({\displaystyle \underset{I=2}{\overset{4}{}}}c_I^{(I)}Q_{\alpha \beta }\right)\vartheta ^\alpha {}_{}{}^{}T_{}^{\beta }`$
$`+Q_{\alpha \beta }{}_{}{}^{}\left({\displaystyle \underset{I=1}{\overset{4}{}}}b_I^{(I)}Q^{\alpha \beta }\right)`$
$`+b_5(^{(3)}Q_{\alpha \gamma }\vartheta ^\alpha )^{}(^{(4)}Q^{\beta \gamma }\vartheta _\beta )]`$
$`{\displaystyle \frac{1}{2\rho }}R^{\alpha \beta }{}_{}{}^{}\left(z_4^{(4)}Z_{\alpha \beta }\right).`$
where
$$a_0,\mathrm{},a_3,b_1,\mathrm{},b_5,c_2,c_3,c_4,z_4,$$
(42)
are dimensionless coupling constants, $`\kappa `$ is the weak and $`\rho `$ the strong gravitational coupling constant. The cosmological constant is denoted by $`\lambda _{\text{cosm}}`$. The signature of spacetime is $`(+++)`$, the volume 4–form $`\eta :={}_{}{}^{}\mathrm{\hspace{0.17em}1}`$, the 2–form $`\eta _{\alpha \beta }:=^{}(\vartheta _\alpha \vartheta _\beta )`$.
The two MAG field equations for electromagnetic matter are given by
$`DH_\alpha E_\alpha `$ $`=`$ $`\mathrm{\Sigma }_\alpha ^{\mathrm{Max}},`$ (43)
$`DH_\beta ^\alpha E_\beta ^\alpha `$ $`=`$ $`0,`$ (44)
with $`\mathrm{\Sigma }_\alpha ^{\mathrm{Max}}`$ as defined in (13). It can be alternatively written as
$$\mathrm{\Sigma }_\alpha ^{\mathrm{Max}}=e_\alpha V_{\mathrm{Max}}+(e_\alpha F)H.$$
(45)
For the torsion and nonmetricity field configurations, we concentrate on the simplest non–trivial case with shear. According to its irreducible decomposition , the nonmetricity contains two covector pieces, namely the dilation piece
$${}_{}{}^{(4)}Q_{\alpha \beta }^{}=Qg_{\alpha \beta }$$
(46)
and the proper shear piece
$${}_{}{}^{(3)}Q_{\alpha \beta }^{}=\frac{4}{9}(\vartheta _{(\alpha }e_{\beta )}\mathrm{\Lambda }\frac{1}{4}g_{\alpha \beta }\mathrm{\Lambda }),\text{with}\mathrm{\Lambda }:=\vartheta ^\alpha e^\beta Q_{\alpha \beta }.$$
(47)
Accordingly, our ansatz for the nonmetricity reads
$$Q_{\alpha \beta }=^{(3)}Q_{\alpha \beta }+^{(4)}Q_{\alpha \beta }.$$
(48)
The torsion, in addition to its tensor piece, encompasses a covector and an axial covector piece. Let us choose only the covector piece as non–vanishing:
$$T^\alpha ={}_{}{}^{(2)}T_{}^{\alpha }=\frac{1}{3}\vartheta ^\alpha T,\text{with}T:=e_\alpha T^\alpha .$$
(49)
Thus we are left with the three non–trivial 1–forms $`Q`$, $`\mathrm{\Lambda }`$, and $`T`$. We shall assume that this triplet of 1–forms shares the spacetime symmetries, that is, its members are proportional to each other . Our ansatz for the nonmetricity is expected to require a nonvanishing post–Riemannian term quadratic in the segmental curvature. This is the term in (41) carrying the coupling constant $`z_4`$ (note that the enumeration of the constants stems from the general Lagrangian mentioned in ).
We assume the following so–called triplet ansatz for our three 1–forms in (48) and (49),
$$Q=k_0\omega ,\mathrm{\Lambda }=k_1\omega ,T=k_2\omega ,$$
(50)
where $`k_0`$, $`k_1`$, and $`k_2`$ are constants. The triplet ansatz (50) reduces the electrovacuum MAG field equations (43)–(44) to an effective Einstein–Proca–Maxwell system:
$`{\displaystyle \frac{a_0}{2}}\eta _{\alpha \beta \gamma }\stackrel{~}{R}^{\beta \gamma }+\lambda _{\mathrm{cosm}}\eta _\alpha `$ $`=`$ $`\kappa \left[\mathrm{\Sigma }_\alpha ^{(\omega )}+\mathrm{\Sigma }_\alpha ^{\mathrm{Max}}\right],`$ (51)
$`d{}_{}{}^{}d\omega +m^2{}_{}{}^{}\omega `$ $`=`$ $`0,`$ (52)
$`dF=0`$ , $`d{}_{}{}^{}F=0.`$ (53)
These are partial differential equations in terms of the coframe $`\vartheta ^\alpha `$, the triplet 1–form $`\omega `$, and the electromagnetic potential 1–form $`A`$; here the tilde $`\stackrel{~}{}`$ denotes again the Riemannian part of the curvature. The energy–momentum current of the triplet field $`\omega `$ reads
$`\mathrm{\Sigma }_\alpha ^{(\omega )}`$ $`:`$ $`={\displaystyle \frac{z_4k_0^2}{2\rho }}\left\{(e_\alpha d\omega \right){}_{}{}^{}d\omega (e_\alpha ^{}d\omega )d\omega `$ (55)
$`+m^2[(e_\alpha \omega ){}_{}{}^{}\omega +(e_\alpha ^{}\omega )\omega ]\};`$
the effective “mass” $`m`$ depends, additionally, on $`\kappa `$ and the strong gravitational coupling constant $`z_4/\rho `$, see .
Therefore, as mentioned above, in the framework of the triplet ansatz, the electrovacuum sector of MAG reduces to an effective Einstein–Proca–Maxwell system. Moreover, by setting $`m=0`$, the system acquires the following constraint among the coupling constants $`k_0`$, $`k_1`$, $`k_2`$ of the triplet ansatz (50) and the constants of the Lagrangian (41):
$`4b_4+{\displaystyle \frac{3}{2}}a_0+{\displaystyle \frac{k_1}{2k_0}}(b_5a_0)+{\displaystyle \frac{k_2}{k_0}}(c_4+a_0)=0.`$ (56)
The coframe we will consider is of the form (2), i.e., it is the same as in the general relativistic case. Note that we changed the name of the function $`H`$ in $`s`$ (cf.(5)) into $``$ in order to distinguish the general relativistic from the MAG case.
Now $``$, representing a combined gravitational MAG plane wave and an electromagnetic wave, has to fulfill the equation
$$_{,\zeta \overline{\zeta }}+\frac{\lambda _{\text{cosm}}}{3}p^2=\frac{2\kappa p}{q}\left[f\overline{f}+g\overline{g}\right],$$
(57)
where $`f=f(\zeta ,\sigma )`$ and $`g=g(\zeta ,\sigma )`$ are arbitrary functions of their arguments.
The general solution of this equation is given by $`_\mathrm{h}+_\mathrm{p}`$ with
$$_\mathrm{h}=\mathrm{\Phi }_{,\zeta }\frac{\lambda _{\text{cosm}}}{3}\frac{\overline{\zeta }}{p}\mathrm{\Phi }+\overline{\mathrm{\Phi }}_{,\overline{\zeta }}\frac{\lambda _{\text{cosm}}}{3}\frac{\zeta }{p}\overline{\mathrm{\Phi }},$$
(58)
and
$$_\mathrm{p}=M_{,\zeta }\frac{\lambda _{\text{cosm}}}{3}\frac{\overline{\zeta }}{p}M+\overline{M}_{,\overline{\zeta }}\frac{\lambda _{\text{cosm}}}{3}\frac{\zeta }{p}\overline{M}.$$
(59)
Here $`M=M(\sigma ,\zeta ,\overline{\zeta })`$ is a solution of the non–homogeneous equation for $``$, which is given by
$$M=\kappa ^{\overline{\zeta }}𝑑\overline{\zeta }p^2^\zeta \frac{d\zeta ^{}}{p^2}^\zeta ^{}𝑑\zeta ^{\prime \prime }\frac{p}{q}\left[f\overline{f}+g\overline{g}\right].$$
(60)
For given functions $`f`$ and $`g`$, one integrates (60) for $`M`$ and obtains $`_\mathrm{p}`$ from (59). The general solution is obtained by adding the homogeneous solution (58), where $`\mathrm{\Phi }`$ is an arbitrary holomorphic function of $`\zeta `$ and $`\sigma `$. The 1–form $`\omega `$ entering the triplet ansatz (50) is given by
$$\omega =\left[^\zeta g(\zeta ^{},\sigma )𝑑\zeta ^{}+^{\overline{\zeta }}\overline{g}(\overline{\zeta }^{},\sigma )𝑑\overline{\zeta }^{}\right]\vartheta ^{\widehat{2}},$$
(61)
where $`g=g(\zeta ,\sigma )`$ represents an arbitrary function of the coordinates. Moreover, the electromagnetic 2–form is given by
$$F=dA=d\left[\left(^\zeta f(\zeta ^{},\sigma )𝑑\zeta ^{}+^{\overline{\zeta }}\overline{f}(\overline{\zeta }^{},\sigma )𝑑\overline{\zeta }^{}\right)\vartheta ^{\widehat{2}}\right]$$
(62)
in terms of the arbitray function $`f=f(\zeta ,\sigma )`$. Inserting this ansatz into the field equations (51)-(53) yields the following additional constraints among the constants of (41):
$`a_0=1,z_4={\displaystyle \frac{\rho }{2k_0}}.`$ (63)
## IV Particular solutions
For better understanding, let us look for certain families of particular solutions of our dynamical system by integrating (57) restricted to $`\alpha =1`$ and $`\beta =0`$. Now the coframe in terms of $`p(\zeta ,\overline{\zeta }),q(\zeta ,\overline{\zeta })`$ and $`(\sigma ,\zeta ,\overline{\zeta })`$ reads
$`\vartheta ^{\widehat{0}}`$ $`=`$ $`{\displaystyle \frac{1}{p}}d\zeta ,\vartheta ^{\widehat{1}}={\displaystyle \frac{1}{p}}d\overline{\zeta },\vartheta ^{\widehat{2}}=d\sigma ,`$ (64)
$`\vartheta ^{\widehat{3}}`$ $`=`$ $`\left({\displaystyle \frac{q}{p}}\right)^2\left[\left({\displaystyle \frac{p}{2q}}(\sigma ,\zeta ,\overline{\zeta }){\displaystyle \frac{\lambda _{\mathrm{cosm}}}{6}}\rho ^2\right)d\sigma +d\rho \right].`$ (65)
Here $`p`$ and $`q`$ take the explicit form:
$$p(\zeta ,\overline{\zeta })=1+\frac{\lambda _{\mathrm{cosm}}}{6}\zeta \overline{\zeta },q(\zeta ,\overline{\zeta })=1\frac{\lambda _{\mathrm{cosm}}}{6}\zeta \overline{\zeta }.$$
(66)
Eq.(57) is a linear equation, therefore, one can look independently for solutions of the non–homogeneous equation for the $`f`$ exitations (associated with the electromagnetic field) and for the $`g`$ exitations (associated with the post–Riemannian pieces). Consequently, the addition of these solutions, corresponding to $`f`$ and $`g`$, will be again a solution. For simplicity, we shall restrict ourselves to the case where $`g(\zeta ,\sigma )`$ and $`f(\zeta ,\sigma )`$ are polynomial functions of $`\zeta `$ and $`\zeta ^1`$. Let us try the cases
$$f(\zeta ,\sigma )=f_0\zeta ^n,n=0,\pm 1,\pm 2,\pm 3,\mathrm{}.$$
(67)
Then one obtains the following branches of solutions for $`_\mathrm{p}`$:
(i) $`n<1`$
$`_\text{p}`$ $`=`$ $`{\displaystyle \frac{2\kappa pf_0^2}{q}}({\displaystyle \frac{(\zeta \overline{\zeta })^{1+n}}{(1+n)^2}}+4\left({\displaystyle \frac{\lambda _{\mathrm{cosm}}}{6}}\right)^{n1}\mathrm{ln}|q|4\left({\displaystyle \frac{\lambda _{\mathrm{cosm}}}{6}}\right)^{n1}\mathrm{ln}|p1|`$ (69)
$`+4{\displaystyle \underset{r=1}{\overset{n1}{}}}{\displaystyle \frac{\left(\frac{\lambda _{\mathrm{cosm}}}{6}\right)^{nr1}}{r\left(\zeta \overline{\zeta }\right)^r}})+{\displaystyle \frac{8\kappa f_0^2\left(\zeta \overline{\zeta }\right)^{n+1}}{(1+n)p}},`$
(ii) $`n=1`$
$$_\text{p}=\frac{2\kappa f_0^2}{p}\left(4q\mathrm{ln}\left|q\right|+\frac{2\lambda _{\mathrm{cosm}}\zeta \overline{\zeta }}{3}\mathrm{ln}\left(\zeta \overline{\zeta }\right)+\frac{q}{2}\mathrm{ln}^2\left(\zeta \overline{\zeta }\right)\right),$$
(70)
(iii) $`n>1`$
$$_\text{p}=\frac{8\kappa f_0^2q}{p}\left(\frac{\lambda _{\mathrm{cosm}}}{6}\right)^{n1}\left(\mathrm{ln}\left|q\right|+\underset{r=1}{\overset{n}{}}\frac{({}_{r}{}^{n})}{r}\left(\left(p2\right)^r(1)^r\right)\right)+\frac{2\kappa f_0^2(\zeta \overline{\zeta })^{n+1}}{p(n+1)^2}(4(n+1)+q).$$
(71)
Similarily one can proceed with solutions for $`g`$,
$$g(\zeta ,\sigma )=g_0\zeta ^l,l=0,\pm 1,\pm 2,\pm 3,\mathrm{}$$
(72)
The form of the different branches of $`_\mathrm{p}`$ do not change, but the substitution $`nl`$ and $`f_0g_0`$ should be performed. Therefore, one can obtain different branches of solutions by combining the $`f`$ branches with the $`g`$ branches of $`_\mathrm{p}`$.
For these particular classes one can choose $`_\mathrm{h}`$ as displayed in (58). Given $`g(\zeta ,\sigma )`$ and $`f(\zeta ,\sigma )`$ it is straightforward to evaluate the 1–form $`\omega `$ of (61) and the electromagnetic 2–form of (62).
This solution was checked by means of the computer algebra system Reduce by applying its Excalc package for treating exterior differential forms.
## V Discussion
We investigated plane–fronted electrovacuum–MAG waves with cosmological constant in the triplet ansatz sector of the theory. These waves carry curvature, nonmetricity, torsion, and an electromagnetic field. Apart from the cosmological constant, the solutions contain four wave parameters, given by the functions $`\alpha (\sigma )`$, $`\beta (\sigma )`$, $`\overline{\beta }(\sigma )`$ and $`_\sigma \mathrm{ln}|q(\sigma ,\zeta ,\overline{\zeta })|`$. Our plane–fronted wave solutions are given in terms of three arbitrary complex functions, i.e. $`\mathrm{\Phi }(\sigma ,\zeta )`$ associated with the Riemannian part, $`g(\sigma ,\zeta )`$ related to the non-Riemannian triplet, and $`f(\sigma ,\zeta )`$ corresponding to the Maxwell field. In this way, we generalize the plane–fronted electrovacuum Ozsvath–Robinson–Rozga waves. In brief, the solution reads:
ansatz for coframe $`\vartheta ^{\widehat{0}},\vartheta ^{\widehat{1}},\vartheta ^{\widehat{2}},\vartheta ^{\widehat{3}}`$ (2.2) arbitrary functions in coframe $`\alpha (\sigma ),\beta (\sigma )`$ MAG Lagrangian $`V_{\mathrm{MAG}}`$ and non-vanishing coupling constants (3.2),(3.3) with (3.16),(3.23) triplet ansatz for nonmetricity and torsion $`QT\mathrm{\Lambda }\omega `$, cf. (3.11) energy-momentum current of the Maxwell field (2.11) resp. (2.13) energy-momentum current of the triplet field (3.15) field equations (3.12)-(3.14) arbitrary function governing the vacuum solution $`_\mathrm{h}`$ $`\mathrm{\Phi }(\sigma ,\zeta )`$, cf. (3.18) arbitrary function in the electromagnetic 2-form $`F`$ $`f(\sigma ,\zeta )`$, cf. (3.22) arbitrary function in the triplet 1-form $`\omega `$ $`g(\sigma ,\zeta )`$, cf. (3.21) solution for the electromagnetic 2-form $`F`$ $`F`$ $`d(f𝑑\zeta +\overline{f}d\overline{\zeta })\vartheta ^{\widehat{2}}`$ solution for the triplet 1-form $`\omega `$ $`\omega (g𝑑\zeta +\overline{g}d\overline{\zeta })\vartheta ^{\widehat{2}}`$ solution for function $``$$`(\sigma ,\zeta ,\overline{\zeta })`$ entering coframe (3.18)-(3.20)
The final form of $`T^\alpha `$ and $`Q_{\alpha \beta }`$ in terms of $`g(\zeta ^{},\sigma )`$ reads,
$`T^\alpha `$ $`=`$ $`{\displaystyle \frac{k_2}{3}}\left[{\displaystyle ^\zeta }g(\zeta ^{},\sigma )𝑑\zeta ^{}+{\displaystyle ^{\overline{\zeta }}}\overline{g}(\overline{\zeta }^{},\sigma )𝑑\overline{\zeta }^{}\right]\vartheta ^\alpha \vartheta ^{\widehat{2}},`$ (73)
$`Q_{\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{4k_1}{9}}\vartheta _{(\alpha }e_{\beta )}[{\displaystyle ^\zeta }g(\zeta ^{},\sigma )d\zeta ^{}+{\displaystyle ^{\overline{\zeta }}}\overline{g}(\overline{\zeta }^{},\sigma )d\overline{\zeta }^{}]\vartheta ^{\widehat{2}}`$ (75)
$`+g_{\alpha \beta }\left({\displaystyle \frac{k_1}{9}}k_0\right)\left[{\displaystyle ^\zeta }g(\zeta ^{},\sigma )𝑑\zeta ^{}+{\displaystyle ^{\overline{\zeta }}}\overline{g}(\overline{\zeta }^{},\sigma )𝑑\overline{\zeta }^{}\right]\vartheta ^{\widehat{2}}.`$
The electromagnetic potential 1-form is given by
$$A=\left(^\zeta f(\zeta ^{},\sigma )𝑑\zeta ^{}+^{\overline{\zeta }}\overline{f}(\overline{\zeta }^{},\sigma )𝑑\overline{\zeta }^{}\right)\vartheta ^{\widehat{2}}.$$
(76)
It is straightforward to perform a detailed classification of the plane–fronted waves in MAG by carrying through a similar analysis as the one done by Sippel and Goenner . We leave this, however, for future work.
## VI Outlook
The theories of modern physics generally involve a mathematical model, defined by a certain set of differential equations, and supplemented by a set of rules for translating the mathematical results into meaningful statements about the physical world. In the case of gravity theories, because they deal with the most universal of the physical interactions, one has an additional class of problems concerning the influence of the gravitational field on other fields and matter. These are often studied by working within a fixed gravitational field, usually an exact solution . In this context our plane–fronted waves solutions contribute to enhance our understanding of some of these questions in the framework of MAG theories, in particular the ones concerned with the gravitational radiation.
Gravitational waves have traveled almost unimpeded through the universe since they were generated at times as early as $`10^{24}`$ sec. after the big bang. This radiation carries information that no electromagnetic radiation can give to us because the electromagnetic radiation is scattered countless times by the dense material surrounding the explosion, losing in the process most of the detailed information it might carry about the explosion. Beyond this, we can be virtually certain that gravitational wave spectrum has surprises for us, clues to phenomena we never suspected. Therefore, it is not surprising, that considerable effort is nowadays being devoted to the development of sufficiently sensitive gravitational wave antennas. Moreover, observing them would provide important constraints on theories of inflation and high–energy physics.
Even though Einstein’s treatment of spacetime as a Riemannian manifold appears almost fully corroborated experimentally, there are several reasons to believe that the validity of such a description is limited to macroscopic structures and to the present cosmological era. The only available finite perturbative treatment of quantum gravity, namely the theory of the quantum superstring , suggests that non–Riemannian features are present on the scale of the Planck length. On the other hand, recent advances in the study of the early universe, as represented by the inflationary model, involve, in addition to the metric tensor, at the very least a scalar dilaton induced by a Weyl geometry, i.e., again an essential departure from Riemannian metricity . Even at the classical cosmological level, a dilatonic field has recently been used to describe the presence of dark matter in the universe, as well as to explain certain cosmological observations which contradicted the fundaments of the standard cosmological model .
Inflation is an attractive scenario for the early universe because it makes the large scale homogeneity of the universe easy to understand. It also provides a mechanism for producing initial density perturbations large enough to evolve into galaxies as the universe expands. These perturbations are accompanied by perturbation of the gravitational field that travel through the universe, redshifting in the same way that photons do. The perturbations arise by parametric amplification of quantum fluctuations in the gravitational wave field that existed before the inflation began. The huge expansion associated with inflation puts energy into these fluctuations, converting them into real gravitational waves with classical amplitudes. Even if inflation did not occur, the perturbations that lead to galaxies must have arisen in some other way, and it is possible that this alternative mechanism also produced gravitational waves.
It is worthwhile to stress the fact that we do not believe that at the present state of the universe the geometry of spacetime is described by a metric–affine one. We rather think, and there is good experimental evidence, that the present–day geometry is metric–compatible, i.e., its nonmetricity vanishes. In earlier epochs of the universe, however, when the energies of the cosmic “fluid” were much higher than today, we expect scale invariance to prevail — and the canonical dilation or scale current of matter, the trace of the hypermomentum current $`\mathrm{\Delta }^\gamma _\gamma `$, is coupled, according to MAG, to the Weyl covector $`Q^\gamma _\gamma `$. By the same token, shear type excitations of the material multispinors (Regge trajectory type of constructs) are expected to arise, thereby liberating the (metric-compatible) Riemann–Cartan spacetime from its constraint of vanishing nonmetricity $`Q_{\alpha \beta }=0`$ . Tresguerres has proposed a simple cosmological model of Friedmann type which carries a metric-affine geometry at the beginning of the universe, the nonmetricity of which dies out exponentially in time. That is the kind of thing we expect.
In full, exact solutions of the type obtained may serve well as starting point for the upcoming analysis of gravitational wave astronomy data. In this sense it might contribute to our understanding of light and gravitational wave propagation in early stages of the universe. Moreover, plane wave solutions contribute to resolve some of the controversies about the existence of such gravitational radiation.
###### Acknowledgements.
We thank Friedrich W. Hehl for useful discussions and literature hints. This research was supported by CONACyT Grants: 28339E, and 32138E.
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# Class numbers of orders in cubic fields
## 1 The main theorem
Let $`𝒪`$ be an order in a number field $`F`$. Let $`I(𝒪)`$ be the set of all finitely generated $`𝒪`$-submodules of $`F`$. According to the Jordan-Zassenhaus Theorem , the set of isomorphism classes $`[I(𝒪)]`$ of elements of $`I(𝒪)`$ is finite. Let $`h(𝒪)`$ be the cardinality of the set $`[I(𝒪)]`$, called the class number of $`𝒪`$.
The multiplicative group of invertible elements $`F^\times `$ of $`F`$ acts on $`I(𝒪)`$ by multiplication: $`\lambda .M=\lambda M=M\lambda `$. Since $`M`$ and $`M\lambda `$ are isomorphic as $`𝒪`$-modules we get a map $`I(𝒪)/F^\times [I(𝒪)]`$ mapping $`MF^\times `$ to its class $`[M]`$. We claim that this map is a bijection. It is clearly surjective. So let $`M`$ and $`N`$ be elements in $`I(𝒪)`$ which are isomorphic. Fix an isomorphism $`T:MN`$. Then $`T`$ extends to an $`F`$-isomorphism $`T_F:F=FMFN=F`$, so there is $`\alpha F^\times `$ such that $`T_F`$ is just multiplication by $`\alpha `$. This implies the claim. It follows that the class number $`h(𝒪)`$ equals the cardinality of $`I(𝒪)/F^\times `$.
The class number of the maximal order is also called the class number of the field. In general the class number of $`𝒪`$ will be larger than that of $`F`$.
A number field $`F`$ of degree $`3`$ over the rationals is also called a cubic field. Consider the set of embeddings of a given cubic field into $``$. This set has three elements and is permuted by complex conjugation. It follows that either all three embeddings are fixed by complex conjugation, i.e. they are real in which case the field is called totally real or two of them are swapped by complex conjugation and the third is fixed, i.e. one is real and the other two are complex, in which case the field is called complex.
Let $`F`$ be a complex cubic field. Since the order of the automorphism group of $`F`$ divides the degree, which is $`3`$, it is either $`1`$ or $`3`$. In the latter case the extension would be galois, so the Galois group would act transitively on the set of embeddings of $`F`$ into the complex numbers. In particular, all embeddings would be either real or complex. Since this is not the case it follows that the automorphism group of a complex cubic field is trivial.
Let $`F`$ be a complex cubic field and let $`𝒪`$ be an order of $`F`$. Then the group of units satisfies
$$𝒪^\times =\pm ϵ^{}$$
for some base unit $`ϵ`$. The image of the base unit under the real embedding will be of the form $`\pm e^{\pm R(𝒪)}`$, where $`R(𝒪)`$ is the regulator of $`𝒪`$ (see ). Set
$$r(𝒪)=e^{3R(𝒪)}.$$
A prime number $`p`$ is called non-decomposed in $`F`$ if there is only one place in $`F`$ lying above $`p`$. Fix a finite set $`S`$ of prime numbers with at least two elements and let $`C(S)`$ be the set of all complex cubic fields $`F`$ such that all primes $`p`$ in $`S`$ are non-decomposed in $`F`$. For $`FC(S)`$ let $`O_F(S)`$ be the set of all orders $`𝒪`$ in $`F`$ which are maximal at all $`pS`$, i.e. are such that the completion $`𝒪_p=𝒪_{}_p`$ is the maximal order of the field $`F_p=F_{}_p`$ for all $`pS`$. Let $`O(S)`$ be the union of all $`O_F(S)`$ where $`F`$ ranges over $`C(S)`$.
Let $`S_i(F)`$ be the set of all $`pS`$ such that $`p`$ is inert in $`F`$. Define
$`\lambda _S(F)`$ $`=`$ $`d^{|S_i(F)|}`$
$`=`$ $`{\displaystyle \underset{pS}{}}f_p(F),`$
where $`f_p(F)`$ is the inertia degree of $`p`$ in $`F`$. For an order $`𝒪`$ in $`F`$ let
$$\lambda _S(𝒪)=\lambda _S(F).$$
The following is our main theorem.
###### Theorem 1.1
For $`x>0`$ let
$$\pi _S(x)=\underset{_{r(𝒪)x}^{𝒪O(S)}}{}h(𝒪)\lambda _S(𝒪).$$
Then, as $`x\mathrm{}`$ we have
$$\pi _S(x)\frac{x}{\mathrm{log}x}.$$
More sharply,
$$\pi _S(x)=li(x)+O\left(\frac{x^{\frac{3}{4}}}{\mathrm{log}x}\right),$$
as $`x\mathrm{}`$ where $`li(x)=_2^{\mathrm{}}\frac{1}{\mathrm{log}t}𝑑t`$ is the integral logarithm as in the usual prime number theorem.
As an immediate consequence we get the following corollary.
###### Corollary 1.4
Let
$$\stackrel{~}{\pi }_S(x)=\underset{\genfrac{}{}{0.0pt}{}{𝒪O(S)}{r(𝒪)x}}{}h(𝒪),$$
then
$$\underset{x\mathrm{}}{lim\; sup}\frac{\stackrel{~}{\pi }_S\mathrm{log}x}{x}1,$$
and
$$\underset{x\mathrm{}}{lim\; inf}\frac{\stackrel{~}{\pi }_S\mathrm{log}x}{x}\frac{1}{3^{|S|}}.$$
Proof: For any $`𝒪`$ we have $`1\lambda _S(𝒪)3^{|S|}`$, which implies the corollary.
$`\frac{}{}`$ Q.E.D.
The theorem will be proved in the following sections.
To clarify the range of the theorem we note:
###### Proposition 1.7
For every complex cubic field $`F`$ there are infinitely many primes $`p`$ which are non-decomposed in $`F`$. Every order $`𝒪`$ in $`F`$ is maximal at almost all primes $`p`$.
Proof: The second assertion is well known and so it simply remains to prove the first. Let $`F`$ be a complex cubic field. We will show that there are infinitely many primes $`p`$ which are non-decomposed in $`F`$. Let $`E/`$ be the galois hull of $`F`$. Then $`\mathrm{Gal}(E/)`$ is the permutation group $`S_3`$ in three letters. Let $`\sigma `$ be the generator of $`\mathrm{Gal}(E/F)/2`$ and let $`p`$ be a prime which is unramified in $`E`$ and such that there is a prime $`𝔭_1`$ of $`E`$ over $`p`$ with Frobenius $`\tau `$ of order $`3`$. By the Tchebotarev density theorem there are infinitely many such $`p`$. Then the stabilizer of $`𝔭_1`$ in $`S_3`$ is $`\tau `$ of order $`3`$, so $`p𝒪_E=𝔭_1𝔭_2`$, say. As $`\sigma `$ does not stabilize $`𝔭_1`$, it interchanges $`𝔭_1`$ and $`𝔭_2`$, so $`𝔭_1F=𝔭_2F`$ is a prime in $`𝒪_F`$, i.e. $`p`$ is non-decomposed in $`F`$.
$`\frac{}{}`$ Q.E.D.
## 2 Division algebras of prime degree
Let $`d`$ be a prime $`>2`$. Let $`M()`$ denote a division algebra of degree $`d`$ over $``$. The dimension of $`M()`$ is $`d^2`$. Fix a maximal order $`M()M()`$; for any commutative ring with unit $`R`$ we will write $`M(R)=M()_{}R`$. Further, let $`M(R)^\times `$ denote the multiplicative group of invertible elements in $`M(R)`$. We say that $`M()`$ splits over a prime $`p`$ if $`M(_p)`$ is not a division algebra. Since $`d`$ is a prime we then get an isomorphism $`M(_p)\mathrm{Mat}_d(_p)`$ (see ). Further, since there are no associative division algebras over the reals of degree $`d`$ it follows that $`M()\mathrm{Mat}_d()`$. For any ring $`R`$ the reduced norm induces a map $`\mathrm{det}:M(R)R`$. Let $`𝒢(R)=\{xM(R)|\mathrm{det}(x)=1\}`$. Then $`𝒢`$ is a simple linear algebraic group over $``$. Let $`G=𝒢()`$; then $`G`$ is isomorphic to $`SL_d()`$. Let $`S`$ be the set of places of $``$ where $`M()`$ does not split. The set $`S`$ always is finite, has at least two elements, and contains only finite places. For any prime $`p`$ we have that $`M(_p)`$ is a maximal $`_p`$-order of $`M(_p)`$.
Let $`F/`$ be a finite field extension which embeds into $`M()`$. By the Skolem-Noether Theorem (, Thm 7.21), any two embeddings $`\sigma _1,\sigma _2:FM()`$ are conjugate by $`M()^\times `$, i.e. there is a $`uM()^\times `$ such that $`\sigma _2(x)=u\sigma _1(x)u^1`$ for any $`xF`$.
A prime $`p`$ is called non-decomposed in the extension $`F/`$ if there is only one place $`v`$ of $`F`$ lying above $`p`$.
###### Lemma 2.10
Let $`F/`$ be a nontrivial finite field extension. Then $`F`$ embeds into $`M()`$ if and only if $`[F:]=d`$ and every prime $`p`$ in $`S`$ is non-decomposed in $`F`$.
Proof: Assume that $`F`$ embeds into $`M()`$; then $`[F:]`$ divides $`d`$, which is a prime, so that $`[F:]=d`$. It follows that $`F`$ is a maximal subfield of $`M()`$. By Proposition 13.3 in it follows that $`F`$ is a splitting field of $`M()`$. Thus, by Theorem 32.15 of it follows that for every $`pS`$ and any place $`v`$ of $`F`$ over $`p`$ we have $`[F_v:_p]=d`$, which by the degree formula implies that there is only one $`v`$ over $`p`$. Conversely, assume that there is only one prime above each $`pS`$. Let $`pS`$ and let $`e`$ be the ramification index of $`p`$ in $`F`$ and $`f`$ the inertia degree. Then $`d=[F:]=ef=[F_v:Q_p]`$ by the global, respectively local, degree formula. Since this holds for any $`pS`$, , Theorem 32.15 implies that $`F`$ is a splitting field for $`M()`$. By Proposition 13.3 of we infer that $`F`$ embeds into $`M()`$.
$`\frac{}{}`$ Q.E.D.
Let $`F/`$ be a field extension of degree $`d`$ which embeds into $`M()`$. Then for any embedding $`\sigma :FM()`$ the set
$$𝒪_\sigma =\sigma ^1(\sigma (F)M())$$
is an order of $`F`$. For $`pS`$ let $`v_p`$ denote the unique place of $`F`$ over $`p`$.
###### Lemma 2.11
Let $`\sigma :FM()`$ be an embedding of the field $`F`$. For any $`pS`$ the completion $`𝒪_{\sigma ,v_p}`$ is a maximal order in $`F_{v_p}`$. Conversely, let $`𝒪F`$ be an order such that for any $`pS`$ the completion $`𝒪_{v_p}`$ is maximal. Then there is an embedding $`\sigma :FM()`$, such that $`𝒪=𝒪_\sigma `$.
Proof: Let $`pS`$; then $`M(_p)`$ is a maximal $`_p`$-order in the division algebra $`M(_p)`$, and hence by Theorem 12.8 of it coincides with the integral closure of $`_p`$ in $`M(_p)`$. Therefore $`𝒪_{\sigma ,v_p}=\sigma ^1(\sigma (F_{v_p}M(_p))`$ is the integral closure of $`_p`$ in $`F_{v_p}`$, which is the maximal order of $`F_{v_p}`$.
For the converse, let $`𝒪`$ be an order of $`F`$ such that the completion $`𝒪_{v_p}`$ is maximal for each $`pS`$. Fix an embedding $`FM()`$ and consider $`F`$ as a subfield of $`M()`$. For any $`uM()^\times `$ let $`𝒪_u=Fu^1M()u`$. We will show that there is a $`uM()^\times `$ such that $`𝒪=𝒪_u`$. This will prove the proposition since one can then take $`\sigma `$ to be the conjugation by $`u`$.
Let $`𝒪_1=FM()`$. Since $`𝒪`$ and $`𝒪_1`$ are orders, they both are maximal at all but finitely many places. So there is a finite set of primes $`T`$ with $`TS=\mathrm{}`$ and such that with $`T_F`$ denoting the set of places of $`F`$ lying over $`T`$ we have that for any place $`v`$ of $`F`$ with $`vT_F`$ the completion $`𝒪_v`$ is maximal and equals $`𝒪_{1,v}`$. Let $`pT`$ and fix an isomorphism $`M(_p)\mathrm{Mat}_d(_p)`$. Fix a $`_p`$-basis of $`𝒪_{}_p`$. This basis then induces an embedding $`\sigma _p:F_{}_p\mathrm{Mat}_d(_p)=M(_p)`$, such that
$$\sigma _p^1(\sigma _p(F_p)\mathrm{Mat}_d(_p))=𝒪_p.$$
By the Noether-Skolem Theorem 7.21 there is a $`\stackrel{~}{u}_pM(_p^\times )`$ such that $`𝒪_p=F_p\stackrel{~}{u}_p^1M(_p)\stackrel{~}{u}_p`$. For $`pT`$ set $`\stackrel{~}{u}_p=1`$ and let $`\stackrel{~}{u}=(\stackrel{~}{u}_p)M(𝔸_{fin})`$, where $`𝔸_{fin}`$ is the ring of the finite adeles over $``$, i.e. the restricted product over the local fields $`_p`$, where $`p`$ ranges over the primes. By strong approximation $`M()^\times `$ is dense in $`M(𝔸_{fin})`$ and so there is $`uM()^\times `$ such that $`M(\widehat{})u=M(\widehat{})\stackrel{~}{u}`$, where $`\widehat{}=_p_p`$. It follows that $`𝒪=𝒪_u`$ for this $`u`$.
$`\frac{}{}`$ Q.E.D.
Let $`F`$ be a field extension of $``$ of degree $`d`$ which embeds into $`M()`$. Let $`𝒪`$ be an order of $`F`$, which is maximal at each place in $`S`$. By Lemma 2.11 we know that there is an embedding $`\sigma `$ of $`F`$ into $`M()`$ such that $`𝒪=𝒪_\sigma `$. Let $`uM()^\times `$ and let $`{}_{}{}^{u}\sigma =u\sigma u^1`$. Then $`𝒪_{{}_{}{}^{u}\sigma }=𝒪_\sigma `$, so the group $`M()^\times `$ acts on the set $`\mathrm{\Sigma }(𝒪)`$ of all $`\sigma `$ with $`𝒪=𝒪_\sigma `$.
###### Lemma 2.12
The quotient $`\mathrm{\Sigma }(𝒪)/M()^\times `$ is finite and has cardinality equal to the product $`h(𝒪)\lambda _S(𝒪)`$.
Proof: Fix an embedding $`FM()`$ and consider $`F`$ as a subfield of $`M()`$ such that $`𝒪=FM()`$. For $`uM()^\times `$ let
$$𝒪_u=Fu^1M()u.$$
Let $`U`$ be the set of all $`uM()^\times `$ such that
$$𝒪=FM()=Fu^1M()u.$$
Then $`F^\times `$ acts on $`U`$ by multiplication from the right and $`M()^\times `$ acts by multiplication from the left. It is clear that
$$\left|M()^\times \backslash U/F^\times \right|=\left|M()^\times \backslash \mathrm{\Sigma }(𝒪)\right|.$$
So we only have to show that the left hand side equals $`h(𝒪)\lambda _S(𝒪)`$. For $`uU`$ let
$$I_u=FM()u.$$
Then $`I_u`$ is a finitely generated $`𝒪`$-module in $`F`$. We claim that the map $`\mathrm{\Psi }:`$
$`M()^\times \backslash U/F^\times `$ $``$ $`I(𝒪)/F^\times `$
$`u`$ $``$ $`I_u,`$
is surjective and $`\lambda _S(𝒪)`$ to one. We will show this through localization and strong approximation. So, for a prime $`p`$ let $`U_p`$ be the set of all $`u_pM(_p)^\times `$ such that $`𝒪_p=F_pM(_p)=F_pu_p^1M(_p)u_p`$. We have to show the following:
* For $`pS`$ the localized map $`\mathrm{\Psi }_p:M(_p)^\times \backslash U_p/F_p^\times I(𝒪_p)/F_p^\times `$ is injective.
* For $`pS`$ the map $`\mathrm{\Psi }_p`$ is $`f_p(F)`$ to one.
* The map $`\mathrm{\Psi }`$ is surjective.
For ‘1.’ let $`pS`$, let $`u_p,v_pU_p`$, and assume
$$F_pM(_p)u_p=F_pM(_p)v_p.$$
Let $`z_p=v_pu_p^1`$. Elementary divisor theory implies that there are $`x,yM(_p)^\times =Mat_d(_p)^\times `$ such that
$$z_p=x\mathrm{diag}(p^{k_1},\mathrm{},p^{k_d})y,$$
where diag denotes the diagonal matrix and $`k_1k_2\mathrm{}k_d`$ are integers. Replacing $`u_p`$ by $`yu_p`$ and $`v_p`$ by $`x^1v_p`$ we may assume that $`z`$ equals the diagonal matrix. The assumptions then easily imply $`k_1=\mathrm{}=k_d=0`$, which gives the first claim.
For ‘2.’ let $`pS`$ and recall that $`F_p`$ is a local field and so $`h(𝒪_p)=1`$. So the claim is equivalent to
$$\left|M(_p)^\times \backslash M(_p)^\times /F_p^\times \right|=f_p(F).$$
By Proposition 17.7 of it follows that
$$\left|M(_p)^\times \backslash M(_p)^\times \right|=d=e(M(_p)/_p),$$
where $`e`$ denotes the ramification index. If $`F_p/Q_p`$ is ramified, then $`f_p(F)=1`$ and $`\left|M(_p)^\times \backslash M(_p)^\times /F_p^\times \right|=1`$. If $`F_p/_p`$ is unramified , then $`f_p(F)=d`$ and $`\left|M(_p)^\times \backslash M(_p)^\times /F_p^\times \right|=d`$ as claimed.
For the surjectivity of $`\mathrm{\Psi }`$ let $`I𝒪`$ be an ideal. We shall show that there is a $`uM()^\times `$ such that
$$Fu^1M()u=FM()$$
and
$$I=I_u=FM()u.$$
We shall do this locally. First note that, since $`I`$ is finitely generated, there is a finite set of primes $`T`$ with $`TS=\mathrm{}`$ such that for any $`pTS`$ the completion $`I_p`$ equals $`𝒪_p`$ which is the maximal order of $`F_p`$. For these $`p`$ set $`\stackrel{~}{u}_p=1`$.
Next let $`pS`$. Let $`v_p`$ be the unique place of $`F`$ over $`p`$. Then $`𝒪_p=𝒪_{v_p}`$ is maximal, so it is the valuation ring to $`v_p`$ and $`I_p=\pi _p^k𝒪_p`$ for some $`k0`$, where $`\pi _p`$ is a uniformizing element in $`𝒪_p`$. It follows that at this $`p`$, the element $`\stackrel{~}{u}_p=\pi _p^kId`$ will do the job.
Next let $`pT`$. Then $`M(_p)=\mathrm{Mat}_d(_p)`$. Let $`\overline{𝒪_p}=𝒪_p/p𝒪_p`$ and $`\overline{I_p}=I_p/pI_p`$. Then $`\overline{𝒪_p}`$ is a commutative algebra over the field $`𝔽_p`$ with $`p`$ elements, which implies that $`\overline{𝒪_p}_{i=1}^sF_i`$, where each $`F_i`$ is a finite field extension of $`𝔽_p`$. Let $`n_i`$ be the degree of $`F_i/𝔽_p`$. Then there is an embedding $`\overline{𝒪_p}\mathrm{Mat}_d(𝔽_p)`$ whose image lies in $`\mathrm{Mat}_{n_1}(𝔽_p)\times \mathrm{}\times \mathrm{Mat}_{n_s}(𝔽_p)\mathrm{Mat}_d(𝔽_p)`$. According to the Noether-Skolem Theorem there is a $`\overline{S}GL_d(𝔽_p)`$ such that $`\overline{S}\overline{𝒪_p}\overline{S}^1\mathrm{Mat}_{n_1}(𝔽_p)\times \mathrm{}\times \mathrm{Mat}_{n_s}(𝔽_p)\mathrm{Mat}_d(𝔽_p)`$. The $`\overline{𝒪_p}`$-ideal $`\overline{I_p}`$ must be of the form
$$\overline{I_p}=\underset{i1}{\overset{s}{}}ϵ_iF_i,$$
where $`ϵ_i\{0,1\}`$. Let $`S`$ be a matrix in $`GL_d(_p)`$ which reduces to $`\overline{S}`$ modulo $`p`$ and let $`\stackrel{~}{u}_p=S^1(p^{ϵ_1}Id_{n_1}\times \mathrm{}\times p^{ϵ_s}Id_{n_s})S`$ in $`\mathrm{Mat}_d(_p)`$. By abuse of notation we also write $`\stackrel{~}{u}_p`$ for its reduction modulo $`p`$. Then we have
$$\overline{I_p}=\overline{𝒪_p}\mathrm{Mat}_d(𝔽_p)\stackrel{~}{u}_p.$$
Let
$$I_{\stackrel{~}{u}_p}=FM(_p)\stackrel{~}{u}_p.$$
Then it follows that
$$\overline{I_p}\overline{I_{\stackrel{~}{u}_p}}=I_{\stackrel{~}{u}_p}/pI_{\stackrel{~}{u}_p}$$
and by Theorem 18.6 of it follows that $`I_pI_{\stackrel{~}{u}_p}`$, which implies that there is some $`\lambda F_p`$ such that $`I_p=I_{\stackrel{~}{u}_p}\lambda `$. Replacing $`\stackrel{~}{u}_p`$ by $`\stackrel{~}{u}_p\lambda `$ and setting $`\stackrel{~}{u}=(\stackrel{~}{u}_p)_pM(𝔸_{fin})`$ we get
$$I=FM()\stackrel{~}{u}.$$
By strong approximation there is a $`uM()^\times `$ such that $`M(\widehat{})u=M(\widehat{})\stackrel{~}{u}`$ and therefore $`I=I_u`$.
$`\frac{}{}`$ Q.E.D.
We will summarize the results of this section in the following proposition.
###### Proposition 2.17
Let $`d`$ be a prime $`>2`$ and let $`F/`$ be an extension of degree $`d`$. Then $`F`$ embeds into the division algebra $`M()`$ of degree $`d`$ if and only if every prime $`p`$ at which $`M()`$ does not split is non-decomposed in $`F`$. Every embedding $`\sigma :FM()`$ gives by intersection with $`M()`$ an order $`𝒪_\sigma `$ in $`F`$. Every order $`𝒪`$ of $`F`$, which is maximal at each $`p`$ where $`M()`$ is non-split, occurs in this way. The number of $`M()^\times `$-conjugacy classes of embeddings giving rise to the same order $`𝒪`$ is equal to $`h(𝒪)\lambda _S(𝒪)`$.
## 3 Analysis of the Ruelle zeta function
From now on we restrict to the case $`d=3`$. Let $`\mathrm{\Gamma }=𝒢()`$; then $`\mathrm{\Gamma }`$ forms a discrete subgroup of $`G=𝒢()SL_3()`$. Since $`M()`$ is a division algebra it follows that $`𝒢`$ is anisotropic over $``$ and so $`\mathrm{\Gamma }`$ is cocompact in $`G`$. Let $`𝔤_0`$ be the real Lie algebra of $`G`$ and let $`𝔤=𝔤_0_{}`$ be its complexification. A subgroup $`\mathrm{\Sigma }`$ of $`G`$ is called weakly neat if it is torsion free and for each $`\sigma \mathrm{\Sigma }`$ the adjoint $`Ad(\sigma )GL(𝔤)`$ does not have a root of unity except $`1`$ as an eigenvalue. In other words, a torsion free group $`\mathrm{\Sigma }`$ is weakly neat if for every $`\sigma \mathrm{\Sigma }`$ and every $`n`$ the connected component $`G_\sigma ^0`$ of the centralizer of $`\sigma `$ coincides with the connected component $`G_{\sigma ^n}^0`$ of the centralizer of $`\sigma ^n`$. Further an element $`x`$ of $`G`$ is called regular if its centralizer is a torus.
###### Lemma 3.16
The group $`\mathrm{\Gamma }`$ is torsion-free and every $`\gamma \mathrm{\Gamma }`$ with $`\gamma 1`$ is regular. In particular, it follows that $`\mathrm{\Gamma }`$ is weakly neat.
Proof: Let $`\gamma \mathrm{\Gamma }`$, $`\gamma 1`$; then the centralizer $`M()_\gamma `$ of $`\gamma `$ in $`M()`$ is a division subalgebra of $`M()`$ and so it is either $``$, a cubic number field or $`M()`$. It cannot be $``$ since it contains $`\gamma `$ and if $`\gamma `$ is in $``$, it is central so its centralizer is $`M()`$. It can neither be $`M()`$, since then $`\gamma `$ would be in $`1`$, say $`\gamma =r1`$ and then $`1=\mathrm{det}(\gamma )=r^3`$, thus $`\gamma =1`$ which was excluded. So it follows that $`M()_\gamma `$ is a cubic field $`F`$. Note that this implies that $`\gamma `$ cannot be a root of unity, since a cubic field does not contain roots of unity other than $`\pm 1`$. This shows that $`\mathrm{\Gamma }`$ is torsion free. Moreover, we get $`M()_\gamma =M()_\gamma _{}`$ so $`G_\gamma =(F)^1`$, the norm one elements, this is a torus, so $`\gamma `$ is regular.
$`\frac{}{}`$ Q.E.D.
We have $`𝔤_0=sl_3()`$ and $`𝔤=sl_3()`$. Let $`b`$ be the Killing form on $`𝔤`$. Then
$$b(X,Y)=\mathrm{tr}\mathrm{ad}(X)\mathrm{ad}(Y)=6\mathrm{tr}(XY).$$
Let $`KG`$ be the maximal compact subgroup $`SO(3)`$. Let $`𝔨_0𝔤_0`$ be its Lie algebra and let $`𝔭_0𝔤_0`$ be the orthogonal space of $`𝔨_0`$ with respect to the form $`b`$. Then $`b`$ is positive definite on $`𝔭_0`$ and thus defines a $`G`$-invariant metric on $`X=G/K`$, the symmetric space attached to $`G`$.
Since $`\mathrm{\Gamma }`$ is torsion-free it acts without fixed points on the contractible space $`X`$, so $`X_\mathrm{\Gamma }=\mathrm{\Gamma }\backslash X`$ is the classifying space of $`\mathrm{\Gamma }`$, in particular, it follows that $`\mathrm{\Gamma }`$ is the fundamental group of $`X_\mathrm{\Gamma }`$. We thus obtain a natural bijection
$$\left\{\begin{array}{c}\mathrm{free}\mathrm{homotopy}\mathrm{classes}\\ \mathrm{of}\mathrm{maps}S^1X_\mathrm{\Gamma }\end{array}\right\}\left\{\begin{array}{c}\mathrm{conjugacy}\mathrm{classes}\\ [\gamma ]\mathrm{in}\mathrm{\Gamma }\end{array}\right\}$$
For each conjugacy class $`[\gamma ]`$ let $`X_\gamma `$ denote the union of all geodesics lying in the homotopy class $`[\gamma ]`$ of $`\gamma `$. It is known that all geodesics in $`[\gamma ]`$ have the same length $`l_\gamma `$ and $`X_\gamma `$ is a submanifold of $`X_\mathrm{\Gamma }`$ diffeomorphic to $`\mathrm{\Gamma }_\gamma \backslash G_\gamma /K_\gamma `$, where $`G_\gamma `$ and $`\mathrm{\Gamma }_\gamma `$ are the centralizers of $`\gamma `$ and $`K_\gamma `$ is a maximal compact subgroup of $`G_\gamma `$.
An element $`\gamma \mathrm{\Gamma }`$ is called primitive if for $`\sigma \mathrm{\Gamma }`$ and $`n`$ the equation $`\sigma ^n=\gamma `$ implies that $`n=1`$. Since every closed geodesic is a positive power of a unique primitive one it is easy to see that every $`\gamma \mathrm{\Gamma }`$ with $`\gamma 1`$ is a positive power of a unique primitive element $`\gamma _0`$. We write $`\gamma =\gamma _0^{\mu (\gamma )}`$ and call $`\mu (\gamma )`$ the multiplicity of $`\gamma `$. Clearly primitivity is a property of the conjugacy class.
Up to conjugacy the group $`G`$ has two Cartan subgroups, namely the group of diagonal matrices and the group $`H=AB`$, where
$$A=\left\{\left(\begin{array}{ccc}a& & \\ & a& \\ & & a^2\end{array}\right)\right|a>0\}$$
and
$$B=\left(\begin{array}{cc}SO(2)& \\ & 1\end{array}\right).$$
Let $`P`$ denote the parabolic $`\left(\begin{array}{cc}& \\ \mathrm{0\; 0}& \end{array}\right)`$. It has a Langlands decomposition $`P=MAN`$ and $`B`$ is a compact Cartan subgroup of
$$MSL_2^\pm ()=\{xMat_2()|\mathrm{det}(x)=\pm 1\}.$$
Let
$$H_1=\frac{1}{6}\left(\begin{array}{ccc}1& & \\ & 1& \\ & & 2\end{array}\right)𝔞_0=\mathrm{Lie}A.$$
Then it follows that $`b(H_1)=b(H_1,H_1)=1`$. Let $`A^{}=\{\mathrm{exp}(tH_1)|t>0\}`$ and let $`_P(\mathrm{\Gamma })`$ be the set of all conjugacy classes $`[\gamma ]`$ in $`\mathrm{\Gamma }`$ such that $`\gamma `$ is conjugate in $`G`$ to an element $`a_\gamma b_\gamma `$ of $`A^{}B`$, and let $`_P^p(\mathrm{\Gamma })`$ the set of all primitive elements therein. For $`s`$ with $`\mathrm{Re}(s)>>0`$ let
$$R_\mathrm{\Gamma }(s)=\underset{[\gamma ]_P^p(\mathrm{\Gamma })}{}(1e^{sl_\gamma })$$
be the Ruelle-zeta function attached to $`P`$. In it is shown that $`R_\mathrm{\Gamma }(s)`$ converges for $`\mathrm{Re}(s)>>0`$ and that it extends to a meromorphic function of finite order on the plane. We will show:
###### Theorem 3.23
The function $`R_\mathrm{\Gamma }(s)`$ has a simple zero at $`s=1`$. Apart from that, all poles and zeros of $`R_\mathrm{\Gamma }(s)`$ are contained in the union of the interval $`[1,\frac{3}{4}]`$ with the three vertical lines given by $`\frac{1}{2}+i`$, $`i`$ and $`\frac{1}{2}+i`$.
Proof: For any finite dimensional representation $`\sigma `$ of $`M`$ let
$$Z_{P,\sigma }(s)=\underset{[\gamma ]_P^p(\mathrm{\Gamma })}{}\underset{n0}{}\mathrm{det}(1e^{sl_\gamma }\sigma (b_\gamma )S^n(a_\gamma b_\gamma |𝔫)),$$
where $`𝔫=Lie_{}(N)`$ and $`S^n(a_\gamma b_\gamma |𝔫)`$ is the $`n`$-th symmetric power of the adjoint action of $`a_\gamma b_\gamma `$ on $`𝔫`$. In , Theorem 2.1 it is shown that $`Z_{P,\sigma }`$ extends to a meromorphic function and that all its poles and zeros lie in $`(\frac{1}{2}+i)`$. Note that $`M`$ is isomorphic with $`SL_2^\pm ()`$, the group of real $`2\times 2`$ matrices of determinant $`\pm 1`$. Let $`\sigma _0:MGL_2`$ denote the standard representation. In , Theorem 4.1 it is shown that
$$R_\mathrm{\Gamma }(s)=\frac{Z_{P,1}(s)Z_{P,1}(s+1)}{Z_{P,\sigma _0}(s+\frac{1}{2})}.$$
So to complete the proof, it suffices to show the following proposition.
###### Proposition 3.24
For $`\sigma =1`$ the poles and zeros of $`Z_{P,\sigma }(s)`$ lie in $`[\frac{1}{4},\frac{3}{4}](\frac{1}{2}+i)\{0,1\}`$ and the function $`Z_{P,\sigma }(s)`$ has a simple zero at $`s=1`$.
For $`\sigma =\sigma _0`$ the poles and zeros of $`Z_{P,\sigma }`$ all lie in $`[0,1](\frac{1}{2}+i)`$.
Proof: Let $`\widehat{G}`$ denote the set of all isomorphism classes of irreducible unitary representations of $`G`$. The group $`G`$ acts on the Hilbert space $`L^2(\mathrm{\Gamma }\backslash G)`$ by translations from the right. Since $`\mathrm{\Gamma }\backslash G`$ is compact this representation decomposes discretely:
$$L^2(\mathrm{\Gamma }\backslash G)=\underset{\pi \widehat{G}}{}N_\mathrm{\Gamma }(\pi )\pi ,$$
with finite multiplicities $`N_\mathrm{\Gamma }(\pi )`$. For $`\pi \widehat{G}`$ let $`\pi _K`$ denote the $`(𝔤,K)`$-module of $`K`$-finite vectors in $`\pi `$. Then the Lie algebra $`𝔫`$ acts on $`\pi _K`$ and we denote by $`H^q(𝔫,\pi _K)`$ the corresponding Lie algebra cohomology . Let $`𝔪`$ denote the complexified Lie algebra of $`M`$ and let $`𝔪=𝔨_M𝔭_M`$ be its polar decomposition, where $`𝔨_M`$ is the complexified Lie algebra of $`K_M=KM`$.
In Theorem 2.1 it is shown that all poles and zeros of $`Z_{P,\sigma }`$ lie in $`(\frac{1}{2}+i)`$ and that for $`\lambda 𝔞^{}`$ the (vanishing-) order of $`Z_{P,\sigma }`$ at $`s=\lambda (H_1)`$ is
$$\underset{\pi \widehat{G}}{}N_\mathrm{\Gamma }(\pi )\underset{p,q0}{}(1)^{p+q}dim(H^q(𝔫,\pi _K)^p𝔭_MV_{\stackrel{˘}{\sigma }})_\lambda ^{K_M},$$
where $`(.)_\lambda `$ denotes the generalized $`\lambda `$-eigenspace. Note that the torus $`A`$ acts trivially on all tensor factors except $`H^q(𝔫,\pi _K)`$. Let $`\pi \widehat{G}`$ and let $`_\pi 𝔥^{}`$ be a representative of its infinitesimal character. Corollary 3.23 of says that every character of the $`𝔞`$-action on $`H^q(𝔫,\pi _K)`$ is of the form $`\mu =w_\pi +\rho _P`$ for some $`wW(𝔤,𝔥)`$.
To show the proposition we concentrate on $`\sigma `$ being trivial since the other case is similar. Since $`\rho (H_1)=\frac{1}{2}`$ we have to show that for every $`\pi `$ which has a nonzero contribution, all eigenvalues $`\lambda `$ of $`𝔞`$ on $`H^p(𝔫,\pi _k)`$ satisfy $`\frac{3}{2}\rho _P\mathrm{Re}(\lambda )\frac{1}{2}`$. By the isomorphism of $`AM`$-modules , p. 57:
$$H_p(𝔫,\pi _K)H^{2p}(𝔫,\pi _K)^2𝔫$$
this becomes equivalent to $`\frac{1}{2}\mathrm{Re}(\lambda )\frac{3}{2}\rho _P`$ whenever $`\lambda `$ is an eigenvalue of $`𝔞`$ on $`H_p(𝔫,\pi _K)`$. Now fix $`\pi \widehat{G}`$ and let $`_\pi 𝔥^{}`$ be a representative of its infinitesimal character. Then, according to Corollary 3.32 of we have to show that $`\frac{1}{2}\rho _P\mathrm{Re}(_\pi |_𝔞)\frac{1}{2}\rho _P`$.
In the case when $`\pi `$ is induced from the minimal parabolic it follows that its distributional character $`\mathrm{\Theta }_\pi `$ is zero on $`AB`$. By the construction of the test function in , p. 903, this implies that the contribution of $`\pi `$ is zero. Therefore, by the classification of the unitary dual in , it remains to consider the case of the trivial representation and the case when $`\pi `$ is unitarily induced from $`P=MAN`$. So let $`\pi =\pi _{\xi ,\nu }`$ be induced from $`P`$, where $`\nu `$ is imaginary. Then we may assume that $`\xi `$ is not induced, since otherwise the double induction formula would lead us back to the previous case. Let $`wW`$ and $`_\xi `$ be the infinitesimal character of $`\xi `$. We lift $`_\xi `$ to $`𝔡`$ by defining it to be zero on $`𝔞`$. Then the Weyl group $`W`$ will act on $`_\xi `$. We have to show that
$$\frac{1}{2}\rho _P\mathrm{Re}(w_\xi |_𝔞)\frac{1}{2}\rho _P.$$
Let us start with $`\xi `$ being the trivial representation. Then $`_\xi \left(\begin{array}{cc}t& \\ & t\end{array}\right)`$. Lifting $`_\xi `$ to $`𝔡`$ we get $`_\xi \left(\begin{array}{ccc}a& & \\ & b& \\ & & c\end{array}\right)=\frac{1}{2}(ab)`$. This vanishes on $`A`$, so it’s real part is zero. This deals with the case when $`w=1`$. For $`wW`$ being the transposition interchanging $`b`$ and $`c`$ we get $`w_\xi \left(\begin{array}{ccc}a& & \\ & b& \\ & & c\end{array}\right)=\frac{1}{2}(ac)`$, which, restricted to $`𝔞`$, coincides with $`\frac{1}{2}\rho `$. This implies the claim for this $`w`$. All other Weyl group elements are treated similarly.
It remains to consider the case when $`\xi `$ is a (limit of) discrete series representation, so $`\xi =𝒟_n^+𝒟_n^{}`$, where the notation is as in . Let $`\tau \widehat{K_M}`$ and let $`P_\tau :V_\xi V_\xi (\tau )`$ be the projection onto the $`\tau `$-isotype. For any function $`f`$ on $`G`$ which is sufficiently smooth and of sufficient decay the operator $`\pi (f)`$ is of trace class. Its trace is
$$\underset{\tau \widehat{K_M}}{}_K_{MAN}a^{\nu +\rho }f(k^1mank)\mathrm{tr}P_\tau \xi (m)P_\tau 𝑑man𝑑k.$$
Plugging in the test function $`f=\mathrm{\Phi }`$ constructed in , p.903, this gives
$$_{A^+}a^{\nu +\rho }l_a^{j+1}e^{sl_a}\mathrm{tr}\xi (f_1)𝑑a,$$
where $`f_1`$ is the Euler-Poincaré function on $`M`$ attached to the trivial representation. Then
$$\mathrm{tr}\xi (f_1)=\underset{p=0}{\overset{dim𝔭_M}{}}(1)^pdim(V_\xi \stackrel{p}{}𝔭_M)^{K_M}.$$
Now $`K_MO(2)`$ and the $`K_M`$-types can be computed explicitly. Thus one sees that $`\mathrm{tr}\xi (f_1)`$ can only be nonzero if $`n=1`$ or $`n=2`$. This means that either $`_\xi \left(\begin{array}{cc}t& \\ & t\end{array}\right)`$ equals $`0`$ or $`t`$ respectively, which, in a similar way to the above, implies the claim. In the case $`\sigma =\sigma _0`$ the function $`f_1`$ is replaced by $`f_{\sigma _0}`$ and one proceeds in the same fashion. This takes care of all induced representations. By the classification of the unitary dual of $`GL_3()`$ in it follows that it remains to worry about the trivial representation only, so let $`\pi =triv`$ be the trivial representation of $`G`$; then the space $`H_0(𝔫,\pi _K)=\pi _K/𝔫\pi _K`$ is one dimensional with trivial $`𝔞`$-action. This gives a simple zero at $`s=1`$.
$`\frac{}{}`$ Q.E.D.
## 4 Asymptotics of closed geodesics
For $`\gamma \mathrm{\Gamma }`$ let $`N(\gamma )=e^{l_\gamma }`$ and define for $`x>0`$:
$$\pi (x)=\mathrm{\#}\{[\gamma ]_P^p(\mathrm{\Gamma })|N(\gamma )x\}.$$
The geodesic prime number theorem in our context is
###### Theorem 4.27
For $`x\mathrm{}`$ we have the asymptotic
$$\pi (x)\frac{x}{\mathrm{log}x}.$$
More sharply we have that
$$\pi (x)=li(x)+O\left(\frac{x^{\frac{3}{4}}}{\mathrm{log}x}\right)$$
as $`x\mathrm{}`$ where $`li(x)=_2^{\mathrm{}}\frac{1}{\mathrm{log}t}𝑑t`$ is the integral logarithm.
Proof: To simplify the notation in what follows we write $`\gamma `$ for an element of $`_P(\mathrm{\Gamma })`$ and $`\gamma _0`$ for a primitive element. If $`\gamma `$ and $`\gamma _0`$ occur in the same formula it is understood that $`\gamma _0`$ will be the primitive element underlying $`\gamma `$. Unless otherwise specified, all sums will run either over $`\gamma `$ or $`\gamma _0`$. For $`x>0`$ let
$$\psi (x)=\underset{N(\gamma )x}{}l_{\gamma _0}.$$
For $`\mathrm{Re}(s)>1`$ we have
$`{\displaystyle \frac{R_\mathrm{\Gamma }^{}}{R_\mathrm{\Gamma }}}(s)`$ $`=`$ $`{\displaystyle \underset{\gamma _0}{}}{\displaystyle \frac{l_{\gamma _0}e^{sl_{\gamma _0}}}{1e^{sl_{\gamma _0}}}}`$
$`=`$ $`{\displaystyle \underset{\gamma _0}{}}l_{\gamma _0}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{snl_{\gamma _0}}`$
$`=`$ $`{\displaystyle \underset{\gamma }{}}l_{\gamma _0}e^{sl_\gamma }.`$
From this point on the argumentation is, up to minor changes the same as in , which leads to
$$\psi (x)=x+O(x^{\frac{3}{4}}).$$
From this, the theorem is deduced by standard techniques of analytic number theory.
$`\frac{}{}`$ Q.E.D.
We now finish the proof of the main theorem (1.1). For this we have to find a division algebra $`M()`$ such that for a given set of primes $`S`$ with at least two elements we have $`\pi _S(x)=\pi (x)`$. Firstly, there is a division algebra $`M()`$ of degree $`3`$ such that the set of places at which $`M()`$ is non-split, coincides with the set $`S`$. This algebra is obtained by taking the local Brauer-invariants at $`pS`$ to be equal to $`\frac{1}{3}`$ or $`\frac{2}{3}`$ and zero everywhere else in such a way that they sum to zero in $`/`$ (, Theorem 18.5). Next, we have a bijection
$$_P(\mathrm{\Gamma })O(S),$$
given by
$$[\gamma ]F_\gamma M(),$$
where $`F_\gamma `$ is the centralizer of $`\gamma `$ in $`M()`$. Under this bijection the length $`l_\gamma `$ is transferred to $`\mathrm{log}r(𝒪)`$. The theorem follows.
$`\frac{}{}`$ Q.E.D.
###### Corollary 4.32
We finally note that for $`M()`$ chosen as above the Ruelle zeta function of Theorem 3.23 takes the form:
$$R_\mathrm{\Gamma }(s/3)=\underset{𝒪O(S)}{}\left(1e^{sR(𝒪)}\right)^{h(𝒪)\lambda _S(𝒪)}$$
University of Exeter, Mathematics, Exeter EX4 4QE, Devon, UK
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# Sea Quark Effects on Quarkonia
## I Introduction
The experimental efforts to pin down the parameters of the Standard Model have been paralleled by intense theoretical attempts to provide experimentalists with non-perturbative constraints from Quantum Chromodynamics (QCD). It is hoped that Lattice QCD will ultimately provide such important information. To this end it is crucial to understand and control the systematic errors in numerical calculations, which rely on extrapolations and interpolations to the physical point. This important task is particularly demanding for heavy quark phenomenology, where one has to describe accurately both the light and heavy quarks in the system.
In particular the inclusion of light dynamical fermions in the gluon background is still a daunting task and requires the largest fraction of computational resources. In the past these restrictions led to the so-called quenched approximation, in which only valence quarks are allowed to propagate in a purely gluonic background, whereas the virtual creation of sea quarks is ignored. We have shown in a recent work that this results in systematic deviations in the lattice prediction of the light hadron spectrum from the observed experimental data. More recently it has also been found that the inclusion of two dynamical sea quarks has a significant effect on the light hadron spectrum and quark masses . This is of course analogous to QED, where the inclusion of all higher order effects, which could be made through perturbation theory, resulted in an impressive agreement with experimental observations. A distinctive aspect of QCD is that a proper non-perturbative treatment is in order so as to provide high-precision tests of this theory. Here we take this as our motivation to study heavy quarkonia in the presence (and absence) of dynamical sea quarks.
Heavy quark systems have long been considered an ideal testing ground for QCD and they have triggered the development of static potential models and heavy quark effective field theories . On the lattice, heavy quarks have frequently been studied using a non-relativistic approach to QCD (NRQCD ) or a relativistic formulation promoted by the Fermilab group . Both provide effective descriptions to deal with large scale differences, which are difficult to accommodate on conventional lattices. NRQCD has been quite successful in reproducing the spin-independent spectrum of heavy quarkonia owing to the fact that the quarks within such mesons move with small velocities $`v`$ such that $`v^2c^2`$.
As an effective field theory the predictive power of NRQCD relies on the control of higher dimensional operators, which have to be matched to the non-relativistic expansion of QCD. This has been the subject of many previous studies . As a result of these works it seemed plausible that bottomonium states could be accurately described in the NRQCD approach, whereas the spin structure of charmonium is very sensitive to the higher order relativistic corrections .
Within the NRQCD framework sea quark effects on the heavy quarkonium spectra have previously been studied at a lattice spacing of $`a0.1`$ fm using the Kogut-Susskind or the Wilson action for sea quarks. Here we present results for three lattice spacings in the range $`a0.20.1`$ fm, paying particular attention to the dependence on the sea quark mass and scaling properties. Our gauge configurations are generated with a renormalization-group (RG)-improved gluon action and a tadpole-improved clover quark action for two dynamical flavours . Some measurements are also made on quenched configurations generated with the same gluon action for making direct comparisons of dynamical and quenched results.
In Section II we introduce the actions which we use in our calculation. In Section III we give the details of our simulation, meson operators and fitting methods. Our results are discussed in Section IV and Section V concludes this paper.
## II Actions
### A Glue: RG-Improved Action
Since there is no unique discretisation of the continuum gluon action one can employ a set of operators to cancel some of the discretisation errors in the lattice action. The most common choice is to simply add a rectangular $`1\times 2`$ plaquette, $`\mathrm{Tr}R_{\mu \nu }`$, to the standard Wilson action, $`\mathrm{Tr}P_{\mu \nu }`$,
$$S_\mathrm{g}(g^2)=\frac{1}{g^2}\left\{c_0\mathrm{Tr}P_{\mu \nu }+c_1\mathrm{Tr}R_{\mu \nu }\right\},$$
(1)
where $`\mathrm{Tr}`$ denotes the trace over all indices and $`c_0+8c_1=1`$. All such choices have the same continuum limit, but different discretisation errors. Here we adopt a prescription which is motivated by an RG-analysis of the pure gauge theory ($`c_1=0.331`$ ). This has proven to be a suitable choice compared to, say, the Symanzik-improved action ($`c_1=1/12`$), for reducing scaling violation on coarse lattices. Instead of the coupling constant squared, $`g^2`$, we often quote the value of $`\beta 6/g^2`$.
### B Light Quarks: Clover Action
The standard discretisation of the fermion action removes the doublers at the expense of $`𝒪(a)`$ discretisation errors. It is possible to remove these errors by adding a single operator (the clover term) as first suggested in :
$$S_\mathrm{q}(g^2,m_q)=\overline{q}Qq=\overline{q}(/\mathrm{\Delta }+m_q)q+ar\overline{q}\mathrm{\Delta }^2qc_{sw}(g^2)ar\frac{ig}{4}\overline{q}\sigma _{\mu \nu }F_{\mu \nu }q.$$
(2)
Here the second term removes the doublers a la Wilson and the last is to reduce the resulting $`𝒪(a)`$ errors. We choose $`r=1`$ and $`c_{sw}=(10.1402g^2)^{3/4}`$. For the latter we follow the tadpole prescription of , which has been shown to improve the convergence of lattice perturbation theory significantly. Our choice is based upon the perturbative plaquette values as determined in . To one-loop order our choice differs from the correct value only by $`0.008g^2`$. Hence we expect only small $`𝒪(\alpha a)`$ scaling violations due to radiative corrections from the clover action, and $`𝒪(a^2)`$ errors from the gluon action.
In our simulations we work with two flavours of degenerate quarks of a common mass: $`m_q=m_u=m_d`$. For further reference, it is customary to replace the bare quark mass by the hopping parameter: $`\kappa 1/2(m_qa+4)`$. Since the direct simulation of realistic light Wilson quarks is not feasible on present-day computers we study the spectrum at a sequence of different $`\kappa `$.
### C Heavy Quarks: NRQCD
With the above actions we generated full QCD dynamical configurations on lattices of about 2.5 fm in spatial extent and lattice spacings ranging from approximately 0.1 to 0.2 fm. Such lattices are particularly suited to study light quark physics which is determined by a single scale: $`\mathrm{\Lambda }_{QCD}200`$ MeV. However, for systems containing slow-moving and heavy quarks we have to adjust the theoretical description to take into account all the non-relativistic scales: mass ($`m_Q`$), momentum ($`m_Qv`$) and kinetic energy ($`m_Qv^2`$). In this work we implement the NRQCD formulation to propagate the heavy quarks in a given gluon background. This approach has met with considerable success for b-quarks . Also charm quarks have previously been studied in this framework .
Whereas the relativistic boundary value problem requires several iterations to determine the quark propagator, the NRQCD approach has the numerical advantage to solve the two-spinor theory as a much simpler initial value problem. The forward propagation of the source vector, $`S(𝐱)`$, is described by:
$`G(𝐱,t+1)`$ $`=`$ $`\left(1{\displaystyle \frac{a}{2}}\delta H\right)\left(1{\displaystyle \frac{aH_0}{2n}}\right)^nU_t^{}(x)\left(1{\displaystyle \frac{aH_0}{2n}}\right)^n\left(1{\displaystyle \frac{a}{2}}\delta H\right)G(𝐱,t),t1,`$ (3)
$`G(𝐱,1)`$ $`=`$ $`\left(1{\displaystyle \frac{a}{2}}\delta H\right)\left(1{\displaystyle \frac{aH_0}{2n}}\right)^nU_t^{}(x)\left(1{\displaystyle \frac{aH_0}{2n}}\right)^n\left(1{\displaystyle \frac{a}{2}}\delta H\right)S(𝐱),t=0,`$ (4)
where
$`H_0`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }^2}{2m_Q}},`$ (5)
$`\delta H`$ $`=`$ $`c_0{\displaystyle \frac{\mathrm{\Delta }^4}{8m_Q^3}}c_1{\displaystyle \frac{1}{2m_Q}}\sigma g\stackrel{~}{𝐁}+c_2{\displaystyle \frac{i}{8m_Q^2}}(\stackrel{~}{\mathrm{\Delta }}g\stackrel{~}{𝐄}g\stackrel{~}{𝐄}\stackrel{~}{\mathrm{\Delta }})c_3{\displaystyle \frac{1}{8m_Q^2}}\sigma (\stackrel{~}{\mathrm{\Delta }}\times g\stackrel{~}{𝐄}g\stackrel{~}{𝐄}\times \stackrel{~}{\mathrm{\Delta }})`$ (8)
$`c_4{\displaystyle \frac{1}{8m_Q^3}}\{\mathrm{\Delta }^2,\sigma g\stackrel{~}{𝐁}\}c_5{\displaystyle \frac{1}{64m_Q^3}}\{\mathrm{\Delta }^2,\sigma (\stackrel{~}{\mathrm{\Delta }}\times g\stackrel{~}{𝐄}g\stackrel{~}{𝐄}\times \stackrel{~}{\mathrm{\Delta }})\}c_6{\displaystyle \frac{i}{8m_Q^3}}\sigma g\stackrel{~}{𝐄}\times g\stackrel{~}{𝐄}`$
$`c_7{\displaystyle \frac{a\mathrm{\Delta }^4}{16nm_Q^2}}+c_8{\displaystyle \frac{a^2\mathrm{\Delta }^{(4)}}{24m_Q}}.`$
The improved lattice operators $`\stackrel{~}{\mathrm{\Delta }}_i,\stackrel{~}{𝐄}`$ and $`\stackrel{~}{𝐁}`$ are defined as in . Other discretisations of NRQCD have been suggested in the past but they differ only at higher order in the lattice spacing. Here we follow and employ a formulation which includes all spin terms up to $`O(mv^6)`$ in non-relativistic expansion of QCD. On the coarsest lattice we checked explicitly that our equation (4) gives the same hyperfine splitting as from the asymmetric version employed in . The parameter $`n`$ was introduced to stabilise the evolution equation against high momentum modes. This is standard in such diffusive problems, but one should keep in mind that a change in $`n`$ will have to be accompanied by a change in $`m_Q`$ to simulate the same physical system. Alternatively one could decrease the temporal lattice spacing to prevent the high momentum modes from blowing up .
For a single quark source at point $`𝐲`$ we have $`S(𝐱)=\delta ^{(3)}(𝐱,𝐲)`$, but we also propagate extended objects with the same evolution equation (4). The operator $`H_0`$ is the leading kinetic term and $`\delta H`$ contains the relativistic corrections. The last two terms in $`\delta H`$ are present to correct for lattice spacing errors in temporal and spatial derivatives. For the derivatives we use the improved operators defined in and we also replace the standard discretized gauge field $`F_{\mu \nu }`$ by
$$\stackrel{~}{F}_{\mu \nu }=\frac{5}{3}F_{\mu \nu }\frac{1}{6}(U_\mu (x)F_{\mu \nu }(x+\mu )U_\mu ^{}(x)+U_\mu (x)F_{\mu \nu }(x\mu )U_\mu ^{}(x)(\mu \nu )).$$
(9)
With this prescription we aimed to achieve an accuracy of $`𝒪(a^4)`$ for the heavy quark sector. Of course we also expect terms of $`𝒪(\alpha a^2)`$ due to radiative corrections to this leading order result. In principle, we have to determine all coefficients in (8) by (perturbative or non-perturbative) matching to relativistic continuum QCD. Just as for the light quark sector we rely on a mean-field treatment of all gauge links to account for the leading radiative corrections. However, there is no unique prescription for such an improvement and several different schemes have been employed in the past. More recently it has been suggested that the average link variable in the Landau gauge should be less sensitive to radiative corrections since the gauge fields in Landau gauge have less UV fluctuations . Here we adopt this view and divide all links by the appropriate tadpole coefficient at each value of $`(\beta ,\kappa )`$:
$$U_\mu (x)U_\mu (x)/u_{0L},u_{0L}\frac{1}{3}\mathrm{Tr}U_\mu _{LG}.$$
(10)
An alternative and gauge-invariant implementation of mean-field improvement that has frequently been used in the past forces the average plaquette, $`P_{\mu \nu }`$, to unity:
$$U_\mu (x)U_\mu (x)/u_{0P},u_{0P}\frac{1}{3}\mathrm{Tr}P_{\mu \nu }^{1/4}.$$
(11)
In some selected cases we have also implemented this method to estimate the effect of unknown radiative corrections to the renormalisation coefficients, $`c_i`$. In all applications of (8) we set them to their tree-level value 1. We denote as $`𝒪(mv^6,a^2)`$ the evolution equation which includes all spin-dependent operators up to $`𝒪(mv^6)`$ and where all operators have been improved to reduce the $`𝒪(a^2)`$ errors.
## III Simulation
### A Updates, Trajectories and Autocorrelations
The gauge field configurations with two dynamical sea quarks used for the present study were generated on the CP-PACS supercomputer at the Center for Computational Physics, University of Tsukuba. They can be classified by two parameters, $`(\beta ,\kappa )`$, which determine the lattice spacing and the sea quark mass. A standard hybrid Monte Carlo algorithm is used to incorporate the effects of the fermion determinant. For the matrix inversion we implemented the BiCGStab algorithm. To reduce the autocorrelations between separate measurements we only used every fifth or tenth trajectory and binned the remaining data until the statistical error was independent of the bin size. In Tab. I we list the number of trajectories we generated for each set of couplings along with the other simulation parameters and the actual number of configurations we used in this study; for more details see . The subsequent determination of the quarkonium spectrum is subject of this work.
Since there has been no previous study of heavy quarkonia using the RG action for the gluon sector, we also supplemented our calculation by a comparative quenched calculation. The coupling constant $`\beta `$ of these quenched configurations were chosen so that the string tension of the static quark-antiquark potential matches that of one of the dynamical runs. The parameters of these runs are also given in Tab. I.
### B Meson Operators
To extract meson masses we calculate two-point functions of operators with appropriate quantum numbers. In a non-relativistic setting gauge-invariant meson operators can be constructed from the two-spinors $`\chi ^{}(x)`$, $`\mathrm{\Psi }(y)`$ and a Wilson line, $`W(x,y):\chi ^{}(x)W(x,y)\mathrm{\Psi }(y)`$. The construction of meson states with definite $`J^{PC}`$ from those fundamental operators is standard (on the lattice $`J`$ labels the irreducible representations of the octahedral group ($`J=A_1,A_2,T_1,T_2,E`$)).
Since here we are only interested in S and P-states it is sufficient to consider $`\chi ^{}(x)\mathrm{\Psi }(x)`$ and $`\chi ^{}(x)\mathrm{\Delta }_i\mathrm{\Psi }(x)`$, respectively. The corresponding spin-triplets can be constructed by inserting the Pauli matrices into those bilinears. We also sum over different polarisations to increase the statistics.
The overlap of those simplistic operators with the state of interest can be further improved upon. One way is to employ wavefunctions, which try to model the ground state more accurately. This requires assumptions about the underlying potential and gauge-fixed configurations. We decided to use a gauge-invariant smearing technique, which regulates the extent of the lattice operator, with a single parameter $`ϵ`$:
$$\chi ^{}(x)\widehat{O}\mathrm{\Psi }(x)\chi ^{}(x)\widehat{O}(1ϵ\mathrm{\Delta }^2)^{10}\mathrm{\Psi }(x).$$
(12)
For computational ease we limited this procedure to 10 smearing iterations and implemented it only at the source. From such operators we obtain the meson correlators as a Monte Carlo average over all configurations
$$\mathrm{C}_ϵ(x,y)=\mathrm{𝐭𝐫}\left[G^{}(x,y)\widehat{O}G(x,y)(1ϵ\mathrm{\Delta }^2)^{10}\widehat{O}^{}\right],$$
(13)
where $`\mathrm{𝐭𝐫}`$ denotes contraction over all internal degrees of freedom.
For the smeared propagator we solve (4) with $`S(𝐱,𝐲)=\delta (𝐱,𝐲)(1ϵ\mathrm{\Delta }^2)^{10}\widehat{O}^{}`$. We fix the origin at some (arbitrary) lattice point $`y=(𝐲,0)`$. This creates a meson state with all possible momenta. In practice we employed up to 8 spatial sources on every configuration. This is permissible since heavy quarkonia are small compared to the lattice extent of about 2.5 fm. We explicitly checked the independence of such measurements by binning. At the sink, $`x=(𝐱,t)`$, we perform a Fourier transformation to project the correlator onto a given momentum eigenstate:
$$\mathrm{C}_{ϵ,t}(𝐩)=\underset{𝐱}{}C_ϵ(𝐱,t)\mathrm{exp}(i\mathrm{𝐩𝐱}).$$
(14)
In the trivial case of zero momentum this amounts to simply summing over all spatial $`𝐱`$.
### C Fitting
Since there is no backward propagation of the heavy quark in our framework, we can fit the meson propagators to the exponential form:
$$y_{ϵ,t}(a_i,E_i)=\underset{i=1}{\overset{\mathrm{n}_{\mathrm{fit}}}{}}a_i(\alpha ,𝐩,ϵ)e^{E_i(\alpha ,𝐩)t}.$$
(15)
This is the theoretical prediction for a multi-exponential decay of a state with momentum $`𝐩`$ and quantum numbers $`\alpha `$ along Euclidean time $`t`$. Different choices for the smearing parameter $`ϵ`$ will result in different overlaps with the ground state and the higher excited states of the same quantum number. In practice it is difficult to project directly onto a given state, so we chose to extract the ground state from multi-exponential correlated fits. In some cases we were also able to extract the first excited states reliably from our data. The simplest way to visualise our data is by means of effective mass plots, which are expected to display a plateau for long Euclidean times. In Fig. 1 we show a representative plot for one set of simulation parameters.
Correlations between different times, $`t`$, and different smearing radii, $`ϵ`$, of the meson propagator $`C_{ϵ,t}`$ are taken into account by employing the full covariance matrix for the $`\chi ^2`$-minimisation:
$$\chi ^2(\pi )\underset{r,s=1}{}(C_ry_r(\pi ))\mathrm{\Gamma }_{rs}^1(C_sy_s(\pi )),\frac{\chi ^2}{\pi _k}\left(\overline{\pi }\right)=0.$$
(16)
Here, in order to ease the notation, we introduced multi-indices ($`r=[ϵ,t]`$) for the data points and $`\pi =[a_i,E_i]`$ for the parameters.
Our statistical ensembles are large enough to determine the covariance matrix, $`\mathrm{\Gamma }_{rs}^1`$, with sufficient accuracy. Therefore the inversion of this matrix did not present a problem. For the spin splittings we applied (15) also to the ratio of two propagators, $`C(\alpha _1)/C(\alpha _2)`$. In this way we utilized correlations between states of different quantum numbers to obtain improved estimates for their energy difference.
Statistical fluctuations in the data cause fluctuations in the fit results determined by (16). We estimate the covariance matrix, $`\mathrm{\Delta }_{kl}`$, of the fitted parameters, $`\pi _k`$, from the inverse of $`(^2\chi ^2)/(\pi _k\pi _l)`$. We have checked these errors against bootstrap errors which gave consistent results. We also require consistent results as we change the fit ranges $`(t_{min},t_{max})`$ or the number of exponentials to be fitted. This is illustrated in Fig. 2, where we show the $`t_{min}`$-plot for the S-state of Fig. 1. The goodness of each fit is quantified by its Q-value and we demand an acceptable fit to have $`Q>0.1`$. Finally we subjected those results to a binning analysis, which takes into account autocorrelations of the same measurement at different times in the update process.
## IV Results
We now present our new results for elements of the spectrum of heavy quarkonia. Our data from two quenched lattices is given in Tab. II and the dynamical data is collected in Tabs. III to VI. For notational ease we use dimensionless lattice quantities throughout the remainder of this paper, unless stated otherwise. To convert the lattice predictions into dimensionful quantities we take the experimental $`1P1S`$ splitting to set the scale.
### A Heavy Mass Dependence and Kinetic Mass
For a given value of the gauge coupling (lattice spacing) and the mass of the two degenerate sea quarks there is only one remaining parameter to choose – the mass of the heavy valence quark. On the lattice we are free to simulate every arbitrary value, but in order to obtain physical results we tune the ratio $`M_{\mathrm{kin}}/(1P1S)`$ of the kinetic mass of a quarkonium state and the $`1P1S`$ mass splitting to its experimental value. The determination of $`1S`$ and $`1P`$ masses has already been described in Sec. III C. The kinetic mass, $`M_{\mathrm{kin}}`$, is defined through the dispersion relation of the quarkonium state:
$$E(𝐩)E(\mathrm{𝟎})=\frac{𝐩^2}{2M_{\mathrm{kin}}}+\mathrm{},𝐩=\frac{2\pi }{L}(n_1,n_2,n_3).$$
(17)
For each heavy quark mass, $`m_Q`$, we project the $`{}_{}{}^{1}S_{0}^{}`$-state and the $`{}_{}{}^{3}S_{1}^{}`$-state onto 5 different momentum eigenstates: $`(n_1,n_2,n_3)=(0,0,0);(1,0,0);(1,1,0);(1,1,1)`$ and $`(2,0,0)`$. We obtain $`E(𝐩)E(\mathrm{𝟎})`$ from ratio fits and determine the kinetic mass by fitting the data to the dispersion relation. To this end we have also included higher terms in the expansion of (17) and find consistent results for $`M_{\mathrm{kin}}`$. However, such fits tend to have larger errors and the coefficient of $`𝐩^4`$ is not well determined. For better accuracy we normally restrict ourselves to a linear fit in $`𝐩^2`$ for the lowest four momenta. An example of this procedure is given in Fig. 3. We have plotted the fitted values of $`M_{\mathrm{kin}}`$ against the heavy quark mass in Fig.4. Large discretisation errors can be seen for almost all masses, but once we include all $`O(a^2)`$ correction terms in (4), the discrepancy due to higher order relativistic corrections is small. Comparing the relative changes in Fig. 3 due to $`O(a^2)`$ improvement at different momentum scales, we can also estimate the size of higher order corrections and we expect them to be small.
In this paper we tuned the bare quark mass on all our lattices $`(\beta ,\kappa )`$, so as to reproduce the experimental value of the mass of $`\mathrm{{\rm Y}}`$ ($`M_{\mathrm{kin}}=9.46`$ GeV). In some selected cases we interpolated the spectrum to this physical point.
### B Scale Determination and $`1P1S`$ Splitting
It has been noticed in the past that the tuning of $`m_Q`$ can be done efficiently since the spin-averaged splitting is not very sensitive to the quark mass parameter. However, with our newly achieved accuracy we could also resolve a slight mass dependence of $`1P1S`$ in the range from charmonium to the bottomonium system as shown in Fig. 5. The experimental values for the $`1P1S`$ splitting show a 4% decrease when going from charmonium (458 MeV) to bottomonium (440 MeV), which should be compared to a reduction of about 10% from our simulation at $`N_f=2`$ and an unphysical sea quark mass. This larger change is related to the running of the strong coupling between the two scales, which still does not fully match the running coupling in nature. It is expected that the modified short-range potential will result in a ratio $`(1P1S)_{c\overline{c}}/(1P1S)_{b\overline{b}}`$ bigger than in experiment.
While the heavy quark mass can be tuned to its physical value as described in the previous section, this is not possible for the light quark mass and one has to rely on extrapolations to realistic quark masses, where the ratio $`m_\pi /m_\rho `$ equals the experimental value. Here we are mainly interested in the behaviour of physical quantities as we approach the chiral limit. We take $`m_\pi ^2`$ as a measure of the light quark mass and extrapolate quadratically in this parameter. This is a common procedure but we will demonstrate below that the physical dependence on the sea quark mass may indeed be difficult to disentangle from unphysical scaling violations. In taking the naive chiral limit we hope to account for at least a fraction of the spectral changes towards smaller sea quark masses. At $`\beta =2.10`$ we only have results from two values of $`\kappa `$ and take a linear estimate for the chiral limit. The chiral behaviour of the $`1P1S`$ splitting is shown in Fig. 6 for all values of $`\beta `$ in our study.
In quenched simulations there exist uncertainties when setting the scale from different physical quantities. We expect these uncertainties to be reduced in our simulations incorporating two light dynamical flavours. To examine this point we compare in Fig. 7 our results for $`1P1S`$ splitting with the data for $`m_\rho `$ as a representative example from the light quark sector. If it were not for quenching effects and lattice spacing artefacts, one would expect the ratio $`m_\rho /(1P1S)`$ to equal its experimental value.
It is encouraging to see that the dynamical calculations are always and significantly closer to the experimental value of 1.75 than the corresponding quenched simulations. This demonstrates the importance of dynamical over quenched simulations. The scaling violations in this ratio do not fully cancel, however; we observe a 10% shift in $`m_\rho /(1P1S)`$ over $`a0.20.1`$ fm. Keeping in mind that we are working on rather coarse lattices with $`a>0.1`$ fm, the remaining scaling violations are perhaps not too surprising.
Looking at the ratio $`(1P1S)/\mathrm{\Lambda }_{\mathrm{QCD}}`$ it is clear that our data does not satisfy the strict criterion of asymptotic scaling, see Fig. 8. In this plot we determined $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ from the 2-loop formula in the $`\overline{MS}`$ scheme,
$$\mathrm{\Lambda }_{\mathrm{QCD}}=\pi \left(\frac{\alpha b_0}{4\pi }\right)^{(b_1/2b_0^2)}\mathrm{exp}\left(\frac{2\pi }{b_0\alpha }\right)\left(1+\alpha \frac{b_1^2b_2b_0}{8\pi b_0^3}\right),$$
(18)
where the $`\overline{MS}`$ coupling constant $`\alpha =\alpha _{\overline{MS}}(\pi /a)`$ is estimated with a tadpole-corrected one-loop relation defined by
$$\frac{1}{\alpha _{\overline{MS}}(\pi /a)}=\frac{(3.648P2.648R)}{\alpha _0}+4\pi (0.0589+0.0218N_f).$$
(19)
Here $`\alpha _0=g^2/4\pi `$ denotes the bare coupling, and $`P`$ and $`R`$ are, respectively, $`1\times 1`$ and $`1\times 2`$ Wilson loops normalized to unity for $`U_\mu (x)=1`$.
Within the effective approach of NRQCD, we cannot extrapolate such scaling violations away and it is crucial to find other ratios in which the scaling violations cancel each other already at finite lattice spacing. In Fig. 9 we show a test of this nature for the string tension, which shows a better scaling. Here we plot as open symbols the data obtained from 4 different sea quark masses. At $`\beta =2.20(2.10)`$ we only measured the lightest (and heaviest) sea-quark mass, corresponding to $`m_\pi /m_\rho 0.60(0.80)`$. This figure also suggests that the string tension, in units of the $`1P1S`$ mass splitting, is smaller for 2-flavour QCD when compared to the quenched ($`N_f=0`$) theory.
### C Hyperfine Splitting
Quenching effects are also expected to show up in short-range quantities, since they are particularly sensitive to the shape of the QCD potential. In this difference has been demonstrated explicitly by observing a change in the Coulomb coefficient of the static potential. In the context of heavy quarkonia, the hyperfine splitting is such a UV-sensitive quantity which should be particularly susceptible to changes in the number of flavours and the sea quark mass. The prediction from potential models is
$${}_{}{}^{3}S_{1}^{}{}_{}{}^{1}S_{0}^{}=\frac{32\pi \alpha _s(q)}{9m_Q^2}|\mathrm{\Psi }(0)|^2.$$
(20)
In our study this is the most accurately measured quantity and it is clearly very sensitive to the value of the heavy quark mass, see Fig. 10. As has been noticed previously, higher order relativistic and radiative corrections are equally important for precision measurements of the hyperfine splitting in bottomonium and even more so for charmonium . Here we employ $`𝒪(mv^6,a^2)`$ as the standard accuracy throughout this paper.
The equation (20) involves a direct dependence on both the strong coupling and the wavefunction at the origin, which makes the hyperfine splitting an ideal quantity to study quenching effects. Here we also observe a clear rise of the hyperfine splitting as we decrease the sea quark mass, see Fig. 11.
In Fig. 12 we collected all our dynamical results for the hyperfine splitting over the range of $`0.10.2`$ fm. Here we plotted the data from each sea-quark mass as open symbols and used the experimental $`1P1S`$ splitting to convert lattice data into MeV. One should keep in mind that these points correspond to unphysical bottomonium in a world of different sea quark masses. We also plot as filled symbols the results of our naive chiral extrapolation as described in the previous section.
At around 0.10 fm we notice a very good agreement with the only previous calculation . These authors have performed a dynamical simulation at a single lattice spacing using Wilson glue and unimproved sea-quarks. For the bottom quarks they used an NRQCD formulation with the same accuracy, $`O(mv^6,a^2)`$, as in this study.
An unpleasant feature with our results in Fig. 12 is lack of scaling; for both the full and quenched case we find scaling violations of about 100 MeV/fm for the hyperfine splitting. Nonetheless, we do find several indications for sea quark effects in our results. First we notice that, if it were not for sea quark effects, then all points in Fig. 12 would lie on a universal curve which is not the case. This is a strong indication that for this quantity we have to expect effects of the order of 3-5 MeV when going from zero to two flavour QCD.
To substantiate this observation we make a direct comparison of quenched and dynamical calculations at the same lattice spacing of 0.14 fm in Fig. 13. For the $`{}_{}{}^{3}S_{1}^{}{}_{}{}^{1}S_{0}^{}`$ splitting replotted from Fig. 12, a clear increase of around 5 MeV (20%) represents more than a $`5\sigma `$-effect, which reflects the accuracy in our determination of this quantity. On the other hand, the hyperfine splitting in $`P`$-states is reduced as we approach a more realistic description of QCD. Within potential models, states with $`L0`$ are not sensitive to the contact term of the spin-spin interaction. However the perturbative expression for a higher order radiative corrections gives a $`{}_{}{}^{3}P{}_{}{}^{1}P_{1}^{}`$ splitting opposite in sign to our values. Experimentally, the spin-triplet states are well established, but $`{}_{}{}^{1}S_{0}^{}`$ and $`{}_{}{}^{1}P_{1}^{}`$ have yet to be confirmed for bottomonium.
We comment that our data for charmonium (Tab. VI) also indicates a rise in the hyperfine splitting towards the chiral limit. It is, however, apparent that such a rise can not explain the discrepancy between the NRQCD predictions and the experimentally observed spin structure. We confirm an earlier observation that the velocity expansion is not well controlled for charmonium where $`v_c^20.3`$.
### D Fine Structure
In the continuum, the fine structure in quarkonia is due to the different ways in which the spin can couple to the orbital angular momentum of the bound state. In our approach, the spin-orbit term and the tensor term of potential models can be traced back to the $`c_3`$ term in (8). A correct description of the fine structure will therefore require a proper determination of $`c_3`$ and the corrections to this term.
On the lattice we have also additional splittings with no continuum analogue. For example, the $`{}_{}{}^{3}P_{2E}^{}{}_{}{}^{3}P_{2T}^{}`$ splitting is known to be a pure discretisation error since the lattice breaks the rotational invariance of the continuum and causes the $`J=2`$ tensor to split into two irreducible representations of the orthogonal group: $`T_2`$ and $`E`$. Indeed, for both the dynamical calculations as well as the quenched data, we observe a significant reduction of this splitting when the lattice spacing is decreased; the splitting diminishes from about 1.5 MeV at $`a0.2`$ fm to 0.5 MeV at $`a0.1`$ fm for the dynamical case.
In Fig. 14 we show our results for the fine structure. For $`{}_{}{}^{3}P_{2}^{}`$ and $`{}_{}{}^{3}P_{1}^{}`$ we observe no clear dependence on the sea quark mass. This is not totally unexpected since $`P`$-state wavefunctions vanish at the origin and should not be as strongly dependent on changes in the UV-physics. In any case such small changes would be difficult to resolve within our statistical errors.
From Fig. 14 we can also see a better scaling behaviour of the $`P`$ states, apart perhaps from the $`{}_{}{}^{3}P_{0}^{}`$, where scaling violations still obscure the chiral behaviour. The latter has $`J=0`$ and therefore we may expect that for this state restoration of rotational invariance is particularly important.
On our finer lattices we observe an increase of the $`{}_{}{}^{3}P_{0}^{}{}_{}{}^{3}P`$ splitting, closer towards the experimental value of $`40`$ MeV. We take this as an indication that a better control of the lattice spacing errors and radiative corrections is necessary to reproduce this quantity in NRQCD lattice calculations. The other splittings, $`{}_{}{}^{3}P_{2}^{}{}_{}{}^{3}P`$ and $`{}_{}{}^{3}P_{1}^{}{}_{}{}^{3}P`$, deviate only by a few MeV from their experimental values of 12 MeV and $`8`$ MeV, which could be due to missing dynamical flavours ($`N_f=3`$), higher order relativistic effects and radiative corrections.
We take the fine structure ratio, $`R_{fs}=({}_{}{}^{3}P_{2}^{}{}_{}{}^{3}P_{1}^{})/({}_{}{}^{3}P_{1}^{}{}_{}{}^{3}P_{0}^{})`$, as a particularly sensitive quantity to measure the internal consistency of the $`P`$-triplet structure. This quantity should be less sensitive to radiative corrections of the NRQCD coefficients away from their tree-level values. Previous NRQCD calculations had measured this quantity to be much larger than 1, compared to the experimental value of 0.66(4). We believe that this discrepancy is due to lattice spacing artefacts as it is very sensitive to the implementation of $`𝒪(a^2)`$ improvement in the NRQCD formalism. It is encouraging to see that this value is further reduced on our finer lattices, see Fig. 15. Notably, we do not observe any difference between our dynamical results and the quenched data.
### E $`2S1S`$ Splitting
Another spectroscopic quantity which has attracted much attention is the $`2S1S`$ splitting, since it should also be sensitive to the short-range potential. On conventional lattices such higher excitations are difficult to resolve and requires delicate tuning to minimise the mixing of the $`2S`$ with the ground state. Given our rather coarse lattices we did not attempt to perform a systematic study of this quantity, but in the context of this section it is important to notice that we do not observe any chiral dependence of the ratio, $`R_{2S}=(2S1S)/(1P1S)`$.
In Fig. 16 we compiled representative data from other groups along with our new results from the RG-action. Within the large errors we cannot resolve a discrepancy between the experiment and the lattice data. This result is in contrast to the previous determinations of this quantity which claim to see deviations due to missing sea quarks .
Observing such deviations is certainly plausible as this ratio is thought to be sensitive to differences in the underlying short-range potential. However, for the same reason we should also expect large scaling violations. Interestingly, on our coarsest dynamical lattices we even observe smaller values of $`R_{2S}`$, which we take as an indication of large discretisation errors. Apart from this very coarse lattice data, we cannot resolve either scaling violations or quenching effects. We feel that it requires a much better resolution of the higher excited states, which is hard to achieve on isotropic lattices. Future lattice studies will need optimised meson operators or finer temporal discretisations to observe these effects.
## V Conclusion
We have demonstrated that dynamical sea quarks have a significant effect on the spectrum of heavy quarkonia. Namely the hyperfine splitting, $`{}_{}{}^{3}S_{1}^{}{}_{}{}^{1}S_{0}^{}`$, is raised by almost $`20\%`$ when going from zero to two dynamical flavours. The efficiency of the NRQCD approach has played an important role to establish such effects, but the numerical simplicity of this approach is offset by additional systematic errors, which have to be controlled. The sensitivity of the spin-structure to relativistic, $`𝒪(mv^6)`$, and radiative, $`𝒪(\alpha )`$, corrections was well known before we started this work. Here we demonstrated that quenching errors are equally important for precision measurements of the spectrum of heavy quarkonia. Perhaps more worrying are scaling violations, which we could resolve in many quantities. Without a proper control of lattice spacing artefacts it is not possible to make predictions for such UV-sensitive quantities as the hyperfine splitting on the lattices we used here.
While the lattice predictions for $`{}_{}{}^{3}P_{1}^{}{}_{}{}^{3}P`$ and $`{}_{}{}^{3}P_{2}^{}{}_{}{}^{3}P`$ agree well with their experimental values, the determination of the $`{}_{}{}^{3}P_{0}^{}`$ is more problematic and we still observed large deviations from the experimental value when the $`1P1S`$ splitting is used to set the scale. Clearly much work remains to be done to reduce the systematic errors in heavy quark physics to the same degree as the statistical ones. We feel that this may be difficult to achieve within the non-relativistic framework. A better description of the NRQCD coefficients or a relativistic treatment is in order to describe the spin structure in quarkonia accurately. From this work it is apparent that full QCD simulations are also necessary to achieve such a goal.
For less accurate quantities such as $`2S1S`$ it is more difficult to reach a similar conclusion and we leave a decisive observation of both scaling violations and sea quark effects to future studies with refined methods.
Acknowledgments
The calculations were done using workstations and the CP-PACS facilities at at the Center for Computational Physics at the University of Tsukuba. This work is supported in part by the Grants-in-Aid of Ministry of Education (Nos. 09304029, 10640246, 10640248, 10740107, 11640250, 11640294, 11740162). TM and AAK are supported by the JSPS Research for the Future Program (Project No. JSPS-RFTF 97P01102). SE, KN and HPS are JSPS Research Fellows.
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# Theorem 1
A Note on Cyclic Gradients
Dan Voiculescu
To the memory of Gian-Carlo Rota
The cyclic derivative was introduced by G.-C. Rota, B. Sagan and P. R. Stein in as an extension of the derivative to noncommutative polynomials. Here we show that there are simple necessary and sufficient conditions for an $`n`$-tuple of polynomials in $`n`$ noncommuting indeterminates to be a cyclic gradient (see Theorem 1) and similarly for a polynomial to have vanishing cyclic gradient (see Theorem 2). Our interest in cyclic gradients stems from free probability theory and random matrices (see the Remark at the end) ,,,,. This note should also reduce the paucity of results on cyclic derivatives in several variables pointed out in \[3, page 73\].
Let $`K_n=KX_1,\mathrm{},X_n`$ be the ring of polynomials in noncommuting indeterminates $`X_1,\mathrm{},X_n`$ with coefficients in the field $`K`$ of characteristic zero. The partial generalized difference quotients are the derivations
$$_j:K_nK_nK_n$$
such that $`_jX_k=0`$ if $`jk`$ and $`_jX_j=11`$. The $``$ here is over $`K`$ and $`K_nK_n`$ is given the bimodule structure such that $`a(bc)=abc`$, $`(bc)d=bcd`$.
The partial cyclic derivatives are then
$$\delta _j=\stackrel{~}{\mu }_j:K_nK_n$$
where $`\stackrel{~}{\mu }(ab)=ba`$.
We shall denote by $`N:K_nK_n`$ the “number operator”, i.e. the linear map so that $`N1=0`$, $`NX_{i_1}\mathrm{}X_{i_k}=kX_{i_1}\mathrm{}X_{i_k}`$. Also, $`CK_n`$ will denote the cyclic subspace, i.e. the vector subspace spanned by all cyclic symmetrizations of monomials
$$CX_{i_1}\mathrm{}X_{i_p}=\underset{1jp}{}X_{i_{j+1}}\mathrm{}X_{i_p}X_{i_1}\mathrm{}X_{i_j},p1\text{and}C1=0$$
(the constants are not in the cyclic subspace).
###### Theorem 1
Let $`P_1,\mathrm{},P_nK_n`$. The following conditions are equivalent:
(i) there is $`PK_n`$ such that $`\delta _jP=P_j`$ $`(1jn)`$.
(ii) $`{\displaystyle \underset{1jn}{}}[X_j,P_j]=0`$.
(iii) $`{\displaystyle \underset{1jn}{}}X_jP_jCK_n`$.
(iv) $`\delta _k\left({\displaystyle \underset{1jn}{}}X_jP_j\right)=(N+I)P_k`$. (Here $`I`$ denotes the identity map of $`K_n`$ to itself.)
Proof. It is easily seen that it suffices to prove the theorem for homogeneous $`P_1,\mathrm{},P_n`$ of the same degree, i.e. we may assume $`NP_j=sP_j`$ $`(1jn)`$ for some $`s0`$. Also the case of constants being obvious we will concentrate on $`s1`$.
(i) $``$ (ii) To check that $`{\displaystyle \underset{1jn}{}}[X_j,\delta _jP]=0`$ it suffices to do so when $`P`$ is a monomial $`X_{i_0}\mathrm{}X_{i_s}`$. Then
$$\delta _jP=\underset{i_p=j}{}X_{i_{p+1}}\mathrm{}X_{i_s}X_{i_0}\mathrm{}X_{i_{p1}}$$
so that
$$[X_j,\delta _jP]=\underset{i_p=j}{}\left(X_{i_p}\mathrm{}X_{i_s}X_{i_0}\mathrm{}X_{i_{p1}}X_{i_{p+1}}\mathrm{}X_{i_s}X_{i_0}\mathrm{}X_{i_p}\right)$$
and hence
$$\underset{1jn}{}[X_j,\delta _jP]=\underset{1ps}{}\left(X_{i_p}\mathrm{}X_{i_s}X_{i_0}\mathrm{}X_{i_{p1}}X_{i_{p+1}}\mathrm{}X_{i_s}X_{i_0}\mathrm{}X_{i_p}\right)=0.$$
(ii) $``$ (iii) Let
$$P_j=\underset{i_1,\mathrm{},i_s}{}c_{i_1\mathrm{}i_s}^jX_{i_1}\mathrm{}X_{i_s}$$
The coefficient of $`X_{i_0}\mathrm{}X_{i_s}`$ in $`{\displaystyle \underset{1jn}{}}[X_j,P_j]`$ is
$$c_{i_1\mathrm{}i_s}^{i_0}c_{i_0\mathrm{}i_{s1}}^{i_s}.$$
Hence (ii) gives $`c_{i_1\mathrm{}i_s}^{i_0}=c_{i_0\mathrm{}i_{s1}}^{i_s}`$. On the other hand, if $`c_{i_0\mathrm{}i_s}`$ denotes the coefficient of $`X_{i_0}\mathrm{}X_{i_s}`$ in $`{\displaystyle \underset{ijn}{}}X_jP_j`$, clearly $`c_{i_0\mathrm{}i_s}=c_{i_1\mathrm{}i_s}^{i_0}`$, so that (ii) implies $`c_{i_0\mathrm{}i_s}=c_{i_si_0\mathrm{}i_{s1}}`$, i.e. cyclicity.
(iii) $``$ (iv) As before, let $`c_{i_1\mathrm{}i_s}^j`$ and $`c_{i_0\mathrm{}i_s}`$ denote the coefficients of $`P_j`$ and $`{\displaystyle \underset{j}{}}X_jP_j`$ respectively. Then $`c_{i_1\mathrm{}i_s}^{i_0}=c_{i_0\mathrm{}i_s}`$ and the cyclicity condition gives $`c_{i_0\mathrm{}i_s}=c_{i_si_0\mathrm{}i_{s1}}`$. We have
$`\delta _k\left({\displaystyle \underset{k}{}}X_jP_j\right)`$
$`={\displaystyle \underset{i_0\mathrm{}i_s}{}}{\displaystyle \underset{\{r:i_r=k\}}{}}c_{i_0\mathrm{}i_s}X_{i_{r+1}}\mathrm{}X_{i_s}X_{i_0}\mathrm{}X_{i_{r1}}`$
$`={\displaystyle \underset{i_0\mathrm{}i_s}{}}{\displaystyle \underset{\{r:i_r=k\}}{}}c_{i_{r+1}\mathrm{}i_si_0\mathrm{}i_{r1}}^kX_{i_{r+1}}\mathrm{}X_{i_s}X_{i_0}\mathrm{}X_{i_{r1}}`$
$`={\displaystyle \underset{0rs}{}}{\displaystyle \underset{i_0\mathrm{}i_{r1}i_{r+1}\mathrm{}i_s}{}}c_{i_{r+1}\mathrm{}i_si_0\mathrm{}i_{r1}}^kX_{i_{r+1}}\mathrm{}X_{i_s}X_{i_0}\mathrm{}X_{i_{r1}}=(s+1)P_k.`$
(iv) $``$ (i) Since the $`P_j`$ are homogeneous of the same degree and the field characteristic is zero, this is obvious. $`\mathrm{}`$
There is also a simple description of the noncommutative polynomials with vanishing cyclic gradient.
###### Theorem 2
We have
$$\mathrm{Ker}\delta =\underset{1kn}{}[X_k,K_n]+1=1+[K_n,K_n]=\mathrm{Ker}C$$
Proof. (i) Ker $`\delta \text{Ker}C`$. We have
$$Cp=\underset{1jm}{}X_j\delta _jp=0$$
(ii) Clearly,
$$\underset{1kn}{}[X_k,K_n]+1[K_n,K_n]+1$$
Also, since $`1\text{Ker}C`$ and $`[X_{i_1}\mathrm{}X_{i_r},X_{i_{r+1}}\mathrm{}X_{i_{r+s}}]`$ is the difference of two cyclic permutations of $`X_{i_1}\mathrm{}X_{i_{r+s}}`$, we have $`[K_n,K_n]+1\text{Ker}C`$.
To see that Ker $`C_{1kn}[X_k,K_n]+1`$, remark that Ker $`C`$ is spanned by homogeneous elements and that $`Cp=0`$, where $`p`$ is homogeneous of degree $`m`$ iff $`p`$ is a linear combination of differences $`X_{i_1}\mathrm{}X_{i_m}X_{i_2}\mathrm{}X_{i_m}X_{i_0}`$.
(iii) To see that $`{\displaystyle \underset{1kn}{}}[X_k,K_n]+1\text{Ker}\delta `$, it suffices to show that
$`[X_k,X_{i_1}\mathrm{}X_{i_s}]\text{Ker}\delta `$. This is clearly so, since
$$[X_k,X_{i_1}\mathrm{}X_{i_s}]=X_kX_{i_s}\mathrm{}X_{i_s}X_{i_1}\mathrm{}X_{i_s}X_k$$
and the cyclically equivalent elements $`X_kX_{i_1}\mathrm{}X_{i_s}`$, $`X_{i_1}\mathrm{}X_{i_s}X_{i_k}`$ have the same cyclic gradient. $`\mathrm{}`$
Putting together the two theorems, we have an exact sequence
$$0[K_n,K_n]K_n\stackrel{\delta }{}(K_n)^n\stackrel{\theta }{}K_n$$
where $`\theta ((P_j)_{1jn})={\displaystyle \underset{j}{}}[X_j,P_j]`$.
Remark. The motivation for this note is from free entropy and large deviations for random matrices. Let $`(M,\tau )`$ be a von Neumann algebra with normal faithful trace-state $`\tau `$ and $`X_k=X_k^{}M`$ $`(1kn)`$ which are algebraically free.
Let $`𝒥_k=𝒥(X_k:X_1,\mathrm{},X_{k1},X_{k+1},\mathrm{},X_n)`$ be the noncommutative Hilbert transforms defined in , in connection with free entropy. On the other hand, the upper bound for large deviations for $`n`$-tuples of random matrices found in fits well with free entropy except for a term involving cyclic gradients and about which it is not known whether it is not actually zero. The precise question is, whether the $`n`$-tuple $`(𝒥_k)_{1kn}`$ (when it exists) is a limit in 2-norm of cyclic gradients of polynomials in the noncommuting variables $`X_1,\mathrm{},X_n`$ ? The theorem we proved here provides a partial affirmative answer:
> If $`(𝒥_k)_{1kn}`$ are noncommutative polynomials in $`X_1,\mathrm{},X_n`$, then there is a noncommutative polynomial $`P`$ in $`X_1,\mathrm{},X_n`$ such that $`𝒥_k=\delta _kP`$ $`(1kn)`$.
Indeed, by Corollary 5.12 in we have $`{\displaystyle \underset{k}{}}[𝒥_k,X_k]=0`$. Hence the commutator condition (ii) in the Theorem is satisfied.
Acknowledgments. This research was conducted by the author for the Clay Mathematics Institute. Partial support was also provided by National Science Foundation Grant DMS95–00308.
References
P.Biane, R.Speicher. Free diffusions, free entropy and free Fisher information, preprint.
T.Cabanal-Duvillard, A.Guionnet. Large deviations upper bounds and noncommutative entropies for some matrices ensembles, preprint.
G.-C.Rota, B.Sagan, P.R.Stein. A cyclic derivative in noncommutative algebra.
Journal of Algebra 64, 54–75 (1980).
D.Voiculescu. The analogues of entropy and of Fisher’s information measure in free probabilitiy theory, V: noncommutative Hilbert transforms. Invent. Math. 32, no. 1, 189–227 (1998).
D.Voiculescu. The analogues of entropy and of Fisher’s information measure in free probabilitiy theory, VI: liberation and mutual free information. Advances in Mathematics 146, 101–166 (1999).
D.Voiculescu. Lectures on free probability theory. Notes for a course at the Saint-Flour Summer School on Probability Theory, preprint (1998).
Department of Mathematics University of California Berkeley, California 94720–3840 dvv@math.berkeley.edu
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# X-ray emission from Lindroos binary systems
## 1 Introduction
Pre-main sequence (PMS) late-type stars are known to be X-ray sources (see Walter et al. 1988, Bouvier 1990 and Neuhäuser et al. 1995). In the evolution of these stars to the Main Sequence (MS), there is a state usually defined as Post-T Tauri Stars (PTTS). PTTS were first defined by Herbig (1978) as PMS stars more evolved than Classical T Tauri stars (CTTS) but still contracting to the MS. Given that the stage of CTTS is only a small fraction of the total time of contraction of low-mass stars to the MS, PTTS should be much more abundant than CTTS if star formation has been ongoing for a sufficiently long time. However, it is difficult to find PTTS because they do not show spectroscopic or photometric peculiarities which make them easy to detect. Unlike CTTS, they do not present significant IR or UV excesses and the H$`\alpha `$ line is not observed as a strong emission line. Therefore, the identification of these stars is difficult and relies on the detection of the Li I (6708Å) absorption line and the chromospheric Ca II (H & K) lines in their spectra, as well as on their X-ray detection.
Murphy (1969) first proposed that PTTS could be searched as members of young binary systems. The basis of his idea was that the MS lifetime of high-mass stars is comparable to the contraction timescale of solar-type stars. Hence, there could be binary systems comprised of MS early-type stars physically bound to PTTS. Gahm et al. (1983) carried out photometric and spectroscopic observations of visual double stars with early type primaries. The derived data, together with $`JHKL`$ observations, allowed Lindroos (1986; L86 hereafter) to identify 78 likely physical pairs with several PTTS candidates as secondaries.
A high lithium abundance and a high chromospheric activity level are necessary (although not sufficient) indicators of youth. Martín et al. (1992) and Pallavicini et al. (1992) carried out optical spectroscopy of the Lindroos late-type companions, detecting the Li I(6708Å) absorption line and the Ca II (H & K) emission lines in the spectra of several PTTS candidates. Ray et al. (1995) took this sample of “genuine” PTTS and looked for circumstellar matter around them. For this purpose, they analyzed the IRAS database (Point Source Catalogue and Faint Source Catalogue) and also searched for continuum 1.1 mm emission. While IR excesses were found for most of the sources, no mm dust continuum was detected (see also Gahm et al. 1994 and Jewitt 1994).
The X-ray emission from Lindroos binary systems was first studied by Schmitt et al. (1993). After the analysis of seven pairs comprised of late B-type stars and later-type companions, the main result was the detection of X-rays from both members of the pair. In the case of late-type stars it is well-known that they produce X-rays in their hot coronae. However, this is not the case of late-B type stars. Theoretically, early-type stars between B4 and A7 are not expected to be X-ray emitters: they do not possess the strong winds thought to be responsible of the X-ray emission in O- and early B-type stars (Lucy & White 1980), nor significant convection zones thought to be necessary to sustain a magnetic dynamo to power a corona. Although the detection of X-rays from late-B and early-A type stars have been reported by several authors (i.e. Caillault & Zoonematkermani 1989, Schmitt et al. 1993, Berghöfer & Schmitt 1994, Berghöfer et al. 1996, Simon et al. 1995, Panzera et al. 1999), there is no clear mechanism that explains the origin of this emission. The most accepted explanation is related to the presence of otherwise unknown unresolved late-type companions of these stars.
The aim of this paper is to study the X-ray emission of all L86 binary systems with PTTS candidates which were observed by ROSAT. For this purpose, we have selected binary systems with late-type stars as secondaries. We will study the X-ray emission from the PTTS candidates as well as the emission from the early-type stars. The characteristics of the sample are described in Sect. 2. Sect. 3 provides the details related to the source detection and identification. The processed X-ray data are analyzed in Sections 4, 5 and 6. The conclusions are drawn in Sect. 7.
## 2 The binary sample
The Lindroos catalogue (L86) contains 78 binary systems. We have selected those binaries in which the secondary member is a F, G, or K-type star (note that no M-type stars are present in the Lindroos Catalogue). Our final sample consists of 47 systems. Two of them (HD 113791 and HD 106983) were not included in L86 but in Gahm et al. (1983). We have also included the binary system HR 4796 (HD 109573, TWA 11), a possible member of the TW Hya association (see Webb et al. 1999) given that its stellar properties are in agreement with those of the Lindroos sample. This is the only binary system in our sample with an M-type star as a secondary. Note that we have not rejected those pairs that are classified as likely optical pairs by Pallavicini et al. (1992), having in view that the X-ray emission could provide important information related to the nature of these systems.
Stellar properties of our binary sample are shown in Table 1. The name of the source as well as the components of the binary system (A for the primary, B for the secondary and C or X for companions in multiple systems) are shown in column 1. The spectral types of both stars are given in column 2. Column 3 shows the optical position of the primary star while column 4 shows the projected separation between both members. The distance to the primary star, the visual magnitude of both components and the visual extinction to the pair are given in columns 5, 6 and 7. We have also included in column 8 an orientative ‘flag’ which is related to the nature of the binary system: according to the spectroscopic survey carried out by Pallavicini et al. (1992), the systems are classified as LO (likely optical pair), PP (probably physically bound) and CP (certainly physically bound).
## 3 ROSAT observations and data reduction
The X-ray telescope and the instrumentation onboard the ROSAT satellite are described in detail by Trümper (1983), Pfeffermann et al. (1988) and David et el. (1996). Two main detectors are available: the Position Sensitive Proportional Counter (PSPC), which can be used either in survey or in pointed mode, and the High Resolution Imager (HRI).
We have studied all the available data for our sample, i.e., PSPC (survey and pointed mode) and HRI data. The source detection and identification have been carried out using the source detection routines provided by the Extended Scientific Analysis System (EXSAS; Zimmermann et al. 1997) which are based on a Maximum Likelihood (ML) technique (Cruddace et al. 1988). For all detected X-ray sources we have looked for optical counterparts to check the reliability of the detection. We will briefly describe the data in the following subsections.
### 3.1 PSPC observations
The ROSAT All-Sky Survey (RASS) was performed with the PSPC. The diameter of the field of view is 2° and each object is observed up to $`30`$ times separated by $`90`$ minutes, with up to $``$ 30 sec per scan.
The spectral resolution of the PSPC (43% at 0.93 keV) allows spectral analysis in three energy bands:
\- Soft = 0.1 to 0.4 keV
\- Hard 1 = 0.5 to 0.9 keV
\- Hard 2 = 0.9 to 2.0 keV
While in some pointed observations the signal to noise are large enough to carry out detailed spectral analysis, this is not possible for RASS data. Note that the RASS exposure times of our sample range from 75 sec to 1236 sec. However, we can obtain spectral information of our sources studying the X-ray hardness ratios (HR) defined as follows:
$$HR1=\frac{(H1+H2S)}{(H1+H2+S)}\text{and}HR2=\frac{(H2H1)}{(H2+H1)}$$
(1)
where H1 and H2 are the counts observed in the Hard 1 and Hard 2 bands, and S are the counts observed in the Soft band. Hence, HR values can range from $`1`$ to $`+1`$. If no counts are detected in one of the bands only an upper or lower limit to HR1 and HR2 is available. Neither HR values nor limits are available for undetected stars.
The nominal positional accuracy of the ROSAT PSPC detector in pointing mode is $``$ 25″ at 1 keV (note that it is energy and off-axis angle dependent). However, this accuracy is reduced to $``$ 1′ in the survey phase. If the offset between the X-ray source and the optical position exceeds the spatial resolution of the RASS data ($``$ 1′), the identification of the X-ray source with the optical counterpart is doubtful. Hence, for all the RASS observations we have to select a maximal displacement between the X-ray and the optical position within which the detections are reliable. Following Neuhäuser et al. (1995), we have taken a value of 40″ as the maximal distance between the X-ray detection and the optical counterpart.
The RASS data (detections and upper limits) of our sample are shown in Table 2. The name of the source and the position of the X-ray detection are given in columns 1 and 2. As mentioned before, the spatial resolution of the PSPC is too low to resolve most of the binary systems. In fact, in most of the cases we have only obtained a single detection displaced from both components of the system. Hence, we show in column 3 the displacements of the X-ray detection with respect to the optical positions of both members of the binary system. There is only one system in which the projected separation is so large that the X-ray emission can be clearly attributed to the secondary star: HD 87901. The total number of counts in the broad band, the exposure time and the hardness ratios HR1 and HR2 are given in columns 4, 5, 6 and 7. In the case of non-detections we have computed the number of counts at the position of both components, A and B, respectively. An estimate of the probability of the detection by the ML procedure is given in the last column; a value of ML=5 corresponds to a 2.7$`\sigma `$ signal over the background. Note that 11 sources listed in Table 1 do not appear in Table 2. Most of them (HD 8803, HD 27638, HD 36013, HD 36779, HD 53755, HD 71510, HD 123445 and HD 162082) are not detected and are located so close to another X-ray sources that a computation of an upper limit is not possible. Two of them (HD 70309, HD 76566 and HD 120641) show bad quality exposure maps so no X-ray data can be derived.
The same information as in Table 2 is given in Table 3 for PSPC pointed observations. Only four of our binary systems were observed in PSPC pointings, with just two of them detected. Columns 1 and 2 provide the name of the system and the ROSAT Observation Request (ROR) number. The coordinates of the X-ray detection are given in column 3. As in the case of the RASS data, the spatial resolution is not high enough to detect both components of the binary systems. As a consequence, we show the displacement $`\mathrm{\Delta }`$ of the X-ray source with respect to both components of the binary system in column 4. Column 5 shows the displacement of the source with respect to the center of the image (off-axis angle). Note that the sensitivity of the detector degrades with increasing off-axis angle. The total number of counts, the exposure time and the hardness ratios, HR1 and HR2, are given in columns 6, 7, 8 and 9. The ML coefficient is provided in the last column.
### 3.2 HRI observations
The HRI detector allows high resolution imaging of X-rays sources. The nominal spatial resolution of the detector is 1.7″ but a bore-sight correction as much as 10″ should be applied. The projected separations between the members of our binary systems range from 4.4″ to 177″. Most of them show separations larger than 10″ so, in principle, they can be resolved by the HRI.
The spectral resolution of the HRI is worse than that of the PSPC. However, the HRI pulse height distribution can be used to compute a 2-band (H and S) hardness ratio. We have defined the HRI hardness ratio as:
$$HR=\frac{HS}{H+S}$$
(2)
To obtain the most realiable S- and H-band definitions, several ROSAT sources known from their PSPC pointed observations to be either extremely soft (HR1 $`<`$ -0.8) or extremely hard (HR1 $`>`$ +0.8) were selected ($``$20). For these sources the HRI hardness ratios were calculated according to Eq. (2) for the following different soft and hard band definitions: S1 = 0-5, H1 = 6-15; S2 = 0-4, H2 = 5-15; S3 = 0-3, H3 = 4-15; and S4 = 0-2, H4 = 3-15, where the numbers give the corresponding HRI energy channels. For each of the corresponding four HRI hardness ratios the mean difference to the PSPC hardness ratios (HR1), was computed. The HRI bands for which the selected sources show the smallest deviation to the PSPC hardness ratio were S3,H3. Therefore, the HRI bands are defined as the soft band (S-band) corresponding to channels 0 to 3 (E $``$ 0.3 keV) and the hard band (H-band) corresponding to channels 4 to 15 (E $``$ 0.3 keV). This HRI hardness ratio shows a smaller dynamical range than the PSPC hardness ratio, HR1, but clearly identifies soft sources with negative values and hard sources with positive values. For more details, see Supper et al. (in prep.).
Table 4 shows the HRI detections of Lindroos binaries. Most of the HRI pointed observations were carried out by one of us (HZ) as Principle Investigator (PI). Columns 1 to 3 give the name of the source, the ROR number, and the component of the system detected: A+B when both stars are not resolved and there is only a single detection, and A or B (or C or X for secondaries in multiple systems) when the members of the binary system have been resolved. We have also added an ’e’ to designate the elongated sources. The position of the X-ray sources are shown in column 4. Column 5 shows the difference between the optical and X-ray position. Note that the optical coordinates used as a reference for the B-components are those of the secondaries (taken from SIMBAD database). Column 6 shows the displacement of the sources with respect to the axis of the telescope. Finally, the total counts in the broad band, the exposure time, the hardness ratio and the ML coefficient are given in columns 7, 8, 9 and 10. The upper limits of the undetected sources are also shown.
Table 4 includes three binary systems (HD 86388, HD 109573 and HD 113703) for which it is not clear which of the two members of the system is responsible for the X-ray emission. After the comparison of the optical and the X-ray position, we believe that the late-type companion is most likely the X-ray emitter in the three cases. This is clearer for the HD 86388 system because the X-ray detection is not elongated and it is closer to the optical position of the secondary. In the case of HD 109573 and HD 113703 the X-ray detections in the broad band are slightly elongated although much closer to the optical position of the secondary star (see Table 4). Hence, we will assume that the X-ray emission comes from the late-type secondaries in the three cases under study.
In order to illustrate different observations of the Lindroos binary systems, we have shown in Fig. 1 several HRI images of different pairs.
### 3.3 X-ray fluxes and luminosities
The X-rays fluxes can be computed by multiplying the observed count rates by an energy conversion factor (ECF). The ECF depends on the detector response and the underlying model for the X-ray spectrum. In the case of our sample we have assumed a 1-T Raymond-Smith thermal spectrum (Raymond & Smith 1977), which implies that the ECF mainly depends on the temperature of the emitting plasma and the interstellar absorption. The assumed temperature for the late-type stars is $`kT_x=1\mathrm{keV}`$ (with $`k=`$ Boltzmann’s constant) which is suitable for active late-type stars such as TTS (see Neuhäuser et al. 1995). In the case of the B-type stars we have adopted a mean value of $`kT_x=0.5\mathrm{keV}`$ (see Berghöfer & Schmitt 1994 and Berghöfer et al. 1996).
In order to correct for the interstellar absorption, we have converted the visual extinction $`A_v`$ to our sources into hydrogen column densities, $`N_H`$, following Paresce (1984):
$$\frac{N_H}{\mathrm{cm}^2}=5.510^{21}\frac{E(BV)}{mag}=\frac{5.5}{3.1}10^{21}\frac{A_v}{mag}$$
(3)
For those cases for which no $`A_v`$ is available, we have adopted a lower limit of log $`(N_H/cm^2)=18`$.
The PSPC ECF for different temperatures and interstellar absorption column densities are provided by Neuhäuser et al. (1995) for late-type stars, and by Berghöfer et al. (1996) for early-type stars. In the case of the HRI observations, we have computed the ECF’s following the Technical Appendix to the ROSAT Call for proposals.
Once we obtain the X-ray fluxes, the X-ray luminosities are given by
$$L_x=4\pi d^2f_x$$
(4)
with $`d`$ being the distance to the star. We have made use of the Hipparcos parallaxes of the primary stars to estimate the distances to our sources (see Table 1).
The X-ray fluxes and luminosities for the whole sample are given in Table 5. The HD number and the binary component (A, B, C or X) are shown in columns 1 and 2. Column 3 provides the computed hydrogen column density. The type of observation, R for RASS, P for PSPC and H for HRI, is given in column 4. The X-ray luminosities with their respective errors and the upper limits for non-detections are finally listed in the last column. If the binary system is observed and resolved by the HRI, we provide the fluxes and luminosities for both members of the pair. If unresolved, we just can provide one X-ray luminosity associated with the pair. In these cases we have adopted a temperature of 1 $`keV`$ to compute the ECF’s although, in principle, both stars could be contributing to the total emission. The last method has also been applied to those systems with only RASS or PSPC pointed observations, given that they are always unresolved. In the case of RASS or PSPC non-detections, we have estimated two upper limit luminosities, corresponding to the count rates computed at the optical positions of the primary and the secondary star, respectively (see Table 2).
## 4 Interpretation of the data
The PSPC is not able to resolve most of the binary systems of the sample (the only exception is HD 87901). Therefore, the HRI data are the most appropriate to study the X-ray emission of individual stars in Lindroos systems. We have based the following analysis on the HRI data of resolved pairs. As we will show in Section 5, most of the resolved pairs show evidences of being physically bound. Only two binary systems are definitely optical pairs (HD 123445 and HD 127971), while another two systems are doubtful (HD 40494 and HD 87901).
### 4.1 X-ray luminosity functions
The X-ray luminosity function (XLF) of a sample can be derived with Kaplan-Meier estimators using the statistical package ASURV (see Feigelson & Nelson 1985, Schmitt, 1985 and Isobe et al. 1986), which allows to take into account both detections and upper limits.
To check whether the two samples of early-type primaries and late-type secondaries are statistically different we have performed a two-sample test with ASURV. We do not include those objects which are either unresolved or whose identification is not clear, namely: HD 560 and HD 1438, HD 33802, HD 109573, HD 113703 and HD 86388. The result of the test is that our samples are statistically similar with a probability of 0.8 (see upper panel of Fig. 2). Given the different spectral types of the stars in the two samples, this conclusion may appear surprising. In order to check this result, we have repeated the statistical analysis after removing those sources that could be contaminating both samples. In the case of the late-type stars we have removed the X-ray upper limits of two probably unbound sources, HD 123445 B, HD 127971 B. In the case of early-type stars we have removed one source with an unreliable X-ray detection ( HD 87901 A, see Sect. 4.2), and three sources for which we found indications of unresolved late-type companions (Sect. 4.3): HD 123445 A, HD 127971 A and HD 129791 A. For this reduced sample, the probability of both groups of stars to be statistically similar is reduced to 0.4 but this value is higher than the threshold (0.05) to reject the null hypothesis of our two samples to be equivalent. (see lower panel of Fig. 2).
Note, however, that the X-ray luminosity alone cannot provide information about the nature of the X-ray emission. This could only be studied through the spectral analysis of the emission or through the hardness ratios (HR’s).
### 4.2 Hardness ratios
In the case of the HRI observations, we could obtain individual hardness ratios (HR) using the two sets of channels defined in Sect. 3.2. As shown in Table 4, the computed hardness ratios generally differ from early- to late-type stars. In most of the cases, late-type stars show positive HR’s while early-type stars show negative values. This means that late-type stars emit most of their X-ray radiation in the H-band, which is consistent with the presence of an energetic corona, while early-type stars mainly emit in the S-band. In fact, there are several stars which are only detected in one of these bands, showing upper or lower limits to the HR’s. However, there are some stars which do not follow this trend: HD 87901A, HD 123445A, HD 127971A and HD 129791A, all of which are early-type stars so they are not supposed to have an energetic corona which could explain their positive HR’s.
Berghöfer et al. (1999) has shown that the HRI detector is sensitive to ultraviolet (UV) radiation below 4000Å. The contamination of the UV light to the final count rate is mainly concentrated to pulse height channels 1-3 ($``$ our S-band range). Our primary stars are bright sources in the UV range so, in principle, the computed S-band rate could be just a response of the detector to the UV light. In order to check the reliability of the S-band counts for the early-type stars, we have estimated the contribution of the photospheric UV light to the total S-band rate. Following Berghöfer et al. (1999):
$$HRI_{UV}=10^{(1.022\pm 0.003)(0.555\pm 0.005)U}(cts/s)$$
(5)
with U being the U-magnitudes of the observed stars. This equation was deduced considering the emission of the sources from channels 1 to 8 of the HRI detector. In order to be conservative, we have also computed the S-band rates in these channels. Table 6 shows the results. Column 1 provides the name of the star. The U-magnitudes, taken from The Bright Star Catalogue (Hoffleit & Jaschek 1991) are shown in column 2. For three sources U-magnitudes are not available, so we deduced them from their spectral types and V-magnitudes following Kenyon & Hartmann (1995). The S-band rate and the UV rate, both computed in channels 1-8 of the detector, are shown in columns 3 and 4, while their ratio is provided column 5. We finally show the derived $`\mathrm{lg}(L_\mathrm{x}/L_{\mathrm{bol}})`$ ratio in the last column, in order to check the reliability of the X-ray detections. The bolometric luminosities have been computed using the stellar data provided in Table 1 and the bolometric corrections from Schmidt-Kaler (1982).
As we can see from Table 6, most of the B-type stars of our sample have S-band rates significantly higher than the estimated UV-rates. There is only one source for which both values are comparable: HD 87901A. Note that the HRI observations of this star show very large off-axis values and, according to Berghöfer et al. (1999), Eq. (1) is not reliable for these cases. However, this system was resolved by the PSPC in the RASS survey (see Tab. 2) and, as we mentioned in Sect. 3, no emission was detected from the early-type star but from the secondary.
To further test the reliability of the S-band emission in early-type stars, we have also studied the $`\mathrm{lg}(L_\mathrm{x}/L_{\mathrm{bol}})`$ ratio of these stars to see if it is consistent with the ratio for stars of similar spectral types. According to Table 6, most of the stars show ratios consistent with the reported values for early-type stars (Berghöfer et al. 1997). There are also four stars with unusual ratios: HD 87901 with a lower ratio, and HD 123445, HD 127971 and HD 129791, with ratios closer to those found in late-type stars. Given that these three early-type stars also show positive HR, it is possible that they have unresolved late-type companions. However, in the case of HD 87901 the HR is negative and the $`\mathrm{lg}(L_\mathrm{x}/L_{\mathrm{bol}})`$ is lower than in late-type stars, so the X-ray emission can not be related with an unresolved source. This fact together with its non-detection in the RASS make us think that this emission is due to the HRI UV leak. Therefore, excluding HD 87901 A, we can identify the computed S-band rates with intrinsic X-ray emission from the sources.
The previous test allows to confirm the reliability of the computed HRI HR’s. Given that these HR’s are systematically positive for late-type stars and systematically negative for early-type stars, we can conclude that the nature of the X-ray emission is intrinsically different for our two samples. We will discuss this point more deeply in the following subsection.
### 4.3 X-ray luminosities
The analysis of the HR’s suggests that the nature of the X-ray emission is different for our two samples. A useful way to represent this difference consists of comparing the X-ray luminosities obtained in the two HRI bands, S and H, for all of our detected sources. Fig. 3 shows this comparison. Note that we have converted the S- and H-band rates into luminosities following the procedure described in Sect. 3.3. We have not included HD 33802 because the binary is not completely resolved and it is difficult to obtain reliable measurements from the individual members (see Fig. 1).
As can be seen from Fig. 3, there is a clear separation between early- and late-type stars. Early-type stars are generally softer than the late-type companions. This result is in agreement with that obtained from the HR’s analysis. Note that although most of the early-type stars lie below the 1:1 correlation line, there are three early-type primaries with higher X-ray luminosities in the H-band than in the S-band: HD 123445 A, HD 127971 A and HD 129791 A. As we explained before, the higher H-band luminosities obtained for the later-type stars are in agreement with their condition of active stars (note that all are located above or on the 1:1 correlation line). In contrast, the higher H-band luminosities for the early type stars are not easy to explain, given that these stars are not supposed to have a convective zone able to support a corona. Because these stars lie in the same part of the diagram as the late-type secondaries, the simplest explanation is to relate them to unresolved late-type companions.
After the analysis of the X-ray luminosities in the two bands, we have studied the total X-ray luminosities of the whole sample. Given that the X-ray emission typically assumes a particular value for each spectral type, we have plotted the X-ray luminosity against the bolometric luminosity for each star. The bolometric corrections were adopted from Schmidt-Kaler (1982). Only in the case of the M-type star, HD 109573 B (HR4796 B), the bolometric correction was taken from Kenyon & Hartmann (1995).
As shown in Fig. 4, there is a clear separation between late- and early-type stars. Late-type stars lie in the region of $`\mathrm{lg}(L_{\mathrm{bol}}/erg/s)3334`$ with X-ray luminosities varying between $`\mathrm{lg}(L_\mathrm{x}/erg/s)`$27.5 and 31. On the other hand, the early-type primaries show lower X-ray luminosities that decrease from B1 to A0 spectral types. However, the three late-B type stars which clearly deviate in Fig. 3 also deviate in this figure, showing a higher X-ray emission than expected for their spectral types: B9, B7 and B9.5.
We have compared these results with those previously obtained by Schmitt et al. (1993). The X-ray luminosities derived for their sample of 7 Lindroos systems agree with our results except in one case: HD 113703. While these authors identify the single HRI X-ray detection with the early-type primary, we think it is most probably related to the late-type secondary. As we discussed in Sect. 3, the difference between the X-ray detection and the optical position is smaller for the late-type secondary in both the broad and hard band images.
At this point we must also remark that the young star HD 109573 B (HR4796 B) shows an X-ray luminosity higher than that reported by Jura et al. (1998). These authors compute the X-ray luminosity from the same HRI image but using an ECF that corresponds to the PSPC detector, although these detectors have different sensitivities. Moreover, they take a mean ECF from Neuhäuser et al. (1995) which is deduced from a ROSAT survey on Taurus, a star forming region where most of the stars show visual extinctions larger than those in TWA. Therefore, we think that the value listed in Table 5 is more realistic since it takes into account the ECF from the HRI detector and a negligible absorption to the source.
One of our binary systems, HD 560, was detected but unresolved by the ROSAT HRI (see Fig. 1). Because this system is comprised of a B9 primary and a G5 type companion, we have made a simple test to confirm the results obtained for HD 123445 A, HD 127971 A and HD 129791 A: although we certainly know that HD 560 is a binary system, we have supposed HD 560 to be a single B9 star and not a pair. Therefore, if we assume that the single X-ray detection corresponds to a B9-type star and we plot this source into Fig. 4 (starred symbol), the result is that HD 560 lies in the same part of the diagram as the three sources with unresolved late-type candidates. Hence, this test strengthens the idea of unresolved late-type companions in these late-B type stars.
We have compared our sample of primary stars with a sample of MS B-stars taken from the The RASS catalogue of optically bright OB-type stars (Berghöfer et al., 1996) in Figure 4. In principle, our primary stars are in good agreement with the sample of B-type stars with comparable X-ray luminosities. Note, however, the large scatter in the X-ray values at $`\mathrm{lg}(L_{\mathrm{bol}}/erg/s)`$ 35.5 (B9 stars). HD 123445 A, HD 127971 A and HD 129791 A lie in this region of the diagram with X-ray luminosities ranging from $`10^{29.5}`$ to $`10^{30.8}`$ erg/s.
The wide range of X-ray luminosities found in late-B stars from Berghöfer’s sample together with the position of our late-B primary stars in Fig. 4, suggests the possibility of having unresolved late-type companions in late-B type stars with $`\mathrm{lg}(L_\mathrm{x}/erg/s)>`$29.5. Although it is beyond the scope of this paper, it would be interesting to check if the MS late-B type stars from Berghöfer et al. (1996) with highest X-ray luminosities are suspected to have late-type unresolved companions.
As we mentioned before, the $`\mathrm{lg}(L_\mathrm{x}/L_{\mathrm{bol}})`$ ratio is generally similar for stars of the same spectral type. As an example, a “canonical” relation of $`(L_\mathrm{x}/L_{\mathrm{bol}})10^7`$ has been reported for O- and early B-type stars (Harnden et al. 1979, Long & White 1980, Pallavicini et al. 1981, Sciortino et al. 1990). Given that the $`(L_\mathrm{x}/L_{\mathrm{bol}})`$ ratio is characteristic for each spectral type, we have finally compared this value with the HRI HR’s computed in Sect. 4.2.
Fig. 5 shows a clear separation between the stars of our sample. A boundary line drawn at $`HR=0`$ allows us to classify our stars in soft and hard X-ray emitters. While the softer X-ray emitters lie close to the $`\mathrm{lg}(L_\mathrm{x}/L_{\mathrm{bol}})7`$ ratio reported for O- and early-B type stars, those with positive HR, are spread over a wide range of $`\mathrm{lg}(L_\mathrm{x}/L_{\mathrm{bol}})`$ values. Among the latter, we can also distinguish two groups. On one hand, we find the late-type secondaries lying at the top right corner of the figure. On the other hand, the three late-B type stars with possible unresolved late-type companions occupy a band with $`\mathrm{lg}(L_\mathrm{x}/L_{\mathrm{bol}})`$ ranging from -4.5 to -6. Note that these stars clearly deviate from the group of early-type primaries at the left bottom corner of the figure.
If we include HD 560 in Fig. 5 making the same assumptions as in Fig. 4, ie. assuming that it is a single B9 star responsible for the detected X-ray emission, we can see that the source also lies in the same region as HD 123445 A, HD 127971 A and HD 129791 A.
As a conclusion, we can confirm that B-type stars in the Lindroos systems under study generally show a decrease in the X-ray luminosity for decreasing spectral types (from B0 to B9). The $`(L_\mathrm{x}/L_{\mathrm{bol}})`$ ratio is in agreement with their spectral types and it is well-correlated with their negative HR’s. However, there are three sources, HD 123445 A, HD 127971 A and HD 129791 A which show $`L_\mathrm{x}`$ values higher than those reported for earlier B-type stars. Moreover, when comparing the $`(L_\mathrm{x}/L_{\mathrm{bol}})`$ ratio with the computed HRI HR’s, these three late-B type stars clearly deviate from the sample of primary stars, showing values closer to those reported for the late-type secondaries. Therefore, these three late-B type stars are suspected to have unresolved late-type companions.
## 5 X-ray emission from late-type secondaries: are they Post T Tauri stars?
As mentioned in the introduction, the late-type secondaries of Lindroos systems have been studied in several spectral ranges. In particular, the optical and IR data provide clear evidences of youth among these stars.
In Table 7 we ahow some of the main observational properties of these late-type secondaries. Columns 1, 2 and 3 provide the name of the source, its spectral type and the X-ray luminosity derived from this work. Column 4 shows the equivalent width (EW) of the Li I absorption line while IR data from the IRAS satellite are provided in columns 5 and 6. A ‘flag’ related to the measured radial velocity of the pair is given in column 7. We have finally summarized all these data in the last columns of the table in order to isolate those late-type stars with evidences to be physically bound to the early-type primaries. Apart from the measured radial velocity of the pair and the Li I EW we have also considered other indicators of youth like the X-ray emission and the measured IR excesses.
As we can see from Table 7, all the late-type secondaries in binary systems with similar radial velocities show indications of youth. For the stars with no radial velocity measurements, we have found pairs with clear evidences of youth and hence, most probably physically bound to their primaries: HD 17543, HD 53191, HD 60102, HD 129791 and HD 143939. HD 123445 B and HD 127971 B do not seem to be bound to their primaries because they lack indicators of youth. Finally, HD 40494 B is not detected in X-rays but it shows a strong Li I absorption line (as HD 27638). Therefore, it is not obvious how to classify this source. Radial velocity measurements would be convenient to confirm the nature of this pair.
According to Lindroos (1985), most of these stars show ages lower than 70 Myr but, as discussed in L86, the uncertainties associated with the age determination are large. Note that the ages of the Lindroos systems were first determined from uvby$`\beta `$ photometry of the primaries (Lindroos 1985) and making use of the evolutionary tracks and isochrones by Hejlesen (1980). A new estimation of the ages using modern isochrones plus new optical and IR photometrical data will be postponed to a later publication.
Instead of studying the ages of the systems we have considered the Li I (6708Å) EW and we have related it to the derived X-ray luminosities. As shown in Table 7, the Li I EW has been directly measured in most of the Lindroos secondaries, so we can study if there is a relation between both parameters. We have plotted in Fig. 6 the Li I EW versus the effective temperature of our stars. Note that these effective temperatures have been derived from their spectral types according to the conversion given by Kenyon & Hartmann (1995). We have also considered the upper envelopes to the Li I EW of two young clusters, the Pleiades ($``$ 125 Myr) and IC 2602 ($``$ 35 Myr), so that those secondary stars with Li I EW stronger than in the Pleiades or in IC 2602 are younger than these clusters (data from Soderblom et al. 1993, Randich et al. 1997 and Stauffer et al. 1997).
Fig. 6a shows that most of the strongest X-ray emitters are located above the Pleiades upper envelopes to the Li I. This suggests that a large fraction of the Lindroos secondaries are younger than 125 Myr, with 4 of them showing ages lower than 35 Myr. Among these 4 stars, we can clearly distinguish HD 109573 B (HR 4796 B) in the right middle part of the figure, with a high value of the Li I EW. Its X-ray emission is not as large as this coming from the other three sources above the IC2602 Li I envelope. The different spectral type as well as the different evolutionary state could explain this difference.
Fig. 6a also allows us to confirm the nature of the HD 123445 and HD 127971 pairs. The secondary stars lie at the bottom of the diagram with upper limits in both, the Li I EW and the X-ray emission. Two more stars show non-detections in X-rays although they lie very close to the upper envelope of the Pleiades: HD 27638 and HD 40494. As we have discussed above, the former is probably bound to the primary (see table 7) while the nature of the later one is not still clear.
As a conclusion, we can say that those stars with strong X-ray emission seem to be younger than IC 2602. Note, however that although the uncertainties associated with the Li I EW measurements are generally low (see Table 7), for the stars earlier than $``$K0 the Li I EW does not allow us to reach any conclusion about their ages because F- and G-type PMS together with ZAMS stars still have their initial Li content.
To study the last issue in detail, we have converted the Li I EW into lithium abundances, N(Li) making use of non-LTE curves of growth from Pavlenko & Magazzú (1996) and assuming a surface gravity of $`\mathrm{lg}g_s=4.5`$. Note that we have only considered high-resolution data to make this conversion (see Table 7). We have plotted this value against the effective temperature of our stars in Fig. 6b. As in the left panel of the figure, we have found strong X-ray emitters above the Pleiades and IC 2602 data. Note that most of the Lindroos secondaries (11) show lithium abundances close to 3.0 which is consistent with low mass PMS stars.
## 6 Results and Conclusions
We have reported the X-ray emission from Lindroos binary systems observed with ROSAT. Most of the stars from this sample are detected in the RASS and several binary systems were observed with the PSPC and the HRI. After the analysis of the data, the main conclusions are:
1. Both, early- and late-type stars, show the same distribution of the X-ray Luminosity Function (XLF). Both samples are statistically similar with a probability of a 0.8. When we repeat the analysis without including stars that could be contaminating the samples, that is, without considering most certainly unbound secondaries, early-type primaries with possible unresolved late-type companions and one unreliable X-ray detection, the result does not change. The probability is reduced to 0.4 but it is large enough to statistically confirm the similarity of both samples.
2. A careful study of the HRI hardness ratios of our sample allows to discriminate between soft and hard X-ray emitters. Late-type stars always show positive HR while most of the early-type stars show negative values. Moreover, if we compare the X-ray luminosity in the two energy bands of the HRI, we can see that there is a clear separation between the two samples. Late-type stars show higher H-band than S-band X-ray luminosities. This result can be explained in terms of an energetic corona and activity episodes commonly reported in these stars. In the case of early-type stars, most of them present higher S-band than H-band luminosities. However, there are three X-ray sources formally identified as early-type primaries which show H-band luminosities comparable to those of late-type stars: HD 123445 A, HD 127971 A and HD 129791 A.
3. The $`L_\mathrm{x}L_{\mathrm{bol}}`$ diagram shows a clear separation among sources of different spectral types. This separation is even more prononced when the $`L_\mathrm{x}/L_{\mathrm{bol}}`$ ratio is plotted vs. the HRI hardness ratio. While early-type stars, which generally show negative HR’s, lie close to the $`L_\mathrm{x}/L_{\mathrm{bol}}=\mathrm{\hspace{0.17em}10}^7`$ “canonical” relation, late-type stars show positive HR’s and present higher values of this ratio ($`L_\mathrm{x}/L_{\mathrm{bol}}10^3`$).
For those particular cases of early-type stars with positive HR’s, we have seen that they display X-ray luminosities comparable to those of the late-type companions. When they are plotted into the $`L_\mathrm{x}/L_{\mathrm{bol}}`$ \- HRI HR diagram, they lie closer to the late-type stars than to the B-type group. All these evidences make us conclude that HD 123445 A, HD 127971 A and HD 129791 A are good candidates to have unresolved late-type companions.
5. The computed X-ray luminosities together with the Li I (6708Å) EW and abundances deduced for the Lindroos late-type secondaries have revealed this group of stars to be a good sample of PTTS candidates. A strong X-ray emission is reported for the youngest ones. Although we have no reliable measurements of their ages, most of the Lindroos secondaries seem to be younger than 125 Myr (when compared to the Pleiades) with three of them showing ages lower than 35 Myr (IC 2602 cluster). Also HD 109573 B (HR 4796 B) is located above these two clusters showing a high X-ray emission.
This study also has allowed us to confirm the optical nature of the HD 123445 and HD 127971 pairs. In this sense, two more systems are classified as doubtful: HD 87901 and HD 40494. The former lacks clear indicators of youth while HD 40494, although not detected in X-rays, may be a young star bound to its primary.
###### Acknowledgements.
We would like to thank T. Hearty, M. Fernández, J. Alves, S. Döbereiner and B. König for their assistance and their useful comments. We are very grateful to T. Berghöfer for providing his data. Some archived ROSAT observations investigated here were performed by T. Berghöfer, J. Krautter, J. Puls, T. Simon and J. Staubert as PI’s. The ROSAT project is supported by the Max-Planck-Society and the German Government (DLR/BMBF).
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# Localization Transition in Multilayered Disordered Systems
## I introduction
The understanding of the Anderson transition based on the scaling theory of localization inspired many detailed numerical studies of disordered electronic systems . The universality of the associated critical behavior was tested for various physical models, which include the crucial role of symmetry and added magnetic field . Universal critical transport properties are also expected in the presence of hopping matrix elements which are not the same in the various lattice directions, as it can be seen from computations for weakly coupled chains and coupled planes . It must be emphasized that many of the previous works on anisotropy include site diagonal isotropic disorder and involve anisotropy only in the hopping magnitudes. This kind of anisotropy remains for zero disorder and is manifested in the band structure. However, many realistic materials involve truly anisotropic disorder. For example, attempting to understand the high-$`T_c`$ cuprates within a non-interacting electron picture in the presence of disorder requires the explanation of the contrasting resistivities in the parallel and perpendicular directions . Anisotropic site randomness in a form resembling a random superlattice with lateral inhomogeneities gave anisotropic localization for anisotropy below a critical value, even for arbitrarily small disorder .
The in–plane resistivity for most of the layered high$`T_c`$ materials exhibits metallic behavior, increasing linearly with temperature over a wide temperature range, while the perpendicular out–of–plane resistivity is very high at low temperatures and decreases rapidly as the temperature increases, reminiscent of semiconductors . The contrasting behavior of the parallel and the perpendicular resistivities was observed in Bi<sub>2</sub>Sr<sub>2-x</sub>La<sub>x</sub>CuO<sub>y</sub> far below $`T_c`$, down to the lowest experimental temperature . In the underdoped La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> logarithmic divergencies of the corresponding resistivities accompanied by a nearly constant anisotropy ratio are, instead, observed suggesting an unusual three-dimensional ($`3D`$) insulator. The electronic transport in these materials is expected to arise from scattering in the “insulating” layer between the conducting CuO<sub>2</sub> layers. On the other hand, in almost all high-$`T_c`$ cuprates doping impurities or oxygen vacancies occupy the insulating layers between the conducting “pure” CuO<sub>2</sub> planes which implies interlayer disorder. Although the main aspects of transport in high-$`T_c`$ materials, such as the linear temperature dependence of the in-plane resistivity, are intimately connected with their strongly correlated nature, it is believed that anisotropic transport issues are, somehow, related to their layered structure.
We propose a strongly anisotropic multilayered structure (see Fig. 1) motivated by realistic anisotropic materials. This system involves truly anisotropic interlayer disorder, anisotropic hoppings and the usual isotropic site diagonal disorder. It can be regarded as a very simple model for the cuprates where the CuO<sub>2</sub> planes are believed to be identical without superlattice–like disorder. Our aim is to study both parallel ($``$) and perpendicular ($``$) transport addressing the following main questions: (1) does anisotropic localization occur (for example, localization in the layering direction and delocalization within the layers) for interlayer disorder only? (2) with additional isotropic diagonal disorder is the critical behavior independent of the direction as scaling theory predicts? Firstly, we compute the conductance in the case of interlayer disorder alone to check whether its localization behavior is the same in both directions. Secondly, in the presence of additional isotropic disorder of strength $`W`$ we obtain the critical disorder $`W_c`$ and the localization length critical exponent $`\nu `$ to see if they depend on the direction. We have also analyzed the statistical properties of the critical conductance distributions $`P_c(g)`$ and although we can conjecture that it is a unique single-parameter function the distribution for the logarithm of the critical $`g`$ in the “difficult” $``$ case resembles an insulator.
In Section II we proceed with the definition of the tight-binding Hamiltonian for the multilayered lattice structure. In Section III, we consider the system with only interlayer disorder and estimate the corresponding mean free paths in the two directions. Our results in the absence of isotropic diagonal disorder allow to conclude, in agreement with the scaling theory, that the states remain extended in both directions despite the strongly anisotropic interlayer disorder. However, the metallic conductance is very different for the $``$ case, being insulator-like. In Section IV we review the numerical methods for the computation of the conductance in cubic and long wire systems. We find singular behavior along the layering direction due to the missing bonds. In order to avoid this problem we have developed appropriate numerical algorithms based on transfer matrix and Green function methods. Finally, in Section V we discuss the conclusions of the present study also in connection to realistic systems.
## II random multilayer lattice
We propose a simple $`3D`$ anisotropic multilayered model which consists of parallel lattice planes randomly connected by interplane bonds as in Fig. 1 described by the Hamiltonian
$`H`$ $`=`$ $`{\displaystyle \underset{𝐦l}{}}\epsilon _{𝐦,l}|𝐦,l><𝐦,l|+{\displaystyle \underset{𝐦,𝐦^{},l}{}}(|𝐦,l><𝐦^{},l|`$ (2)
$`+t_{𝐦,l}^{}|𝐦,l><𝐦,l+1|+\text{H. c.}),`$
where $`𝐦`$, $`𝐦^{}`$ denote the two-dimensional site indices in each layer and $`l`$ is the layer index. The first term in Eq. (1) describes diagonal (isotropic) disorder with the site matrix elements $`\epsilon `$ chosen randomly from a box distribution within $`[W/2,W/2]`$, the second term describes nearest-neighbor hopping of unit strength within the layers, which sets the energy unit, and the third term corresponds to interplane hoppings $`t_{𝐦,l}^{}=0`$ or $`t`$, placed with probability $`p`$ at random layer positions $`𝐦`$. The interplane term obeys the binary distribution
$$P(t_{𝐦,l}^{})=p\delta (t_{𝐦,l}^{},t)+(1p)\delta (t_{𝐦,l}^{},0).$$
(3)
The proposed structure has both anisotropic hoppings due to $`t`$ and anisotropic disorder due to $`p`$. The transport characteristics are obtained by calculating the conductance along the $``$ and the $``$ directions.
The missing perpendicular bonds in the layering direction disturb particle migration even the presence of interlayer disorder alone. Unlike a naive expectation we find no critical point when $`W=0`$, for any $`p0`$. We show that the system is metallic independently of the direction, although the behavior of the conductance is very different in the two directions. In the presence of additional diagonal disorder, denoted by $`W`$, a critical disorder $`W_c`$ is obtained for various choices of the density $`p`$ and strength $`t`$. The critical point $`W_c`$ within finite size errors is found to be the same in both directions.
## III mean free paths for interlayer disorder
We consider anisotropic disorder in the perpendicular layering direction represented by $`p`$ and $`t`$, due to the randomly placed bonds among consecutive layers, in the absence of diagonal disorder $`W`$. In this disordered anisotropic lattice one might expect transport to be hindered in the perpendicular direction. It is worth examining whether is present or not. In order to proceed we adopt a convenient layer-diagonal representation since for $`W=0`$ the $`2D`$ layers are perfect planes and can be easily diagonalised. The eigenstates at the $`l`$th layer $`|𝐤_{},l`$ are labelled by the parallel momentum $`𝐤_{}`$ and the Hamiltonian $`H`$ can be expressed in the convenient Bloch-Wannier basis
$$|𝐤_{},l=\frac{1}{\sqrt{N_{}}}\underset{𝐦}{}e^{i𝐤_{}𝐦}|𝐦,l,$$
(4)
with parallel momentum $`𝐤_{}`$, the layer index $`l`$ and $`𝐦`$ summed over all $`N_{}=L^2`$ sites in every layer for a systen with $`L^3`$ sites. For $`p=0`$ the $`2D`$ layers are perfect and $`𝐤_{}`$ is a good quantum number. For $`p=1`$ the system reduces to a perfect $`3D`$ lattice and both $`𝐤_{}`$, $`k_{}`$ become good quantum numbers.
We consider the case of $`p0,1`$ where the translational symmetry in the plane directions is also broken and $`𝐤_{}`$ is no longer a good quantum number. In this mixed representation the Hamiltonian $`H`$ can be expressed as
$`H`$ $`=`$ $`{\displaystyle \underset{𝐤_{},l}{}}ϵ_{}(𝐤_{})|𝐤_{},l𝐤_{},l|`$ (6)
$`+{\displaystyle \underset{l}{}}{\displaystyle \underset{𝐤_{},𝐤_{}^{}}{}}[t_{l,l+1}^{𝐤_{},𝐤_{}^{}}|𝐤_{},l𝐤_{}^{},l+1|+\text{H. c.}],`$
with the parallel kinetic energy
$$ϵ_{}(𝐤_{})=2\mathrm{cos}(k_x)+2\mathrm{cos}(k_y)$$
(7)
and the hopping matrix element between neighboring planes
$$t_{l,l+1}^{𝐤_{},𝐤_{}^{}}=\frac{1}{N_{}}\underset{𝐦}{}e^{i(𝐤_{}^{}𝐤_{})𝐦}t_{l,𝐦}^{}.$$
(8)
In order to investigate the question of localization in the layering direction we define the retarded Green function
$$G(𝐤_{},l;𝐤_{}^{},l^{};t)=i\theta (t)[c_{𝐤_{},l}(t),c_{𝐤_{}^{},l^{}}^{}(0)]_+,$$
(9)
where $`c_{𝐤_{},l}(t)`$ is the time-dependent destruction operator of electron in the state $`|𝐤_{},l`$. Its diagonal element $`G(𝐤_{},l;𝐤_{},l;t)`$ gives the probability for finding an electron on the layer $`l`$ with momentum $`𝐤_{}`$ at time $`t`$, if initially it was on the same layer having the same momentum. The Fourier transformation of the diagonal Green function with respect to time is
$$G(𝐤_{},l;𝐤_{},l;E)=\frac{1}{E\mathrm{\Sigma }(𝐤_{},l,E)}$$
(10)
with the self-energy $`\mathrm{\Sigma }_{}`$ written as
$$\mathrm{\Sigma }(𝐤_{},l,E)=ϵ_{}(𝐤_{})+\underset{n=1}{\overset{\mathrm{}}{}}\underset{j}{}T_j^{(n)}.$$
(11)
The $`n`$-th order term is a sum over all paths $`j`$ of length $`n`$ starting and ending in the same state $`|𝐤_{},l`$
$$T_j^{(n)}=t_{l,l_1}(𝐤_{},𝐤_1)\underset{i=1}{\overset{n}{}}\frac{t_{l_i,l_{i+1}}^{𝐤_i,𝐤_{i+1}}}{Eϵ_{}(𝐤_i)}$$
(12)
with $`𝐤_{n+1}=𝐤_{}`$ and $`l_{n+1}=l`$. From Eq. (3) one has
$$t_{l_i,l_{i+1}}^{𝐤_i,𝐤_{i+1}}=\frac{t}{N_{}}\underset{𝐦_{l_i,l_{i+1}}}{}e^{i(𝐤_i𝐤_{i+1})𝐦}\delta _{l_{i+1},l_i\pm 1},$$
(13)
with the sum for $`𝐦`$ over the set of sites $`_{l_i,l_{i+1}}`$ which have interplane bond connections. It is seen from Eq. (13) that the diagonal in momentum matrix elements ($`𝐤_i=𝐤_{i+1}`$) are exactly $`pt`$. For a given order $`n`$ if $`E`$ lies within the pure $`2D`$ band $`ϵ_{}(𝐤_{})`$ the most divergent term of Eq. (11) comes from the path which has intermediate states $`|𝐤_i,l_i`$ with $`ϵ_{}(𝐤_i)E`$ for all $`i`$. The corresponding term approaches $`\{pt/[Eϵ_{}(𝐤_i)]\}^n`$ which sequentially connects nearest neighbor plane states with the same $`𝐤_i`$. It can be also seen that momentum scattering is always accompanied by interlayer hopping since a change of $`𝐤_{}`$ leads to a change of the layer index.
The above analysis splits the Hamiltonian into two parts $`H=H_0+H_1`$. The “undisturbed” part $`H_0`$ represents a perfect anisotropic $`3D`$ lattice with intralayer (interlayer) hopping $`1`$ ($`pt`$). The rest is a “random” Hamiltonian $`H_1`$ with only off-diagonal matrix elements expressed in the $`H_0`$-diagonal basis $`|𝐤=|𝐤_{},k_z`$ of the form
$`H_1(𝐤,𝐤^{})`$ $`=`$ $`{\displaystyle \frac{t(e^{ik_z}+e^{ik_z^{}})}{N}}{\displaystyle \underset{l,𝐦_{l,l+1}}{}}e^{i(𝐤_{}𝐤_{}^{})𝐦+i(k_zk_z^{})l},`$ (15)
$`\text{for }𝐤_{}𝐤_{}^{},`$
$`H_1(𝐤,𝐤^{})`$ $`=`$ $`0,\text{ for }𝐤_{}=𝐤_{}^{},`$ (16)
where $`N`$ is the total number of lattice sites and $`k_z`$ the perpendicular momentum.
It can be seen from Eq. (15) that the matrix elements of $`H_1`$ are complex numbers of average amplitude value modulo 1 plus a random phase. If the size of the system increases to infinity the phase exhausts all possible values in $`[0,2\pi ]`$ and the average should vanish. In this situation the scattering by the random configurations of the interplane bonds can be well described by perturbation theory with $`𝐤`$-space self-energy
$$\mathrm{\Sigma }(𝐤,E)ϵ(𝐤)+\underset{𝐤^{}}{}\frac{|H_1(𝐤,𝐤^{})|^2}{Eϵ(𝐤^{})i0^+}.$$
(17)
The configuration average $`\text{Im}_𝐤^{}\frac{|H_1(𝐤,𝐤^{})|^2}{Eϵ(𝐤^{})i\mathrm{\Gamma }}`$ can be computed as a function of $`p`$ and the results fitted to a semicircular form as $`\rho (E)t^2p(1p)`$, with $`E`$ within the $`H_0`$ band and $`\rho (E)`$ the corresponding density of states. This allows to estimate the lifetime of states $`\tau \frac{1}{\rho (E)t^2p(1p)}`$ and the corresponding mean free paths
$$\lambda _{}=\frac{\tau u_{}^2}{u}\frac{1}{\rho (E)t^2p(1p)\sqrt{2+p^2t^2}},$$
(18)
$$\lambda _{}=\frac{\tau u_{}^2}{u}\frac{p}{\rho (E)(1p)\sqrt{2+p^2t^2}},$$
(19)
for Fermi velocites $`u_{()}`$.
We observe that for small-$`p`$ the obtained mean free path in the parallel (perpendicular) direction is proportional to $`1/p`$ ($`p`$). This implies that the scattering strength of the interlayer disorder increases with $`p`$ for parallel transport but decreases with $`p`$ for perpendicular transport. This intrinsically anisotropic situation seems in contradiction with the scaling theory of localization because $`p0`$ might be thought to be in favor of localization in the perpendicular direction, while transport is the least affected giving extended states in the parallel direction. However, the application of the one-parameter scaling theory to this situation should give a common critical point in all directions so that transport in the perpendicular direction should be extended as well. This is, indeed, numerically confirmed in Chapter IV where we show that the scaling behavior of the conductance in finite cubic systems with interlayer disorder is in agreement with the scaling theory. It turns out that the interlayer disorder is not sufficient to localize the electrons, even in the “difficult” layering direction, without any additional diagonal disorder $`W`$. The statistical behavior of the conductance, however, is very different in the two directions.
## IV numerical calculation of the conductance
### A The cube
The parallel (perpendicular) conductance $`\frac{e^2}{h}g_{()}(L)`$ can be obtained for a $`L\times L\times L`$ cubic system at the Fermi energy $`E`$ directly from the multichannel Landauer–Buttiker formula
$$g_{()}(L)=\text{Tr}(𝐭_{()}^+𝐭_{()}),$$
(20)
where $`𝐭_{()}`$ is the transmission matrix for electronic propagation along the $``$ ($``$) direction computed by transfer matrix techniques. Two perfect semi-infinite bars are attached to two opposite sides of this cube and hard-wall boundary conditions are used for the other sides. The number of independent channels for the incoming and outgoing leads is $`L^2`$ and the transmission matrix can be calculated from the amplitudes of transmitted waves in the outgoing leads by assigning, at a time, a unit incident wave amplitude for one channel in the incoming leads and zero for the rest. We can easily establish the recursion relations for the corresponding wave function coefficients along the parallel direction to calculate $`g_{}`$. In the $``$ direction computing the matrix $`𝐭_{}`$ is not possible since the recursion relations are singular due to the presence of zero hoppings for the missing interplane bonds. To overcome this difficulty we set up the recursion relations along the direction parallel to the planes, perpendicular to the leads, but with the reflection and transmission coefficients in the channels as the unknown variables. We can solve these recursion relations with hard wall boundary conditions perpendicular to the leads. By this method we obtain the perpendicular transmission matrix $`𝐭_{}`$ avoiding the singularities due to missing bonds. This is a convenient tool to consider propagation in the layering $``$ direction by a transfer matrix product only along the easy $``$ direction. In order to suppress fluctuations we have taken averages over up to a $`5000`$ random cubic configurations in each case .
### B The wire
We can also compute the parallel (perpendicular) dimensionless conductance $`g_{()}(M)`$ for a quasi-one-dimensional $`M\times M\times L`$ geometry, via Green function methods . In the parallel direction $`𝐦=x,y`$ the Hamiltonian of the $`M\times M`$ slice is incorporated into the transfer matrix
$$T_x=\left(\begin{array}{cc}V_{x,x+1}^1& 0\\ 0& 1\end{array}\right)\left(\begin{array}{cc}EH_x& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}1& 0\\ 0& V_{x1,x}\end{array}\right),$$
(21)
where the matrix $`V_{x,x+1}`$ has unit elements. For large length $`L`$ in the $`x`$-direction the product
$$T=\underset{x=1}{\overset{L}{}}T_x,$$
(22)
has eigenvalues $`\mathrm{exp}\gamma _iL`$ and Lyapunov exponents $`\gamma _i`$, $`i=1,2,\mathrm{},M^2`$. The smallest positive Lyapunov exponent $`\gamma _1`$ determines the scaling parameter $`\mathrm{\Lambda }_{}`$, via
$$\mathrm{\Lambda }_{}^1=M\gamma _1=\left(\frac{\xi _M}{M}\right)^1,$$
(23)
with $`\xi _M=\frac{1}{\gamma _1}`$ the largest localization length, which is of interest for finite-size scaling studies when the width $`M^2`$ of the slice is varied. We obtain the critical value $`\mathrm{\Lambda }_c`$ at the point where $`\mathrm{\Lambda }_{}`$ becomes independent of $`M`$ (see Table 1). In order to ensure accuracy of about 1% for $`\mathrm{\Lambda }_{}`$ the length of the studied system is more than about $`200000`$. We find that this length near the critical point varies as $`\mathrm{\Lambda }^1`$. In the case of very strong anisotropy we were unable to obtain this accuracy for all $`M`$ and $`W`$.
This method cannot be used in the $``$-direction. The reason is, again, the zero elements of the hopping matrix $`V_{l,l+1}`$ with probability $`1p`$ so that the inverse matrix $`V_{l,l+1}^1`$ which enters (21) becomes singular. In order to avoid this problem we can alternatively use the Green function $`G(E)=(EH)^1`$ by applying the iterative scheme of based on the two equations
$$G_{1,l+1}^{(l+1)}=G_{1,l}^{(l)}V_{l,l+1}G_{l+1,l+1}^{(l+1)}$$
(24)
and
$$G_{l+1,l+1}^{(l+1)}=[EH_{l+1}V_{l+1,l}G_{l,l}^{(l)}V_{l,l+1}]^1.$$
(25)
The Hamiltonian $`H_l`$ represents the $`l`$th $`M\times M`$ layer, $`V_{l,l+1}`$ the hopping between layers $`l`$, $`l+1`$ and $`G_{l_1,l_2}^{(l)}`$ is the Green function of the system with length $`l`$ between layers $`l_1`$ and $`l_2`$. The diagonal matrix $`V_{l,l+1}`$ of the order of $`M^2`$ has zeros with probability $`1p`$ so that $`V_{l,l+1}^1`$ becomes singular. The main advantage of the formulae (24,25) is that they do not contain the inverse of $`V_{l,l+1}`$. The disadvantage is the necessity to invert a matrix of order $`M^2`$ \[Eq. (25)\] at each iteration step. In this case we restrict the number of iterations to $`L40000`$ and the corresponding peprendicular scaling parameter is defined as
$$\mathrm{\Lambda }_{}^1=\frac{M}{L}\mathrm{log}\mathrm{Tr}G_{1L}^L=\left(\frac{\xi _M}{M}\right)^1.$$
(26)
The critical disorder $`W_c`$ and the critical exponent $`\nu `$ are computed in both directions from the numerical data of $`\mathrm{\Lambda }(M,W)`$. The usual linearization near the critical point
$$\mathrm{log}\mathrm{\Lambda }(M,W)=\alpha _M+\beta _M\mathrm{log}W,$$
(27)
and the independence of $`\mathrm{\Lambda }`$ on $`M`$ at the critical point gives $`W_c`$ from the slope of the linear dependence
$$\alpha _M=\mathrm{log}W_c\times \beta _M+\mathrm{const}.$$
(28)
The obtained critical $`\mathrm{\Lambda }_c`$ and $`\mathrm{\Lambda }_c`$ are different from the value $`\mathrm{\Lambda }_c`$ of the corresponding isotropic system (see however Eq. (32) below). The critical exponent $`\nu `$ is determined from the coefficients $`\beta _M`$ via
$$\nu =\frac{\mathrm{log}M}{\mathrm{log}\beta _M}.$$
(29)
We have also repeated the computations for cubic systems and although the finite size effects became more pronounced the obtained critical values differ very little from those obtained for long wires.
In agreement with the $`W=0`$ case it is reasonable to suppose that the form of $`P_c(g)`$ also depends on the direction. For the $``$ direction the conductance is calculated from the formula
$$g(M)=\underset{i=1}{\overset{M^2}{}}\frac{2}{\mathrm{cosh}^2(z_i/2)}$$
(30)
where the $`z_i`$’s are the logarithms of the $`i`$th eigenvalues of the matrix $`𝐭^{}𝐭`$ and in the limit $`L>>M`$ converge to $`2M\gamma _i`$. In the $``$ direction the singular behavior of the matrix $`V_{l,l+1}`$ does not permit to use the formula (30). To calculate $`\mathrm{log}g_{}`$ in this case we use the fact that $`\mathrm{\Lambda }_c<<1`$ for most of the critical points discussed below. Then, the critical conductance from $`z_1=2/\mathrm{\Lambda }_c`$ is very small and can be estimated from the contribution of the first channel via
$$\mathrm{log}g_{}=\mathrm{Tr}G_{1,L}(E+i\eta )G_{L,1}(Ei\eta ),$$
(31)
with the imaginary part of the energy $`\eta `$ given by the ratio of the bandwidth over the mean level spacing .
## V results
### A $`W=0`$
Fig. 2 shows the scaling behavior of $`g_{()}(L)`$ for the $`L^3`$-site cubic system for various bond densities $`p`$ and anisotropic interplane coupling $`t=0.3`$, in the absence of diagonal disorder $`W`$. The parallel conductance is shown to increase ballistically ($`L^2`$) for small $`L`$ and linearly for higher $`L`$. In the large size ($`L\mathrm{}`$) limit the corresponding scaling function $`\beta (L)=d\mathrm{ln}g/d\mathrm{ln}L`$ becomes positive for $`g_{()}`$, which implies extended states in both directions for any $`p`$, in agreement with the scaling theory which predicts a common critical point in any direction. However, the obtained transport behavior is essentially different in the two directions. In Fig. 3 we show the energy dependence of the conductance for $`p=0.5`$ and $`t=0.3`$ where the ratio of the two conductances is close to the estimate $`g_{}/g_{}(t_{}/t_{})^2`$ for $`t_{}=1`$ and $`t_{}=pt=0.15`$ . A key finding from Fig. 2 is a rather smooth $`g_{}(E)`$ while $`g_{}(E)`$ displays violent oscillations as a function of $`E`$. The dips in $`g_{}(E)`$ can be regarded as due to “minigaps” in the perpendicular direction which might have effect similar to a semiconductor, leading to insulating kind of behavior for the out-of-plane conductivity when the Fermi energy is varied.
### B $`W>0`$
The $`W`$-dependence of $`\mathrm{\Lambda }(M,W)`$ for different $`M`$’s and various parameters $`p,t`$ is presented in Figs 4., 5., 6., 7. The corresponding critical points are calculated by the described fitting procedure. The results are listed in Table 1 and the data are very reliable for the $``$ direction. In the $``$ direction they are much harder to analyze due to finite-size effects. For example, in case D (Fig. 7) the obtained $`M`$-dependence of $`\mathrm{\Lambda }_{}(M)`$ is not monotonic and for $`W=5`$ we find $`\mathrm{\Lambda }(M)`$ which decreases with $`M`$ for small system sizes, imitating insulating behavior. However, for larger $`M>12`$ the values of $`\mathrm{\Lambda }(M)`$ begin to increase and the correct scaling is restored. This $`M`$-dependence is caused by a second irrelevant scaling parameter from the relation $`\mathrm{\Lambda }=aM^{1/\nu }+bM^\beta `$ with $`\beta <0`$. The correct estimation of the critical parameters in the $``$ direction requires either numerical data for larger $`M`$ or possibly more sophisticated fits . Although for strong anisotropy we could not obtain accurately the critical parameters in the $``$ direction we check that our results converge to those of the $``$ direction when $`M`$ grows.
In this direction the critical region is very narrow and we could neither calculate the critical exponent, since larger system sizes are required. Nevertheless, the scaling analysis for A and B gives satisfactory results in both directions which confirm the equal $`W_c`$ and $`\nu `$. Moreover, we find the quantity
$$\mathrm{\Lambda }_c=\left[\mathrm{\Lambda }_c^2\times \mathrm{\Lambda }_c\right]^{1/3}$$
(32)
which gives the critical value of the corresponding isotropic model .
Fig. 8(a), (b) presents the probability distribution of the critical conductance $`P_c(g_{})`$ in the parallel direction. The distribution is shown to be size-independent but depends on the various critical points. For C and D the obtained $`P_c(g_{})`$ and $`\mathrm{\Lambda }_c`$ are the same. This indicates that $`P_c(g_{})`$ is determined only by $`\mathrm{\Lambda }_c`$ . The numerical data for the mean and variance of the conductance are also unique functions of $`\mathrm{\Lambda }_c`$ (see Table 1) supporting this conjecture. In the limit $`t0`$ the critical distribution in the parallel direction is expected to converge to a Gaussian. However, the spectrum of the obtained higher Lyapunov exponents shows square-root behavior similar to isotropic $`3D`$ disordered systems . To display the dramatic differences in parallel and perpendicular transport, we also present the critical distribution of $`\mathrm{log}g_{}`$, calculated for the perpendicular direction for the critical case C. This distribution has all the features of the localized regime since it is log-normal with
$$\mathrm{var}\mathrm{log}g_{}\mathrm{log}g_{}.$$
(33)
The important difference with the insulating regime is the fact that $`P_c(\mathrm{log}g_{})`$ remains system-size invariant. The relation of Eq. (33) is also valid for the critical points B and D where in the $``$ direction $`g_{}<<1`$.
## VI discussion
The random topological multilayered structure studied may be regarded as a first step towards an explanation, via non-interacting electrons, for properties of strongly anisotropic materials. For example, in the case of $`W=0`$ the $`2D`$ layers are perfect and the disorder represented by $`p`$ can be due to impurities or oxygen vacancies in the insulating layer among the $`2D`$ planes of the cuprates. The electrons propagating in the perpendicular direction of this system shall encounter anisotropy in the disorder due to the random interplane links in addition to the value of $`t`$ which can be different to that of the parallel direction. However, the distribution of the critical conductance $`P_c(g)`$ depends both on the choice of the parameters and the direction where the electron moves.
The considered anisotropic structure exhibits rather strange transport properties on a given scale, expressed in the dramatic differences of the critical conductance in the parallel and perpendicular directions. In the perpendicular direction the conductance distribution slightly bellow the critical point is log-normal resembling the statistical properties typical of an insulator. A similar statistical “anomaly” has been described in . However, a strong difference to a “true” insulating regime exists since the conductance still grows with the size. In the large size limit the corresponding distribution reaches a Gaussian. An analogous discussion holds for the parallel direction slightly above the critical point. The main criterion for the specification of the critical regime is the size dependence of the conductance $`g`$ (or $`\mathrm{log}g`$).
The proposed model for $`W=0`$ may have some relation to the strongly anisotropic transport proprties of high–$`T_c`$ cuprates. In these materials, as temperature increases, the inelastic scattering due to phonons, spin waves or other excitations within the $`CuO`$ planes, can cause a decrease of the inelastic scattering length $`l_{in}`$. If the temperature is so high that $`l_{in}`$ becomes smaller than mean free path the transverse conductivity is metallic. Experiments for Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub> and underdoped La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>, YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6+x</sub> give out–of–plane resistivity which has a semiconductor–like temperature dependence at low temperatures (high at small–$`T`$ with a rapid decrease by increasing $`T`$) and a linear–in–$`T`$ behavior at high temperatures. The characteristic crossover temperature between the two regimes $`T^{}`$ decreases by increasing the doping in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> and YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6+x</sub> . We notice that if we relate the bond density $`p`$ with the doping density of the high-$`T_c`$ materials the obtained $`p`$-dependence of the perpendicular mean free path can be used to explain qualitatively the reported behavior. It must be pointed out that the relation between the cuprate doping density and the bond density $`p`$ is natural, since an increase in the number of the doping impurities or the oxygen atoms in the layer between two CuO<sub>2</sub> planes increases the number of hopping paths between the two planes. As $`T^{}`$ decreases further (below $`T_c`$) the out–of–plane normal–state resistivity also becomes metallic, which has been observed in high–quality single crystals of YBA<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> and other high–$`T_c`$ cuprates corresponding to the absence of disorder ($`p1`$ with $`W=0`$ in the proposed model) with almost infinite perpendicular mean free path. This effect occurs only in the perpendicular direction since the parallel mean free path is always much longer and inelasting scattering becomes dominant.
In summary, we have introduced a simple layered lattice model with anisotropic disorder described by the interplane bond density $`p`$, in addition to the usual anisotropic band structure expressed via the interplane hopping $`t`$. In the absence of diagonal disorder we show extended states in both directions but the obtained mean free path and the conductance in the $``$ direction is much smaller than in the $``$ direction. Moreover, $`g_{}`$ fluctuates strongly as a function of energy, which leads to an insulator-like temperature dependence of the conductivity in the $``$ direction. In the presence of additional diagonal disorder of strength $`W`$ we have shown that the critical disorder and the critical exponent $`\nu `$ do not depend on the transport direction. The obtained data for the localization exponent $`\nu `$ agree with recent accurate estimates for the isotropic model and confirm the universality at the metal-insulator transition. The obtained critical conductance distribution $`P_c(g)`$ although independent on the system size depends strongly on the parameters and the direction of transport.
Acknowledgments
This work was supported in part by a TMR network. SNE and SJX thank a Chino-Greek grant and PM the Slovak Grant Agency and NATO. We also like to thank Professors Xing and Economou for many useful discussions.
FIG. 1. A picture of the multilayered structure which consists of $`2D`$ square lattices $`l`$ (layers) connected by perpendicular bonds of strength $`t`$ placed at random positions with probability $`p`$.
FIG. 2. The $`g_{}`$ as a function of the linear system size $`L`$ for a cubic $`L\times L\times L`$ system of parallel planes with $`W=0`$ and randomly placed interplane bonds of density $`p`$ with strength $`t`$=0.3. In the inset $`g_{}`$ for the same system exhibits similar behavior but much smaller values.
FIG. 3. (a) The energy-dependent $`g_{}`$ for a cubic layered system with $`L`$=10,15 and $`W=0`$, $`t`$=0.3, $`p`$=0.5. (b) The $`g_{}`$ is much smaller and exhibits violent oscillations as a function of the energy $`E`$.
FIG. 4. A: The behavior of the scaled localization length $`\xi _M/M`$ for the parallel and perpendicular direction in the $`M\times M\times L`$ system with $`p=0.6`$ and $`t=0.1`$ as a function of $`W`$. The critical point is located at $`W_c14.47`$ in the $``$ direction and $`W_c14.30`$ in the $``$ direction.
FIG. 5. B: The behavior of the scaled localization length $`\xi _M/M`$ for the parallel and perpendicular direction in the $`M\times M\times L`$ system with $`p=0.6`$ and $`t=0.3`$ as a function of $`W`$. The critical point is located at $`W_c10.48`$ in the $``$ and $`W_c10.20`$ in the $``$ direction.
FIG. 6. C: The behavior of the scaled localization length $`\xi _M/M`$ for the parallel and perpendicular direction in the $`M\times M\times L`$ system for $`p=0.1`$ and $`t=0.3`$ as a function of $`W`$. The critical point is displayed at $`W_c7.93`$ in the $``$ and $`W_c7.18`$ in the $``$ direction.
FIG. 7. D: The behavior of the scaled localization length $`\xi _M/M`$ for the $``$ and $``$ direction in the $`M\times M\times L`$ system with $`p=0.6`$ and $`t=0.1`$ as a function of $`W`$. The critical point is located in $`W_c8.05`$ in the $``$ and $`W_c6.80`$ in the $``$ direction. It is seen that the data for smaller $`M=6,8`$ fail to cross at the same point indicating “insulating” behavior.
FIG. 8. $`(𝐚)`$. The critical distribution of the parallel conductance in the case A with $`<g_{}>0.59`$, var$`(g_{})=0.16`$. $`(𝐛)`$. The same as in (a) for the critical points C (full symbols) and D(open symbols) with $`<g_{}>2.6`$, var $`(g_{})=0.7`$. The C,D have the same $`\mathrm{\Lambda }_c`$ and the same critical distribution (see Table 1). $`(𝐜)`$. The critical distribution of $`P(\mathrm{log}g_{})`$ in the perpendicular direction for case C is also shown for comparison.
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# Neutrino Propagation In Color Superconducting Quark Matter
## I Introduction
In this article we study heat diffusion via neutrinos in dense, color superconducting quark matter. Recent theoretical works suggest that quarks form Cooper pairs in medium, a natural consequence of attractive interactions destabilizing the Fermi surface. This would likely affect the early evolution of neutron stars born through a Type II supernova explosion, where the central role of neutrino diffusion through a strongly-interacting medium has long been postulated on theoretical grounds .
Type II (core collapse) supernovae are triggered by the implosion of the inner core of a massive star (M<sub>star</sub> $`820`$ M), when the core mass is on the order of the Chandrashekar mass (M<sub>core</sub> $`1.4`$ M). During the implosion nearly all ($`99`$%) of the enormous gravitational binding energy ($`10^{53}`$ ergs) gained is stored as internal energy of the newly born, proto-neutron star (PNS). The subsequent evolution of the proto-neutron star is driven by neutrino diffusion. Temporal and spectral characteristics of the neutrino emission depend on the rate at which neutrinos diffuse through the imploded PNS which, at this early stage, is composed of hot ($`T2030`$ MeV) and dense ($`n_B23n_0`$ where $`n_0=0.16`$ fm<sup>-3</sup>) strongly-interacting matter.
Neutrino emission during this phase is a directly observable feature of a galactic supernova explosion. The few neutrinos detected from SN 1987A indicate that neutrinos of mean energy $`E_\nu 20`$ MeV are emitted on a diffusion time scale of about $`1020`$ seconds. It is reasonable to expect that neutrino mean free path in the inner, denser regions of the star will strongly influence the temporal characteristics of a supernova neutrino signal. Supernova neutrinos can therefore reveal properties of matter at high baryon density, at temperatures substantially lower than those expected in relativistic heavy ion collisions.
Although the idea of quark pairing in dense matter is not a new one , it has recently seen renewed interest in the context of the phase diagram of QCD . Model calculations, mostly based on four-quark effective interactions, predict the restoration of spontaneously broken chiral symmetry through the onset of color superconductivity at low temperatures . For much higher densities color superconductivity is manifest through perturbative gluon exchange , which can be calculated systematically suggests that the phenomenon is robust. For densities and temperatures relevant to neutron stars, quark matter is therefore expected to be superconducting. Models generally predict an energy gap of $`\mathrm{\Delta }100`$ MeV for a typical quark chemical potential of $`\mu _q400`$ MeV. As in BCS theory, the gap will weaken for $`T>0`$ and at some $`T=T_c`$ there is a (second-order) transition to a “standard” quark-gluon plasma. Thus, when cooling from temperatures greater than critical, the formation of a such a gap in the fermionic excitation spectrum in quark matter at high density will influence various neutron star observables – if neutron stars contain quark matter in their cores at early time<sup>*</sup><sup>*</sup>* This remains unclear. In particular, the existence of a finite electron neutrino chemical potential at very early time has shown to inhibit the appearance of quark matter. See Ref. for a review.. Two examples recently investigated are the effects on magnetic fields and on the thermal evolution of neutron stars at late time , when interior temperatures evolve from $`T1`$ MeV to $`T\stackrel{<}{}\text{ }1`$ KeV. In this work we shall consider the cooling of quark matter at earlier times, when temperatures pass through the critical $`T_c\mathrm{\Delta }`$. Assuming a simplified scenario, we investigate the thermal evolution of the inner core of a proto-neutron star as it cools via neutrino diffusion during its first few seconds.
Our main finding is that the neutrino mean free path in the superconducting phase has a strong dependence on temperature. An energy gap tends to increase the neutrino mean free paths exponentially when $`T\mathrm{\Delta }`$. However in the intermediate regime, when $`T\mathrm{\Delta }`$, the temperature dependence is not exponential. From this we predict a uniquely uneven cooling process for simplified PNS matter, marked by a slowdown of the cooling rate when the system undergoes a second-order phase transition. This is a consequence of the specific heat being peaked at $`T_c`$, a standard characteristic of a phase transition, rather than a large change in the neutrino mean free path. Since the energy gap vanishes at this point, the mean free path is not significantly modified in the neighborhood of $`T_c`$.
In Section II the neutrino mean free path, always denoted $`\lambda `$ in this work, is computed in a background of superconducting quarks. Taking a general, BCS-type model for the energy gap, we relate $`\lambda `$ to in-medium quark polarizations via the differential and total neutrino-quark cross sections. Explicit formulae are derived for the imaginary parts of the vector and axial-vector polarizations of a relativistic Fermi system with an energy gap. After assuming BCS-type mean field behavior of the gap and the specific heat, in Section III we compute a characteristic time for heat diffusion via neutrino emission from a simple model of superconducting quark matter. We then consider some astrophysical consequences and outline the model’s applicability to proto-neutron stars. Section IV contains our conclusions.
## II Neutrino-quark Scattering in a Color Superconductor
The primary process by which heat escapes a proto-neutron star is neutrino diffusion, and so a significant consequence of color superconductivity in this context will be modified neutrino propagation. While noting that the neutrino production rate will also differ from that of normal matter, in the diffusive regime one can see that the dominant critical behavior will be a change in the inelastic quark-neutrino cross section, since here neutrino production rates decouple from the transport equation and depend only on the neutrino mean free path. In this section we calculate the differential and total cross sections and then compute the neutrino mean free path in two-flavor quark matter. The magnitude of the superconducting gap, $`\mathrm{\Delta }`$, is taken as arbitrary within a range of values found in recent literature. Closely following BCS theory, we assume the gap to be a constant, and calculate the response functions and neutrino cross sections in the weak coupling approximation.
The neutral current coupling between neutrinos and quarks, a four-fermion effective interaction for energies $`E_\nu M_Z`$, is written as
$$_W=\frac{G_F}{\sqrt{2}}\overline{\nu }\gamma _\mu (1\gamma _5)\nu \overline{q}\gamma ^\mu (C_VC_A\gamma _5)q,$$
(1)
where $`G_F=1.166\times 10^5`$ GeV<sup>-2</sup> is the Fermi weak coupling constant and $`C_V`$ and $`C_A`$ are the flavor-specific vector and axial vector coupling constants, respectively. The differential neutrino scattering cross section per unit volume in an infinite and homogeneous system of fermions as calculated in linear response theory is
$`{\displaystyle \frac{1}{V}}{\displaystyle \frac{d^3\sigma }{d^2\mathrm{\Omega }_3dE_3}}`$ $`=`$ $`{\displaystyle \frac{G_F^2}{32\pi ^2}}{\displaystyle \frac{E_3}{E_1}}{\displaystyle \frac{\left[1f_\nu (E_3)\right]}{\left[1\mathrm{exp}\left(q_0/T\right)\right]}}\mathrm{Im}(L^{\alpha \beta }\mathrm{\Pi }_{\alpha \beta }),`$ (2)
where $`E_1`$ ($`E_3`$) is the incoming (outgoing) neutrino energy. The factor $`[1\mathrm{exp}(q_0/T)]^1`$ maintains detailed balance and the final state blocking of the outgoing neutrino is enforced by the Pauli blocking factor, $`[1f_\nu (E_3)]`$. The neutrino tensor $`L_{\alpha \beta }`$ is given by
$$L^{\alpha \beta }=8[2k^\alpha k^\beta +(kq)g^{\alpha \beta }(k^\alpha q^\beta +q^\alpha k^\beta )iϵ^{\alpha \beta \mu \nu }k_\mu q_\nu ],$$
(3)
where the incoming four-momentum is $`k^\alpha `$ and the momentum transferred to the medium is $`q^\alpha `$. The minus (plus) sign on the final term applies to neutrino (anti-neutrino) scattering.
The medium is characterized by the quark polarization tensor $`\mathrm{\Pi }_{\alpha \beta }`$. In the case of free quarks, each flavor contributes a term of the form
$$\mathrm{\Pi }_{\alpha \beta }(q)=i\mathrm{Tr}_c\frac{d^4p}{(2\pi )^4}\mathrm{Tr}[S_0(p)\mathrm{\Gamma }_\alpha S_0(p+q)\mathrm{\Gamma }_\beta ],$$
(4)
where $`S_0(p)`$ is the free quark propagator at finite chemical potential and temperature. The outer trace is over color and simplifies to a $`N_c=3`$ degeneracy. The inner trace is over spin, and the $`\mathrm{\Gamma }_\alpha `$ are the neutrino-quark vertex functions which determine the spin channel. Specifically, the vector polarization is computed by choosing $`(\mathrm{\Gamma }_\alpha ,\mathrm{\Gamma }_\beta )=(\gamma _\alpha ,\gamma _\beta )`$. The axial and mixed vector-axial polarizations are similarly obtained from $`(\mathrm{\Gamma }_\alpha ,\mathrm{\Gamma }_\beta )=(\gamma _\alpha \gamma _5,\gamma _\beta \gamma _5)`$ and $`(\mathrm{\Gamma }_\alpha ,\mathrm{\Gamma }_\beta )=(\gamma _\alpha ,\gamma _\beta \gamma _5)`$, respectively.
The free quark propagators in Eq. (4) are naturally modified in a superconducting medium. As first pointed out by Bardeen, Cooper, and Schrieffer several decades ago, the quasi-particle dispersion relation is modified due to the presence of a gap in the excitation spectrum. In calculating these effects, we will consider the simplified case of QCD with two quark flavors which obey SU(2)$`{}_{L}{}^{}\times `$ SU(2)<sub>R</sub> flavor symmetry, given that the light $`u`$ and $`d`$ quarks dominate low-energy phenomena. Furthermore we will assume that, through some unspecified effective interactions, quarks pair in a manner analogous to the BCS mechanism . The relevant consequences of this are the restoration of chiral symmetry (hence all quarks are approximately massless) and the existence of an energy gap at zero temperature, $`\mathrm{\Delta }_0`$, with approximate temperature dependence,
$$\mathrm{\Delta }(T)=\mathrm{\Delta }_0\sqrt{1\left(\frac{T}{T_c}\right)^2}.$$
(5)
The critical temperature $`T_c0.57\mathrm{\Delta }_0`$ is likewise taken from BCS theory; this relation has been shown to hold for perturbative QCD and is thus a reasonable assumption for non-perturbative physics.
Breaking the SU<sub>c</sub>(3) color group leads to complications not found in electrodynamics. In QCD the superconducting gap is equivalent to a diquark condensate, which can at most involve two of the three fundamental quark colors. The condensate must therefore be colored. Since the scalar diquark (in the $`\overline{\mathrm{𝟑}}`$ color representation) appears to always be the most attractive channel, we consider the anomalous (or Gorkov) propagator
$`F(p)_{abfg}`$ $`=`$ $`q_{fa}^T(p)C\gamma _5q_{gb}(p)`$ (6)
$`=`$ $`iϵ_{ab3}ϵ_{fg}\mathrm{\Delta }\left({\displaystyle \frac{\mathrm{\Lambda }^+(p)}{p_o^2\xi _p^2}}+{\displaystyle \frac{\mathrm{\Lambda }^{}(p)}{p_o^2\overline{\xi }_p^2}}\right)\gamma _5C.`$ (7)
Here, $`a,b`$ are color indices, $`f,g`$ are flavor indices, $`ϵ_{abc}`$ is the usual anti-symmetric tensor and we have conventionally chosen 3 to be the condensate color. This propagator is also antisymmetric in flavor and spin, with $`C=i\gamma _0\gamma _2`$ being the charge conjugation operator.
The color bias of the condensate forces a splitting of the normal quark propagator into colors transverse and parallel to the diquark. Quarks of color 3, parallel to the condensate in color space, will be unaffected and propagate freely, with
$$S_0(p)_{af}^{bg}=i\delta _a^b\delta _f^g\left(\frac{\mathrm{\Lambda }^+(p)}{p_o^2E_p^2}+\frac{\mathrm{\Lambda }^{}(p)}{p_o^2\overline{E}_p^2}\right)(p_\mu \gamma ^\mu \mu \gamma _0).$$
(8)
This is written in terms of the particle and anti-particle projection operators $`\mathrm{\Lambda }^+(p)`$ and $`\mathrm{\Lambda }^{}(p)`$ respectively, where $`\mathrm{\Lambda }^\pm (p)=(1\pm \gamma _0\stackrel{}{\gamma }\widehat{p})/2`$. The excitation energies are simply $`E_p=|\stackrel{}{p}|\mu `$ for quarks and $`E_p=|\stackrel{}{p}|+\mu `$ for anti-quarks.
On the other hand, transverse quark colors 1 and 2 participate in the diquark and thus their quasi-particle propagators are given as
$$S(p)_{af}^{bg}=i\delta _a^b\delta _f^g\left(\frac{\mathrm{\Lambda }^+(p)}{p_o^2\xi _p^2}+\frac{\mathrm{\Lambda }^{}(p)}{p_o^2\overline{\xi }_p^2}\right)(p_\mu \gamma ^\mu \mu \gamma _0).$$
(9)
The quasi-particle energy is $`\xi _p=\sqrt{(|\stackrel{}{p}|\mu )^2+\mathrm{\Delta }^2}`$, and for the anti-particle $`\overline{\xi }_p=\sqrt{(|\stackrel{}{p}|+\mu )^2+\mathrm{\Delta }^2}`$.
The appearance of an anomalous propagator in the superconducting phase indicates that the polarization tensor gets contributions from both the normal quasi-particle propagators (9) and anomalous propagator (7). Thus, to order $`G_F^2`$, Eq. (4) is replaced with the two contributions corresponding to the diagrams shown in Fig. 1, and written
$$\mathrm{\Pi }_{\alpha \beta }(q)=i\frac{d^4p}{(2\pi )^4}\left\{\mathrm{Tr}[S_0(p)\mathrm{\Gamma }_\alpha S_0(p+q)\mathrm{\Gamma }_\beta ]+2\mathrm{T}\mathrm{r}[S(p)\mathrm{\Gamma }_\alpha S(p+q)\mathrm{\Gamma }_\beta ]+2\mathrm{T}\mathrm{r}[F(p)\mathrm{\Gamma }_\alpha \overline{F}(p+q)\mathrm{\Gamma }_\beta ]\right\}.$$
(10)
The remaining trace is over spin, as the color trace has been performed. Fig. 1(a) corresponds to the first two terms, which have been decomposed into one term with ungapped propagators (8) and the other with gapped quasi-particle propagators (9). Fig. 1(b) represents the third, anomalous term.
For neutrino scattering we must consider vector, axial, and mixed vector-axial channels, all summed over flavors. The full polarization, to be used in evaluating Eq. (2), may be written
$$\mathrm{\Pi }_{\alpha \beta }=\underset{f}{}\left[(C_V^f)^2\mathrm{\Pi }_{\alpha \beta }^V+(C_A^f)^2\mathrm{\Pi }_{\alpha \beta }^A2C_V^fC_A^f\mathrm{\Pi }_{\alpha \beta }^{VA}\right].$$
(11)
The coupling constants for up quarks are $`C_V^u=\frac{1}{2}\frac{4}{3}\mathrm{sin}^2\theta _W`$ and $`C_A^u=\frac{1}{2}`$ , and for down quarks, $`C_V^d=\frac{1}{2}+\frac{2}{3}\mathrm{sin}^2\theta _W`$ and $`C_A^d=\frac{1}{2}`$, where $`\mathrm{sin}^2\theta _W0.23`$ is the Weinberg angle.
If we specify a frame in which the transfer momentum is $`q_\mu =(q_0,q,0,0)`$ we can separate longitudinal components as
$$\mathrm{\Pi }_L^V=\frac{q_\mu ^2}{q^2}\mathrm{\Pi }_{00}^V=\frac{q_\mu ^2}{q_0^2}\mathrm{\Pi }_{11}^V=\frac{q_\mu ^2}{q_0q}\mathrm{\Pi }_{10}^V,$$
(12)
and the transverse,
$$\mathrm{\Pi }_T^V=\mathrm{\Pi }_{22}^V=\mathrm{\Pi }_{33}^V.$$
(13)
Identical definitions apply to the axial polarizations $`\mathrm{\Pi }_L^A`$ and $`\mathrm{\Pi }_T^A`$. The non-zero mixed correlation function is written
$$\mathrm{\Pi }_{\alpha \beta }^{VA}=iϵ_{\alpha \beta \mu 0}q^\mu \mathrm{\Pi }^{VA}.$$
(14)
Detailed calculations of these polarization functions are given in the Appendix.
As the gap increases, the superconducting quasi-particles naturally become the dominant excitations of the background, a property clearly visible in the neutrino response functions. The left panel of Fig. 2 shows the longitudinal response in the vector channel. The free quark case, shown as a solid line labeled $`\mathrm{\Delta }=0`$, is the standard result describing Pauli-blocking and kinematics of massless single particle excitations . Here, energy-momentum conservation restricts the response to the spacelike region ($`q_0<q`$). Superconductivity modifies this result, as the quasi-quark excitations become suppressed due to the pairing correlations at the Fermi surface. At the same time, the response is enhanced when $`q_02\mathrm{\Delta }`$, signifying the excitation of Cooper pairs. In particular, the results with greater $`\mathrm{\Delta }`$ clearly show the threshold for these excitations at energy transfer $`q_0=2\mathrm{\Delta }`$. Results for $`\mathrm{\Delta }=10,\mathrm{\hspace{0.17em}30}`$ and $`50`$ MeV show the gradual reallocation of response strength from small $`q_0`$ to the region $`q_02\mathrm{\Delta }`$. Since scattering probes only the spacelike region, the $`\mathrm{\Pi }_{00}^V`$ contribution to the cross section is generically suppressed in the superconducting phase. Contributions at $`q_0q`$ will contribute to the neutrino production rate, rather than scattering cross section, when the temperature is near $`\mathrm{\Delta }`$.
Analytic results may be obtained for small transfer energies. When $`q_0`$ is much smaller than all other energy scales, the vector longitudinal response (see Eq. (27) of the Appendix) reduces to
$`\underset{q_00}{lim}\mathrm{Im}\mathrm{\Pi }_L^V(\mathrm{\Delta })`$ $`=`$ $`{\displaystyle \frac{2}{1+e^{\beta \mathrm{\Delta }}}}\mathrm{\Pi }_L^V(\mathrm{\Delta }=0)`$ (15)
$`=`$ $`{\displaystyle \frac{\mu ^2q_0}{2\pi q}}{\displaystyle \frac{1}{1+e^{\beta \mathrm{\Delta }}}}.`$ (16)
From this the weakening of the low energy, vector-longitudinal response can be calculated for a given gap $`\mathrm{\Delta }`$.
The axial-longitudinal response, which physically corresponds to the excitation of spin waves, is shown in the right panel of Fig. 2. As with the vector channel a threshold of $`2\mathrm{\Delta }`$ is apparent but, unlike the previous case, the response as $`q_00`$ is enhanced. This is manifest in the the limit $`q_0T\mathrm{\Delta }`$, where one finds
$$\mathrm{Im}\mathrm{\Pi }_L^A(\mathrm{\Delta })\frac{\mathrm{\Delta }}{2T}\mathrm{sech}^2\left(\frac{\mathrm{\Delta }}{2T}\right)\mathrm{ln}\left(\frac{\mathrm{\Delta }}{q_0}\right)\mathrm{\Pi }_L^A(\mathrm{\Delta }=0),$$
(17)
where in this limit $`\mathrm{\Pi }_L^A(\mathrm{\Delta }=0)=\mu ^2q_0^2/4\pi q`$. This logarithmic enhancement will lead to an integrable peak in the differential cross section at $`q_0=0`$, as will be described below.
The transverse response functions, $`\mathrm{\Pi }_T^V`$ and $`\mathrm{\Pi }_T^A`$ as defined in Eq. (13), exhibit behavior similar to the longitudinal channels. The primary distinction is that superconductive coherence at low $`q_0`$ enhances the vector channel and suppresses the axial. In the former case the interference effects are constructive and in the latter case they are destructive , which is a reversal of the longitudinal results. This difference has been theoretically understood and experimentally observed in electric superconductors, where the absorption of acoustic waves (vector-longitudinal) is suppressed while the infra-red absorption (vector-transverse) is enhanced for small energy transfer There is naturally a vast literature concerning coherence effects in superconducting metals, and we find it both gratifying and reassuring to obtain similar results in this less terrestrial context..
In addition to the vector and axial polarizations, the mixed vector-axial polarization (14) also contributes to the cross section. While this contribution is always much smaller – by at least one order of magnitude – we note that it is enhanced in the superconducting phase, leading to an amplified difference between neutrino and anti-neutrino cross sections in PNS matter.
Once the polarization tensors have been computed it is straightforward to obtain the differential cross section, Eq. (2). Results for neutrinos of energy $`E_\nu =50`$ MeV, ambient matter conditions of $`\mu _q=400`$ MeV and $`T=30`$ MeV, and gaps of varied size are plotted as a function of transfer energy in Fig. 3.
A striking feature of these results is the singular behavior of the differential cross section near $`q_0=0`$ in superconducting matter. This is the (integrable) logarithmic divergence of $`\mathrm{\Pi }_L^A`$ and $`\mathrm{\Pi }_T^V`$ (see (Eq. 17)). The threshold behavior seen in differential cross section, for the case $`\mathrm{\Delta }=10`$ and $`30`$ MeV, at $`q_0=2\mathrm{\Delta }`$ and $`q_0=2\mathrm{\Delta }`$ correspond to excitations of the Cooper pairs.
The total cross section (per unit volume) is obtained by integrating over all neutrino energy transfers and angles. From this the mean free path is determined, since
$$\lambda =\left(\frac{\sigma }{V}\right)^1.$$
(18)
As suggested by the differential cross section in Fig. 3, the total cross section is reduced in the presence of a gap $`\mathrm{\Delta }`$. The logarithmic peak at $`q_0=0`$ has a minimal effect after integration, when $`\mathrm{\Delta }T`$ and when the neutrino energy $`E_\nu \pi T`$ since neutrino probe a significant part of the response outside of this $`q_0=0`$ region.
Results for the neutrino mean free path, $`\lambda `$, are shown in Fig. 4 as a function of incoming neutrino energy $`E_\nu `$ (ambient matter conditions of $`\mu _q=400`$ MeV and $`T=30`$ MeV have again been used). The same energy dependence has been computed previously in the case of free relativistic and degenerate fermionic matter ; it decreases as $`1/E_\nu ^2`$ for $`E_\nu T`$, and $`1/E_\nu `$ at $`E_\nu T`$. The results indicate that this energy dependence is not modified by the presence of a gap when $`\mathrm{\Delta }T`$. Thus the primary effect of the superconducting phase is a much larger mean free path. This is consistent with the suppression found in the vector-longitudinal response function, which dominates the sum polarization sum (11), at $`q_0<q`$.
## III The Cooling of an Idealized Quark Star
The immediate application of the last section’s results is neutrino emission from a proto-neutron star containing quark matter. While we will fall far short of a complete description of the role of color superconductivity in such a complicated environment, in this section we will outline the key ingredients for estimating the possible physical observables. The first subsection describes a simple and general model for the cooling of quark matter through temperatures relevant to proto-neutron stars, and the second addresses the applicability of this scenario to more realistic situations.
### A The Cooling of Superconducting Quark Matter
Having determined the neutrino mean free path in a color superconductor, we now consider the cooling of a macroscopic sphere of quark matter as it becomes superconducting. As stated previously, this toy model is motivated by the possibility that the core of neutron star might contain such matter. Following our preceding calculations, we will consider the relatively simple case of two massless flavors with identical chemical potentials. Furthermore, we will disregard the quarks parallel in color to the condensate; i.e. we consider a background comprised exclusively of quasi-quarks.
The cooling of a spherical system of quark matter from $`TT_c50`$ MeV is driven by neutrino diffusion, for the neutrino mean free path is much smaller than the dimensions of system of astrophysical size, and yet several orders of magnitude larger than the mean free path of the quarks. The diffusion equation for energy transport by neutrinos in a spherical geometry is
$`C_V{\displaystyle \frac{dT}{dt}}={\displaystyle \frac{1}{r^2}}{\displaystyle \frac{L_\nu }{r}},`$ (19)
where $`C_V`$ is the specific heat per unit volume of quark matter, $`T`$ is the temperature, and $`r`$ is the radius. The neutrino energy luminosity for each neutrino type, $`L_\nu `$, depends on the neutrino mean free path and the spatial gradients in temperature and is approximated by
$$L_\nu 6𝑑E_\nu \frac{c}{6\pi ^2}E_\nu ^3r^2\lambda (E_\nu )\frac{f(E_\nu )}{r},$$
(20)
where $`c`$ denotes the speed of light in vacuum. In our analysis we assume that neutrino interactions are dominated by the neutral current scattering common to all neutrino types. Consequently, we take the same neutrino and anti-neutrino mean free path for every neutrino flavor, giving rise to the factor of six in Eq. (20). The equilibrium Fermi distribution, $`f(E_\nu )`$, and the (scattering) mean free path, $`\lambda (E_\nu )`$, are integrated over all neutrino energies, $`E_\nu `$.
The solution to the diffusion equation will depend on the size of the system and its initial temperature gradients. However, we are merely interested in a qualitative description of cooling through a second-order phase transition to superconducting matter. The temporal behavior is characterized by a time scale $`\tau _c`$, which is proportional to the inverse cooling rate and can hence be deduced from Eq. (19). The characteristic time
$$\tau _c(T)=C_V(T)\frac{R^2}{c\lambda (T)},$$
(21)
is a strong function of the function of the ambient matter temperature since it depends on the matter specific heat and the neutrino mean free path. This applies to a system characterized by the radial length $`R`$ and the energy-weighted average of the mean free path, $`\lambda (T)`$. Following our general treatment of the superconducting gap, we assume that the temperature dependence of the specific heat is described by BCS theory. We will then use the results obtained in the previous section to calculate $`\lambda (T)`$. Furthermore, since neutrinos are in thermal equilibrium for the temperatures of interest, we may assume
$$\lambda (T)\lambda (E_\nu =\pi T).$$
(22)
The quantity on the left is dependent on the gap $`\mathrm{\Delta }`$, a dependence computed in the previous section and plotted in the right panel of Fig. 4. The results indicate that for small $`\mathrm{\Delta }/T`$ the neutrino mean free path is not strongly modified, but as $`\mathrm{\Delta }/T`$ increases so too does $`\lambda `$, non-linearly at first and then exponentially for $`\mathrm{\Delta }/T\stackrel{>}{}\text{ }5`$. We note that the diffusion approximation is only valid when $`\lambda R`$ and will thus fail for very low temperatures, when $`\lambda R`$.
The ratio $`\tau _c^\mathrm{\Delta }(T)/\tau _c(T)`$, a measure of the extent to which the cooling rate is changed by a gap, is shown by the solid line in Fig. 5. The ratio $`\lambda /\lambda ^\mathrm{\Delta }`$, plotted with the short-dashed curve, measures the decrease in neutrino interaction rates in the superconducting background. The temperature dependences we have taken from BCS theory, that of the specific heat (dashed curve) and the magnitude of the gap itself (dot-dashed curve), are shown for reference in Fig. 5.
The results shown in Fig. 5 are readily interpreted. The cooling rate around $`T_c`$ is influenced mainly by the peak in the specific heat associated with the second order phase transition, since the neutrino mean free path is not strongly affected when $`\mathrm{\Delta }T`$. Subsequently, as the matter cools, both $`C_V`$ and $`\lambda ^1`$ decrease in a non-linear fashion for $`\mathrm{\Delta }T`$. Upon further cooling, when $`\mathrm{\Delta }T`$, both $`C_V`$ and $`\lambda ^1`$ decrease exponentially. Both of these effects accelerate the cooling process.
We conclude that if it were possible to measure the neutrino luminosity from the hypothetical object described here, a unique temporal profile would be observed. This suggests that if a second order, superconducting transition were to occur inside a proto-neutron star it could be identified by the temporal characteristics of the late time supernova neutrino signal. Specifically, there would be a brief interval during which the cooling would slow around $`TT_c`$, signified by a period of reduced neutrino detection.
### B Relation to Proto-Neutron Stars
The temporal pattern of neutrino emission deduced here would be observable evidence of the onset of color superconductivity in dense matter. However, ours is a simple model, and we must temper our speculations with considerations of more realistic systems.
Neutrinos emitted from the core of a proto-neutron star, be it hadronic or quark matter in any phase, must pass through a large amount of matter in the outer shell. The neutron-rich material is characterized in the outer shell is opaque to neutrinos and will thus even out any sharp temporal features associated with the interior emission. This is the first and foremost of caveats since this will directly impact the possibility of observing a characteristic in the neutrino signal associated with the phase transition.
As discussed in Section II, the scalar diquark condensate in a two-flavor color superconductor is necessarily colored. Therefore quarks of one color, taken conventionally as color 3, will be color-orthogonal to the scalar condensate and can remain ungapped. We have not included these color-3 quarks in our analysis of neutrino scattering since their fate is uncertain ; they could form color-symmetric Cooper pairs, or perhaps a chiral condensate. If we take the simplest scenario and assume that they remain free, their presence in the medium will further dilute any direct effects of the gapped quasi-quarks. Specifically, the physical quantities in plotted in Fig. 5 will not vanish when $`T0`$, instead being reduced to one third of their $`\mathrm{\Delta }=0`$ values. Likewise, the maxima at $`T_c`$ will be reduced relative to a color-neutral background.
The second-order phase transition is taken directly from BCS theory, which we assume rather than derive. While this is the mean-field result for two massless flavors with equivalent chemical potentials, neutron star matter is constrained in two ways. First of all, weak-interaction equilibrium requires $`\mu _d\mu _u=\mu _e`$, where the subscripts refer to down quarks, up quarks, and electrons. But since a finite electron number is required to achieve electric charge neutrality in the star, $`\mu _e`$ cannot vanish and we necessarily have $`\mu _d\mu _u`$. This difference in chemical potentials is likely to drive a first order rather than a BCS-second order transition .
The other notable omission in our study is the strange quark. As before, weak interaction equilibrium requires that $`\mu _s=\mu _d`$, and thus for $`\mu _d\stackrel{>}{}\text{ }m_s`$ strange quarks will be present. Furthermore, the finite strange quark mass implies a mismatch in the Fermi momenta which would also drive a first order transition . A generic consequence of a first order transition is a mixed phase containing both normal and superconducting quark matter, and transport in the heterogeneous mixed phase is qualitatively different from that considered here, for neutrino scattering will depend on the size and nature of the structures (droplets) present. In previous work it was shown that neutrino mean free path in the heterogeneous phase can be greatly reduced due to coherent scattering of droplets in the mixed phase. Combining the results of Section II with neutrino transport in a mixed phase of superconducting and normal quark matter is beyond the scope of this work.
Therefore, while our calculation of the neutrino mean free path will be an essential ingredient in a more realistic and hence more complicated model of neutron star evolution, our toy model only applies to a highly idealized quark core of a neutron star where $`\mathrm{\Delta }|\mu _d\mu _u|`$.
## IV Conclusions
Motivated by the physical relevance of neutrino diffusion in the cooling of strongly-interacting matter, we have analyzed the effects on neutrino-quark scattering arising from a second-order phase transition from normal to color superconducting quark matter. The principal microscopic ingredient is the neutrino mean free path, and this was calculated in linear-response theory with a BCS-type superconducting background. We then applied this result to a schematic model of quark matter at temperatures and densities relevant to the evolution of proto-neutron stars. The modified mean free path for neutrinos in a superconducting background was computed in the simplified approximation of iso-symmetric, two-flavor quark matter. We have enumerated the main shortcomings of these simplifications and realize that a more realistic treatment of PNS evolution would include many other effects of similar importance.
Despite these caveats, our toy model calculation indicates that a superconducting transition in the quark core of a PNS can impact its early cooling and thereby potentially alter the temporal characteristics of the neutrino emission. We view this work as a first step towards an understanding of how the presence of color superconducting matter in the core of a neutron star may affect the early – and observable – neutrino signal. Quark pairing invariably occurs in theoretical treatments of finite-density QCD, and thus our microscopic calculations of the neutrino cross sections are pertinent to transport processes in dense quark matter. Given the real prospect of detecting neutrinos emitted from a future supernova event, such transport processes might someday serve as an probe of the properties of extremely dense matter.
## V Acknowledgments
We thank the organizers of the INT Program on QCD at Finite Baryon Density, during which this work was begun, and G.W.C. thanks the Institute for Nuclear Theory for their hospitality. We also thank David Kaplan, M. Prakash, Krishna Rajagopal, and Martin Savage for critical readings of the manuscript and for useful comments. This work was supported by the US Department of Energy grants DE-FG02-88ER40388 (G.W.C.) and DE-FG03-00-ER41132 (S.R.).
## Quark Polarization Tensors
At densities relevant to this work, only the quark particle-hole excitations are accessible. Thus we discard all anti-particle and anti-hole contributions from the propagators in Eqs. (7) and (9).
The imaginary part of the vector longitudinal polarization is, for each flavor,
$`\mathrm{Im}\mathrm{\Pi }_{00}^V(q_0,q)`$ $`=`$ $`2\pi {\displaystyle \frac{d^3p}{(2\pi )^3}\frac{d^3k}{(2\pi )^3}\delta ^3(pkq)(1+\widehat{p}\widehat{k})}`$ (27)
$`\times \{[n(\xi _p)n(\xi _k)][\delta (q_0+\xi _p\xi _k){\displaystyle \frac{(\xi _p+E_p)(\xi _k+E_k)\mathrm{\Delta }^2}{4\xi _p\xi _k}}`$
$`\delta (q_0\xi _p+\xi _k){\displaystyle \frac{(\xi _pE_p)(\xi _kE_k)\mathrm{\Delta }^2}{4\xi _p\xi _k}}]`$
$`+[1n(\xi _p)n(\xi _k)][\delta (q_0\xi _p\xi _k){\displaystyle \frac{(\xi _pE_p)(\xi _k+E_k)+\mathrm{\Delta }^2}{4\xi _p\xi _k}}`$
$`\delta (q_0+\xi _p+\xi _k){\displaystyle \frac{(\xi _p+E_p)(\xi _kE_k)+\mathrm{\Delta }^2}{4\xi _p\xi _k}}]\}.`$
The axial longitudinal differs only in the sign of the anomalous $`\mathrm{\Delta }^2`$ term, a consequence of coherence effects being different in different channels, and is given by
$`\mathrm{Im}\mathrm{\Pi }_{00}^A(q_0,q)`$ $`=`$ $`2\pi {\displaystyle \frac{d^3p}{(2\pi )^3}\frac{d^3k}{(2\pi )^3}\delta ^3(pkq)(1+\widehat{p}\widehat{k})}`$ (32)
$`\times \{[n(\xi _p)n(\xi _k)][\delta (q_0+\xi _p\xi _k){\displaystyle \frac{(\xi _p+E_p)(\xi _k+E_k)+\mathrm{\Delta }^2}{4\xi _p\xi _k}}`$
$`\delta (q_0\xi _p+\xi _k){\displaystyle \frac{(\xi _pE_p)(\xi _kE_k)+\mathrm{\Delta }^2}{4\xi _p\xi _k}}]`$
$`+[1n(\xi _p)n(\xi _k)][\delta (q_0\xi _p\xi _k){\displaystyle \frac{(\xi _pE_p)(\xi _k+E_k)\mathrm{\Delta }^2}{4\xi _p\xi _k}}`$
$`\delta (q_0+\xi _p+\xi _k){\displaystyle \frac{(\xi _p+E_p)(\xi _kE_k)\mathrm{\Delta }^2}{4\xi _p\xi _k}}]\}.`$
From these expressions, one may obtain $`\mathrm{\Pi }_{11}^{V,A}`$, $`\mathrm{\Pi }_{10}^{V,A}`$, and $`\mathrm{\Pi }_{01}^{V,A}`$ as specified in Eq. (12).
The transverse response functions have similar forms. The vector is
$`\mathrm{Im}\mathrm{\Pi }_{22}^V(q_0,q)`$ $`=`$ $`2\pi {\displaystyle \frac{d^3p}{(2\pi )^3}\frac{d^3k}{(2\pi )^3}\delta ^3(pkq)(1+\widehat{p}\widehat{k}2\widehat{p}_2\widehat{k}_2)}`$ (37)
$`\times \{[n(\xi _p)n(\xi _k)][\delta (q_0+\xi _p\xi _k){\displaystyle \frac{(\xi _p+E_p)(\xi _k+E_k)+\mathrm{\Delta }^2}{4\xi _p\xi _k}}`$
$`\delta (q_0\xi _p+\xi _k){\displaystyle \frac{(\xi _pE_p)(\xi _kE_k)+\mathrm{\Delta }^2}{4\xi _p\xi _k}}]`$
$`+[1n(\xi _p)n(\xi _k)][\delta (q_0\xi _p\xi _k){\displaystyle \frac{(\xi _pE_p)(\xi _k+E_k)\mathrm{\Delta }^2}{4\xi _p\xi _k}}`$
$`\delta (q_0+\xi _p+\xi _k){\displaystyle \frac{(\xi _p+E_p)(\xi _kE_k)\mathrm{\Delta }^2}{4\xi _p\xi _k}}]\},`$
and the axial,
$`\mathrm{Im}\mathrm{\Pi }_{22}^A(q_0,q)`$ $`=`$ $`2\pi {\displaystyle \frac{d^3p}{(2\pi )^3}\frac{d^3k}{(2\pi )^3}\delta ^3(pkq)(1+\widehat{p}\widehat{k}2\widehat{p}_2\widehat{k}_2)}`$ (42)
$`\times \{[n(\xi _p)n(\xi _k)][\delta (q_0+\xi _p\xi _k){\displaystyle \frac{(\xi _p+E_p)(\xi _k+E_k)+\mathrm{\Delta }^2}{4\xi _p\xi _k}}`$
$`\delta (q_0\xi _p+\xi _k){\displaystyle \frac{(\xi _pE_p)(\xi _kE_k)+\mathrm{\Delta }^2}{4\xi _p\xi _k}}]`$
$`+[1n(\xi _p)n(\xi _k)][\delta (q_0\xi _p\xi _k){\displaystyle \frac{(\xi _pE_p)(\xi _k+E_k)\mathrm{\Delta }^2}{4\xi _p\xi _k}}`$
$`\delta (q_0+\xi _p+\xi _k){\displaystyle \frac{(\xi _p+E_p)(\xi _kE_k)\mathrm{\Delta }^2}{4\xi _p\xi _k}}]\}.`$
Finally, there is a small but finite response in the mixed vector-axial channel. Since the neutrino tensor element $`L^{23}`$ is imaginary, we take the real part of the polarization. Antisymmetric in spin, it is
$`\mathrm{Re}\mathrm{\Pi }_{23}^{VA}(q_0,q)`$ $`=`$ $`2\pi {\displaystyle \frac{d^3p}{(2\pi )^3}\frac{d^3k}{(2\pi )^3}\delta ^3(pkq)(\widehat{p}_1\widehat{k}_1)}`$ (47)
$`\times \{[n(\xi _p)n(\xi _k)][\delta (q_0+\xi _p\xi _k){\displaystyle \frac{(\xi _p+E_p)(\xi _k+E_k)+\mathrm{\Delta }^2}{4\xi _p\xi _k}}`$
$`\delta (q_0\xi _p+\xi _k){\displaystyle \frac{(\xi _pE_p)(\xi _kE_k)+\mathrm{\Delta }^2}{4\xi _p\xi _k}}]`$
$`+[1n(\xi _p)n(\xi _k)][\delta (q_0\xi _p\xi _k){\displaystyle \frac{(\xi _pE_p)(\xi _k+E_k)\mathrm{\Delta }^2}{4\xi _p\xi _k}}`$
$`\delta (q_0+\xi _p+\xi _k){\displaystyle \frac{(\xi _p+E_p)(\xi _kE_k)\mathrm{\Delta }^2}{4\xi _p\xi _k}}]\}.`$
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# Numerical Simulations in Cosmology I: Methods
## 1 Introduction
Numerical simulations in cosmology have a long history and numerous important applications. Different aspects of the simulations including history of the subject were reviewed recently by Bertschinger (1998); see also Sellwood (1987). More detailed aspects of simulations were discussed by Gelb (1992), Gross (1997), and Kravtsov (1999).
Numerical simulations play a very significant role in cosmology. It all started in 60s (Aarseth, 1963) and 70s (Peebles, 1970; Press & Schechter, 1974) with simple N-body problems solved using N-body codes with few hundred particles. Later the Particle-Particle code (direct summation of all two-body forces) was polished and brought to the state-of-art (Aarseth, 1985). Already those early efforts brought some very valuable fruits. Peebles (1970) studied collapse of a cloud of particles as a model of cluster formation. The model had 300 points initially distributed within a sphere with no initial velocities. After the collapse and virialization the system looked like a cluster of galaxies. Those early simulations of cluster formation, though producing cluster-like objects, signaled the first problem – simple model of initially isolated cloud (top-hat model) results in the density profile of the cluster which is way too steep (power-law slope -4) as compared with real clusters (slope -3). The problem was addressed by Gunn & Gott (1972), who introduced a notion of secondary infall in an effort to solve the problem. Another keystone work of those times is the paper by White (1976), who studied collapse of 700 particles with different masses. It was shown that if one distributes the mass of a cluster to individual galaxies, two-body scattering will result in mass segregation not compatible with observed clusters. This was another manifestation of the dark matter in clusters. This time it was shown that inside a cluster the dark matter can not reside inside individual galaxies.
Survival of substructures in galaxy clusters was another problem addressed in the paper. It was found that lumps of dark matter, which in real life may represent galaxies, do not survive in dense environment of galaxy clusters. White & Rees (1978) argued that the real galaxies survive inside clusters because of energy dissipation by the baryonic component. That point of view was accepted for almost 20 years. Only recently it was shown the energy dissipation probably does not play a dominant role in survival of galaxies and the dark matter halos are not destroyed by tidal stripping and galaxy-galaxy collisions inside clusters (Klypin et al., 1999; Ghigna et al., 1999). The reason why early simulations came to a wrong result was pure numerical: they did not have enough resolution. But 20 years ago it was physically impossible to make a simulation with sufficient resolution. Even if at that time we had present-day codes, it would have taken about 600 years to make one run.
Generation of initial condition with given amplitude and spectrum of fluctuations was a problem for some time. The only correctly simulated spectrum was the flat spectrum which was generated by randomly distributing particles. In order to generate fluctuations with power spectrum, say $`P(k)k^1`$, Aarseth et al. (1979) placed particles along rods. Formally, it generates the spectrum, but the distribution has nothing to do with cosmological fluctuations. Doroshkevich et al. (1980) and Klypin & Shandarin (1983) were the first to use the Zeldovich (1970) approximation to set initial conditions. Since then this method is used to generate initial conditions for arbitrary initial spectrum of perturbations.
Starting mid 80s the field of numerical simulations is blooming: new numerical techniques are invented, old ones are perfected. The number of publications based on numerical modeling skyrocketed. To large extend, this have changed our way of doing cosmology. Instead of questionable assumptions and waving-hands arguments, we have tools of testing our hypothesis and models. As an example, I mention two analytical approximations which were validated by numerical simulations. The importance of both approximations is difficult to overestimate. The first is the Zeldovich approximation, which paved the way for understanding the large-scale structure of the galaxy distribution. The second is the Press & Schechter (1974) approximation, which gives the number of objects formed at different scales at different epochs. Both approximations cannot be formally proved. The Zeldovich approximation formally is not applicable for hierarchical clustering. It must start with smooth perturbations (truncated spectrum). Nevertheless, numerical simulations have shown that even for the hierarchical clustering the approximation can be used with appropriate filtering of initial spectrum (see Sahni & Coles, 1995, and references therein). The Press-Schechter approximation is also difficult to justify without numerical simulations. It operates with the initial spectrum and the linear theory, but then (a very long jump) it predicts the number of objects at very nonlinear stage. Because it is not based on any realistic theory of nonlinear evolution, it was an ingenious, but a wild guess. If anything, the approximation is based on a simple spherical top-hat model. But simulations show that objects do not form in this way – they are formed in a complicated fashion through multiple mergers and accretion along filaments. Still this a very simple and a very useful prescription gives quite accurate predictions.
This lecture is organized in the following way. Section 2 gives the equations which we solve to follow the evolution of initially small fluctuations. Initial conditions are discussed in section 3. A brief discussion of different methods is given in section 4. Effects of the resolution and some other technical details are also discussed in Section 5. Identification of halos (“galaxies”) is discussed in Section 6.
## 2 Equations of evolution of fluctuations in an expanding universe
Usually the problem of the formation and dynamics of cosmological objects is formulated as $`N`$-body problem: for $`N`$ point-like objects with given initial positions and velocities find their positions and velocities at any later moment. It should be remembered that this just a short-cut in our formulation – to make things simple. While it still mathematically correct in many cases, it does not give a correct explanation to what we do. If we are literally to take this approach, we should follow the motion of zillions of axions, baryons, neutrinos, and whatever else our Universe is made of. So, what it has to do with the motion of those few millions of particles in our simulations? The correct approach is to start with the Vlasov equation coupled with the Poisson equation and with appropriate initial and boundary conditions. If we neglect the baryonic component, which of course is very interesting, but would complicate our situation even more, the system is described by distribution functions $`f_i(𝐱,\dot{𝐱},t)`$ which should include all different clustered components $`i`$. For a simple CDM model we have only one component (axions or whatever it is). For more complicated Cold plus Hot Dark Matter (CHDM) with few different types of neutrinos the system includes one DF for the cold component and one DF for each type of neutrino (Klypin et al., 1993). In the comoving coordinates x the equations for the evolution of $`f_i`$ are:
$`{\displaystyle \frac{f_i}{t}}`$ $`+`$ $`\dot{𝐱}{\displaystyle \frac{f_i}{𝐱}}\varphi {\displaystyle \frac{f_i}{𝐩}}=0,𝐩=a^2\dot{𝐱},`$ (1)
$`^2\varphi `$ $`=`$ $`4\pi Ga^2(\rho (𝐱,t)\rho _{\mathrm{dm}}(t))=4\pi Ga^2\mathrm{\Omega }_{\mathrm{dm}}\delta _{\mathrm{dm}}\rho _{\mathrm{cr}},`$ (2)
$`\delta _{\mathrm{dm}}(𝐱,t)`$ $`=`$ $`(\rho _{\mathrm{dm}}\rho _{\mathrm{dm}})/\rho _{\mathrm{dm}}),`$ (3)
$`\rho _{\mathrm{dm}}(𝐱,t)`$ $`=`$ $`a^3{\displaystyle \underset{i}{}}m_i{\displaystyle d^3𝐩f_i(𝐱,\dot{𝐱},t)}.`$ (4)
Here $`a=(1+z)^1`$ is the expansion parameter, $`𝐩=a^2\dot{𝐱}`$ is the momentum, $`\mathrm{\Omega }_{\mathrm{dm}}`$ is the contribution of the clustered dark matter to the mean density of the Universe, $`m_i`$ is the mass of a particle of $`i`$th component of the dark matter. The solution of the Vlasov equation can be written in terms of equations for characteristics, which look like equations of particle motion:
$`{\displaystyle \frac{d𝐩}{da}}`$ $`=`$ $`{\displaystyle \frac{\varphi }{\dot{a}}},{\displaystyle \frac{d𝐯}{dt}}+2{\displaystyle \frac{\dot{a}}{a}}𝐯={\displaystyle \frac{\varphi }{a^3}}`$ (5)
$`{\displaystyle \frac{d𝐱}{da}}`$ $`=`$ $`{\displaystyle \frac{𝐩}{\dot{a}a^2}},{\displaystyle \frac{d𝐱}{dt}}=𝐯`$ (6)
$`^2\varphi `$ $`=`$ $`4\pi G\mathrm{\Omega }_0\delta _{\mathrm{dm}}\rho _{\mathrm{cr},0}/a,\varphi =a\varphi `$ (7)
$`\dot{a}`$ $`=`$ $`H_0\sqrt{1+\mathrm{\Omega }_0\left({\displaystyle \frac{1}{a}}1\right)+\mathrm{\Omega }_\mathrm{\Lambda }\left(a^21\right)}`$ (8)
In these equations $`\rho _{\mathrm{cr},0}`$ is the critical density at $`z=0`$; $`\mathrm{\Omega }_0`$, and $`\mathrm{\Omega }_{\mathrm{\Lambda },0}`$, are the density of the matter and of the cosmological constant in units of the critical density at $`z=0`$.
The distribution function $`f_i`$ is constant along each characteristic. This property should be preserved by numerical simulations. The complete set of characteristics coming through every point in the phase space is equivalent to the Vlasov equation. We can not have the complete (infinite) set, but we can follow the evolution of the system (with some accuracy), if we select a representative sample of characteristics. One way of doing this would be to split initial phase space into small domains, take only one characteristic as representative for each volume element, and to follow the evolution of the system of the “particles” in a self-consistent way. In models with one “cold” component of clustering dark matter (like CDM or $`\mathrm{\Lambda }`$CDM) the initial velocity is a unique function of coordinates (only “Zeldovich” part is present, no thermal velocities). This means that we need to split only coordinate space, not velocity space. For complicated models with significant thermal component, the distribution in full phase space should be taken into account. Depending on what we are interested in, we might split initial space into equal-size boxes (typical setup for PM or P<sup>3</sup>M simulations) or we could divide some area of interest (say, where a cluster will form) into smaller boxes, and use much bigger boxes outside the area (to mimic gravitational forces of the outside material). In any case, the mass assigned to a “particle” is equal to the mass of the domain it represents. Now we can think of the “particle” either as a small box, which moves with the flow, but does not change its original shape, or as a point-like particle. Both presentations are used in simulations. None is superior to another.
There are different forms of final equations. Mathematically they are all equivalent, but computationally there are very significant differences. There are considerations, which may affect the choice of particular form of the equations. Any numerical method gives more accurate results for a variable, which changes slowly with time. For example, for the gravitational potential we can choose either $`\varphi `$ or $`\varphi `$. At early stages of evolution perturbations still grow almost linearly. In this case we expect that $`\delta _{\mathrm{dm}}a`$, $`\varphi const`$, and $`\varphi a`$. Thus, $`\varphi `$ can be a better choice because it does not change. That is especially helpful, if code uses gravitational potential from previous moment of time as initial “guess” for current moment, as it happens in the case of the ART code. In any case, it is better to have a variable, which does not change much. For equations of motion we can choose, for example, either first equations in eqs.(5– 6) or the second equations. If we choose “momentum” $`p=a^2\dot{x}`$ as effective velocity and take the expansion parameter $`a`$ as time variable, then for the linear growth we expect that the change of coordinates per each step is constant: $`\mathrm{\Delta }x\mathrm{\Delta }a`$. Numerical integration schemes should not have problem with this type of growth. For the $`t`$ and $`v`$ variable, the rate of change is more complicated: $`\mathrm{\Delta }xa^{1/2}\mathrm{\Delta }t`$, which may produce some errors at small expansion parameters. The choice of variables may affect the accuracy of the solution even at very nonlinear stage of evolution as was argued by Quinn et al. (1997).
## 3 Initial Conditions
### 3.1 Zeldovich approximation
The Zeldovich approximation is commonly used to set initial conditions. The approximation is valid in mildly nonlinear regime and is much superior to the linear approximation. We slightly rewrite the original version of the approximation to incorporate cases (like CHDM) when the growth rates $`g(t)`$ depend on the wavelength of the perturbation $`|𝐤|`$. In the Zeldovich approximation the comoving and the lagrangian coordinates are related in the following way:
$$𝐱=𝐪\alpha \underset{𝐤}{}g_{|𝐤|}(t)𝐒_{|𝐤|}(𝐪),𝐩=\alpha a^2\underset{𝐤}{}g_{|𝐤|}(t)\left(\frac{\dot{g}_{|𝐤|}}{g_{|𝐤|}}\right)𝐒_{|𝐤|}(𝐪),$$
(9)
where the displacement vector $`𝐒`$ is related to the velocity potential $`\mathrm{\Phi }`$ and the power spectrum of fluctuations $`P(|𝐤|)`$:
$$𝐒_{|𝐤|}(𝐪)=_q\mathrm{\Phi }_{|𝐤|}(𝐪),\mathrm{\Phi }_{|𝐤|}=\underset{𝐤}{}a_𝐤\mathrm{cos}(\mathrm{𝐤𝐪})+b_𝐤\mathrm{sin}(\mathrm{𝐤𝐪}),$$
(10)
where $`a`$ and $`b`$ are gaussian random numbers with the mean zero and dispersion $`\sigma ^2=P(k)/k^4`$:
$$a_𝐤=\sqrt{P(|𝐤|)}\frac{Gauss(0,1)}{|𝐤|^2},b_𝐤=\sqrt{P(|𝐤|)}\frac{Gauss(0,1)}{|𝐤|^2}.$$
(11)
The parameter $`\alpha `$, together with the power spectrum $`P(k)`$, define the normalization of the fluctuations.
In order to set the initial conditions, we choose the size of the computational box $`L`$ and the number of particles $`N^3`$. The phase space is divided into small equal cubes of size $`2\pi /L`$. Each cube is centered on a harmonic $`𝐤=2\pi /L\times \{i,j,k\}`$, where $`\{i,j,k\}`$ are integer numbers with limits from zero to $`N/2`$. We make a realization of the spectrum of perturbations $`a_𝐤`$ and $`b_𝐤`$, and find displacement and momenta of particles with $`𝐪=L/N\times \{i,j,k\}`$ using eq.(9). Here $`i,j,k=1,N`$.
### 3.2 Power Spectrum
There are approximations of the power spectrum $`P(k)`$ for a wide range of cosmological models. Publicly available COSMICS code (Bertschinger 1996) gives accurate approximations for the power spectrum. Here we follow Klypin & Holtzman (1997) who give the following fitting formula:
$$P(k)=\frac{k^n}{(1+P_2k^{1/2}+P_3k+P_4k^{3/2}+P_5k^2)^{2P_6}}.$$
(12)
The coefficients $`P_i`$ are presented by Klypin & Holtzman (1997) for a variety of models. The comparison of some of the power spectra with the results from COSMICS (Bertschinger 1996) indicate that the errors of the fits are smaller than 5%. Table 1 gives parameters of the fits for some popular models.
The power spectrum of cosmological models is often approximated using a fitting formula given by Bardeen et al. (1986, BBKS):
$$P(k)=k^nT^2(k),T(k)=\frac{\mathrm{ln}(1+2.34q)}{2.34q}[1+3.89q+(16.1q)^2+(5.4q)^3+(6.71q)^4]^{1/4},$$
(13)
where $`q=k/(\mathrm{\Omega }_0h^2\mathrm{Mpc}^1)`$. Unfortunately, the accuracy of this approximation is not great and it should not be used for accurate simulations. We find that the following approximation, which is a combination of a slightly modified BBKS fit and the Hu & Sugiyama (1996) scaling with the amount of baryons, provides errors in the power spectrum smaller than 5% for the range of wavenumbers $`k=(10^440)h\mathrm{Mpc}^1`$ and for $`\mathrm{\Omega }_b/\mathrm{\Omega }_0<0.1`$:
$`P(k)`$ $`=`$ $`k^nT^2(k),`$
$`T(k)`$ $`=`$ $`{\displaystyle \frac{\mathrm{ln}(1+2.34q)}{2.34q}}[1+13q+(10.5q)^2+(10.4q)^3+(6.51q)^4]^{1/4},`$
$`q`$ $`=`$ $`{\displaystyle \frac{k(T_{\mathrm{CMB}}/2.7K)^2}{\mathrm{\Omega }_0h^2\alpha ^{1/2}(1\mathrm{\Omega }_b/\mathrm{\Omega }_0)^{0.60}}},\alpha =a_1^{\mathrm{\Omega }_b/\mathrm{\Omega }_0}a_2^{(\mathrm{\Omega }_b/\mathrm{\Omega }_0)^3}`$
$`a_1`$ $`=`$ $`(46.9\mathrm{\Omega }_0h^2)^{0.670}[1+(32.1\mathrm{\Omega }_0h^2)^{0.532}],a_2=(12\mathrm{\Omega }_0h^2)^{0.424}[1+(45\mathrm{\Omega }_0h^2)^{0.582}]`$ (14)
### 3.3 Multiple masses: high resolution for a small region
In many cases we would like to set initial conditions in such a way that inside some specific region(s) there are more particles and the spectrum is better resolved. We need this when we want to have high resolution for a halo, but we also need the environment of the halo. This is done in a two-step process. First, we run a low resolution simulation which has a sufficiently large volume to include the effects of the environment. For this run all the particles have the same mass. A halo is picked for rerunning with high resolution. Second, using particles of the halo, we identify region in the lagrangian (initial) space, where the resolution should be increased. We add high-frequency harmonics, which are not present in the low resolution run. We then add contributions of all the harmonics and get initial displacements and momenta (eq. 9). Let’s be more specific. In order to add the new harmonics, we must specify (1) how we divide the phase space and place the harmonics, and (2) how we sum the contributions of the harmonics.
The simplest way is to divide the phase space into many small boxes of size $`2\pi /L`$, where $`L`$ is the box size. This is the same devision, which we use to set the low resolution run. But now we extend it to very high frequencies up to $`2\pi /L\times N/2`$, where $`N`$ is the new effective number of particles. For example, we used $`N=64`$ for the low resolution run. For high resolution run we may choose $`N=1024`$. Simply replace the value and run the code again. Of course, we really can not do it because it would generate too many particles. Instead, in some regions, where the resolution should not be high, we combine particles together (by taking average coordinates and average velocities) and replace many small-mass particles with fewer larger ones. Left panel in Figure 1 gives an example of mass refinement. Note that we try to avoid too large jumps in the mass resolution by creating layers of particles of increasing mass.
This approach is correct and relatively simple. It may seem that it takes too much cpu to get the initial conditions. In practice, cpu time is not much of an issue because initial conditions are generated only once and it takes only few cpu hours even for $`1024^3`$ mesh. For most of applications $`1024^3`$ particles is more then enough. The problem arises when we want to have more then $`1024^3`$ particles. We simply do not have enough computer memory to store the information for all the harmonics. In this case we must decrease the resolution in the phase space. It is a bit easier to understand the procedure, if we consider phase space diagrams like one presented in Figure 3. The low resolution run in this case was done for $`32^3`$ particles with harmonics up to $`16\times 2\pi /L`$ (small points). For the high resolution run we choose a region of size 1/8 of the original large box. Inside the small box we place another box, which is twice smaller. Thus, we will have three levels of mass refinement. For each level we have corresponding size of the phase space block. The size is defined by the size of real space box and is equal to $`2\pi /L\times K`$, $`K=1,8,16`$. Harmonics from different refinements should not overlap: if a region in phase space is represented on lower level of resolution, it should not appear in the the higher resolution level. This is why rows of the highest resolution harmonics (circles) with $`K_x=16`$ and $`K_y=16`$ are absent in the Figure 3: they have been already covered by lower resolution blocks marked by stars. Figure 3 clearly illustrate that matching harmonics is a complicated process: we failed to do the match because there are partially overlapping blocks and there are gaps. We can get much better results, if we assume different ratios of the sizes of the boxes. For example, if instead of box ratios $`1:(1/8):(1/16)`$, we chose ratios $`1:(3/32):(5/96)`$, the coverage of the phase space is almost perfect as shown in Figure: 2.
## 4 Codes
There are many different numerical techniques to follow the evolution of a system of many particles. For earlier reviews see Hockney & Eastwood (1981); Sellwood (1987), and Bertschinger (1998). Most of the methods for cosmological applications take some ideas from three techniques: Particle Mesh (PM) code, direct summation or Particle-Particle code, and the TREE code. For example, the Adaptive Particle-Particle/Particle-Mesh (AP<sup>3</sup>M) code (Couchman, 1991) is a combination of the PM code and the Particle-Particle code. The Adaptive-Refinement-Tree code (ART) (Kravtsov et al., 1997; Kravtsov, 1999) is an extension of the PM code with the organization of meshes in the form of a tree. All methods have their advantages and disadvantages.
PM code. It uses a mesh to produce density and potential. As the result, its resolution is limited by the size of the mesh. There are two advantages of the method: i) it is fast (the smallest number of operations per particle per time step of all the other methods), ii) it typically uses very large number of particles. The later can be crucial for some applications. There are few modifications of the code. “Plain-vanilla” PM was described by (Hockney & Eastwood, 1981). It includes Cloud-In-Cell density assignment and 7-point discrete analog of the Laplacian operator. Higher order approximations improve the accuracy on large distances, but degrades the resolution (e.g., Gelb, 1992). The PM code is available (Klypin & Holtzman, 1997)
P<sup>3</sup>M code is described in detail in Hockney & Eastwood (1981) and Efstathiou et al. (1985). It has two parts: PM part, which takes care of large-scale forces, and PP part, which adds small-scale particle-particle contribution. Because of strong clustering at late stages of evolution, PP part becomes prohibitively expensive once large objects start to form in large numbers. Significant speed is achieved in modified version of the code, which introduces sub-grids (next levels of PM) in areas with high density (Couchman, 1991). With modification the code is as fast as the TREE code even for heavily clustered configurations. The code express the inter-particle force as a sum of a short range force (computed by direct particle-particle pair force summation) and the smoothly varying part (approximated by the particle-mesh force calculation). One of the major problems for these codes is the correct splitting of the force into a short-range and a long-range part. The grid method (PM) is only able to produce reliable inter particle forces down to a minimum of at least two grid cells. For smaller separations the force can no longer be represented on the grid and therefore one must introduce a cut-off radius $`r_e`$ (larger than two grid cells), where for $`r<r_e`$ the force should smoothly go to zero. The parameter $`r_e`$ defines the chaining-mesh and for distances smaller than this cutoff radius $`r_e`$ a contribution from the direct particle-particle (PP) summation needs to be added to the total force acting on each particle. Again this PP force should smoothly go to zero for very small distances in order to avoid unphysical particle-particle scattering. This cutoff of the PP force determines the overall force resolution of a P<sup>3</sup>M code.
The most widely used version of this algorithm is currently the adaptive P<sup>3</sup>M (AP<sup>3</sup>M) code of Couchman (1991), which is available for public.. The smoothing of the force in this code is connected to a $`S_2`$ sphere, as described in Hockney & Eastwood (1981).
TREE code is the most flexible code in the sense of the choice of boundary conditions (Appel, 1985; Barnes & Hut, 1986; Hernquist, 1987). It is also more expensive than PM: it takes 10-50 times more operations. (Bouchet & Hernquist, 1988) and (Hernquist et al., 1991) extended the code for the periodical boundary conditions, which is important for simulating large-scale fluctuations. Some variants of the TREE are publicly available. There are variants of the code modified for massively parallel computers. There code variants with variable time stepping, which is vital for extremely high resolution simulations.
ART code. Multigrid methods were introduced long ago, but only recently they started to show a potential to produce real results. It worth of paying attention if a “multigrid” code is really a fully adaptive multigrid code. An example of this type of the codes is the Adaptive Refinement Tree code (ART; Kravtsov et al. 1997), which reaches high force resolution by refining all high-density regions with an automated refinement algorithm. The refinements are recursive: the refined regions can also be refined, each subsequent refinement having half of the previous level’s cell size. This creates an hierarchy of refinement meshes of different resolutions covering regions of interest. The refinement is done cell-by-cell (individual cells can be refined or de-refined) and meshes are not constrained to have a rectangular (or any other) shape. This allows the code to refine the required regions in an efficient manner. The criterion for refinement is the local overdensity of particles the code refines an individual cell only if the density of particles (smoothed with the cloud-in-cell scheme; Hockney & Eastwood 1981) is higher than $`n_{th}`$ particles, with typical values $`n_{th}=25`$. The Poisson equation on the hierarchy of meshes is solved first on the base grid using FFT technique and then on the subsequent refinement levels. On each refinement level the code obtains the potential by solving the Dirichlet boundary problem with boundary conditions provided by the already existing solution at the previous level or from the previous moment of time. There is no particle-particle summation in the ART code and the actual force resolution is equal to $`2`$ cells of the finest refinement mesh covering a particular region.
Figure 4 (courtesy of A. Kravtsov) gives an example of mesh refinement for hydro-dynamical version of the ART code. The code produced this refinement mesh for spherical strong explosion (Sedov solution).
The refinement of the time integration mimics spatial refinement and the time step for each subsequent refinement level is two times smaller than the step on the previous level. Note, however, that particles on the same refinement level move with the same step. When a particle moves from one level to another, the time step changes and its position and velocity are interpolated to appropriate time moments. This interpolation is first-order accurate in time, whereas the rest of the integration is done with the second-order accurate time centered leap-frog scheme. All equations are integrated with the expansion factor $`a`$ as a time variable and the global time step hierarchy is thus set by the step $`\mathrm{\Delta }a_0`$ at the zeroth level (uniform base grid). The step on level $`L`$ is then $`\mathrm{\Delta }a_L=\mathrm{\Delta }a_0/2^L`$.
What code is the best? Which one to choose? There is no unique answer – everything depends on the problem, which we are addressing. For example, if we are interested in explanation of the large-scale structure (filaments, voids, Zeldovich approximation, and so on), PM code with 256<sup>3</sup> mesh is sufficient. It takes only one night to make a simulation on a (good) workstation. There is a very long list of problems like that. But if you intent to look for the structure of individual galaxies in the large-scale environment, you must have a code with much better resolution with variable time stepping, and with multiple masses. In this case the TREE or ART codes are the choices.
## 5 Effects of resolution
As the resolution of simulations improves and the range of their applications broaden, it becomes increasingly important to understand the limits of the simulations. Knebe et al. (1999) made detailed comparison of realistic simulations done with three codes: ART, AP<sup>3</sup>M, and PM. Here we present some of their results and main conclusions. The simulations were done for the standard CDM model with the dimensionless Hubble constant $`h=0.5`$ and $`\mathrm{\Omega }_0=1`$. The simulation box of $`15h^1\mathrm{Mpc}`$ had $`64^3`$ equal-mass particles, which gives the mass resolution (mass per particle) of $`3.55\times 10^9h^1M_{}`$. Because of the low resolution of the PM runs, we show results only for the other two codes. For the ART code the force resolution is practically fixed by the number of particles. The only free parameter is the number of steps on the lowest (zero) level of resolution. In the case of the AP<sup>3</sup>M, besides the number of steps, one can also request the force resolution. Parameters of two runs with the ART code and five simulations with the AP<sup>3</sup>M are given in Table 2.
Figure 5 shows the correlation function for the dark matter down to the scale of $`5h^1\mathrm{kpc}`$, which is close to the force resolution of all our high-resolution simulations. The correlation function in runs AP<sup>3</sup>M<sub>1</sub> and ART<sub>2</sub> are similar to those of AP<sup>3</sup>M<sub>5</sub> and ART<sub>1</sub> respectively and are not shown for clarity. We can see that the AP<sup>3</sup>M<sub>5</sub> and the ART<sub>1</sub> runs agree to $`\genfrac{}{}{0pt}{}{_<}{^{}}10\%`$ over the whole range of scales. The correlation amplitudes of runs AP<sup>3</sup>M<sub>2-4</sub>, however, are systematically lower at $`r\genfrac{}{}{0pt}{}{_<}{^{}}5060h^1\mathrm{kpc}`$ (i.e., the scale corresponding to $`1520`$ resolutions), with the AP<sup>3</sup>M<sub>3</sub> run exhibiting the lowest amplitude. The fact that the AP<sup>3</sup>M<sub>2</sub> correlation amplitude deviates less than that of the AP<sup>3</sup>M<sub>3</sub> run, indicates that the effect is very sensitive to the force resolution.
Note that the AP<sup>3</sup>M<sub>3</sub> run has formally the best force resolution. Thus, one would naively expect that it should gives the largest correlation function. At scales $`\genfrac{}{}{0pt}{}{_<}{^{}}30h^1\mathrm{kpc}`$ the deviations of the AP<sup>3</sup>M<sub>3</sub> from the ART<sub>1</sub> or the AP<sup>3</sup>M<sub>5</sub> runs are $`100200\%`$. We attribute these deviations to the numerical effects: high force resolution in AP<sup>3</sup>M<sub>3</sub> was not adequately supported by the time integration. In other words, the AP<sup>3</sup>M<sub>3</sub> had too few time-steps. Note that it had a quite large number of steps (6000), not much smaller than the AP<sup>3</sup>M<sub>5</sub> (8000). But for its force resolution, it should have many more steps. The lack of the number of steps was devastating.
Figure 6 presents the density profiles of four of the most massive halos in our simulations. We have not shown the profile of the most massive halo because it appears to have undergone a recent major merger and is not very relaxed. In this figure, we present only profiles of halos in the high-resolution runs. Not surprisingly, the inner density of the PM halos is much smaller than in the high-resolution runs and their profiles deviate strongly from the profiles of high-resolution halos at the scales shown in Figure 6. A glance at Figure 6 shows that all profiles agree well at $`r\genfrac{}{}{0pt}{}{_>}{^{}}30h^1\mathrm{kpc}`$. This scales is about eight times smaller than the mean inter-particle separation. Thus, despite the very different resolutions, time steps, and numerical techniques used for the simulations, the convergence is observed at a scale much lower than the mean inter-particle separation, argued by Splinter et al. (1998) to be the smallest trustworthy scale.
Nevertheless, there are systematic differences between the runs. The profiles in two ART runs are identical within the errors indicating convergence (we have run an additional simulation with time steps twice smaller than those in the ART<sub>1</sub> finding no difference in the density profiles). Among the AP<sup>3</sup>M runs, the profiles of the AP<sup>3</sup>M<sub>1</sub> and AP<sup>3</sup>M<sub>5</sub> are closer to the density profiles of the ART halos than the rest. The AP<sup>3</sup>M<sub>2</sub>, AP<sup>3</sup>M<sub>3</sub>, and AP<sup>3</sup>M<sub>4</sub>, despite the higher force resolution, exhibit lower densities in the halo cores, the AP<sup>3</sup>M<sub>3</sub> and AP<sup>3</sup>M<sub>4</sub> runs being the most deviant.
These results can be interpreted, if we examine the trend of the central density as a function of the ratio of the number of time steps to the dynamic range of the simulations (see Table 2). The ratio is smaller when either the number of steps is smaller or the force resolution is higher. The agreement in density profiles is observed when this ratio is $`\genfrac{}{}{0pt}{}{_>}{^{}}2`$. This suggests that for a fixed number of time steps, there should be a limit on the force resolution. Conversely, for a given force resolution, there is a lower limit on the required number of time steps. The exact requirements would probably depend on the code type and the integration scheme. For the AP<sup>3</sup>M code our results suggest that the ratio of the number of time steps to the dynamic range should be no less than one. It is interesting that the deviations in the density profiles are similar to and are observed at the same scales as the deviations in the DM correlation function (Fig. 5) suggesting that the correlation function is sensitive to the central density distribution of dark matter halos.
## 6 Halo identification
Finding halos in dense environments is a challenge. Some of the problems that any halo finding algorithm faces are not numerical. They exist in the real Universe. We select a few typical difficult situations.
1. A large galaxy with a small satellite. Examples: LMC and the Milky Way or the M51 system. Assuming that the satellite is bound, do we have to include the mass of the satellite in the mass of the large galaxy? If we do, then we count the mass of the satellite twice: once when we find the satellite and then when we find the large galaxy. This does not seem reasonable. If we do not include the satellite, then the mass of the large galaxy is underestimated. For example, the binding energy of a particle at the distance of the satellite will be wrong. The problem arises when we try to assign particles to different halos in an effort to find masses of halos. This is very difficult to do for particles moving between halos. Even if a particle at some moment has negative energy relative to one of the halos, it is not guaranteed that it belongs to the halo. The gravitational potential changes with time, and the particle may end up falling onto another halo. This is not just a precaution. This actually was found very often in real halos when we compared contents of halos at different redshifts. Interacting halos exchange mass and lose mass. We try to avoid the situation: instead of assigning mass to halos, we find the maximum of the “rotational velocity”, $`\sqrt{GM/R}`$, which is observationally a more meaningful quantity.
2. A satellite of a large galaxy. The previous situation is now viewed from a different angle. How can we estimate the mass or the rotational velocity of the satellite? The formal virial radius of the satellite is large: the big galaxy is within the radius. The rotational velocity may rise all the way to the center of the large galaxy. In order to find the outer radius of the satellite, we analyze the density profile. At small distances from the center of the satellite the density steeply declines, but then it flattens out and may even increase. This means that we reached the outer border of the satellite. We use the radius at which the density starts to flatten out as the first approximation for the radius of the halo. This approximation can be improved by removing unbound particles and checking the steepness of the density profile in the outer part.
3. Tidal stripping. Peripheral parts of galaxies, responsible for extended flat rotation curves outside of clusters, are very likely tidally stripped and lost when the galaxies fall into a cluster. The same happens with halos: a large fraction of halo mass may be lost due to stripping in dense cluster environments. Thus, if an algorithm finds that 90% of mass of a halo identified at early epoch is lost, it does not mean that the halo was destroyed. This is not a numerical effect and is not due to “lack of physics”. This is a normal situation. What is left of the halo, given that it still has a large enough mass and radius, is a “galaxy”.
There are different methods of identifying collapsed objects (halos) in numerical simulations.
Friends-Of-Friends (FOF) algorithm was used a lot and still has its adepts. If we imagine that each particle is surrounded by a sphere of radius $`bd/2`$, then every connected group of particles is identified as a halo. Here $`d`$ is the mean distance between particles, and $`b`$ is called linking parameter, which typically is 0.2. Dependence of groups on $`b`$ is extremely strong. The method stems from an old idea to use percolation theory to discriminate between cosmological models. Because of that, FOF is also called percolation method, which is wrong because the percolation is about groups spanning the whole box, not collapsed and compact objects. FOF was criticized for failing to find separate groups in cases when those groups were obviously present (Gelb, 1992). The problem originates from the tendency of FOF to “percolate” through bridges connecting interacting galaxies or galaxies in high density backgrounds.
DENMAX tried to overcome the problems of FOF by dealing with density maxima (Gelb, 1992; Bertschinger & Gelb, 1991). It finds maxima of density and then tries to identify particles, which belong to each maximum (halo). The procedure is quite complicated. First, density field is constructed. Second, the density (with negative sign) is treated as potential in which particles start to move as in a viscous fluid. Eventually, particles sink at bottoms of the potential (which are also maxima density). Third, only particles with negative energy (relative to their group) are retained. Just as in the case of FOF, we can easily imagine situations when (this time) DENMAX should fail. For example, two colliding galaxies in a cluster of galaxies. Because of large relative velocity they should just pass each other. In the moment of collision DENMAX ceases to “see” both galaxies because all particle have positive energies. That is probably a quite unlikely situation. The method is definitely one of the best at present. The only problem is that it seems to be too complicated for present state of simulations. DENMAX has two siblings – SKID (Stadel et al.) and BDM (Klypin & Holtzman, 1997) – which are frequently used.
“Overdensity 200”. There is no name for the method, but it is often used. Find density maximum, place a sphere and find radius, within which the sphere has the mean overdensity 200 (or 177 if you really want to follow the top-hat model of nonlinear collapse).
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# UNITARITY CONSTRAINTS ON NEUTRINO MASS AND MIXINGS
## 1 Introduction
Recent observations of atmospheric neutrinos and especially their zenith-angle dependence $`^\mathrm{?}`$, strongly suggest that muon neutrinos maximally mix with the tau neutrinos. Motivated by this observation and recent theoretical work on neutrino mass models $`^\mathrm{?}`$, we explored $`^\mathrm{?}`$ the implications of imposing the constraint that two neutrino flavors (which for definiteness we take to be $`\nu _\mu `$ and $`\nu _\tau `$) are similarly coupled to the mass basis in addition to the unitarity constraints.
Although the invisible width of the Z particle constraints the number of active neutrino flavors to be three, it is nevertheless worthwhile to consider the possibility of the existence of sterile neutrino states for a number of reasons: i) The possibility of oscillation of atmospheric muon neutrinos into sterile states is not completely ruled out. ii) If the LSND results $`^\mathrm{?}`$ are confirmed, since the analysis of LSND, atmospheric $`^\mathrm{?}`$ and solar $`^\mathrm{?}`$ neutrinos point out to different mass scales, one needs to introduce sterile neutrinos. iii) Serious problems such as the abundance of alpha particles that arise when core-collapse supernovae with neutrino-driven wind are considered as sites of r-process nucleosynthesis can be avoided by the oscillations of active neutrinos into sterile ones $`^{\mathrm{?},\mathrm{?}}`$. Even though cosmological and astrophysical bounds rule out heavier sterile states $`^\mathrm{?}`$, the effect of the lighter sterile neutrinos on big-bang nucleosynthesis is controversial $`^\mathrm{?}`$.
Hence we consider three active flavors and an arbitrary number (which could be taken to be zero) of sterile neutrinos. The $`N\times N`$ neutrino mixing matrix will be denoted by $`U_{\alpha i}`$ where $`\alpha `$ denotes the flavor index and $`i`$ denotes the mass index:
$$|\nu _\alpha =\underset{i}{}U_{\alpha i}|\nu _i.$$
(1)
We impose the constraint that $`U_{\mu i}`$ and $`U_{\tau i}`$ are proportional for all but one mass eigenstate, which we choose for definiteness to be the third mass eigenstate:
$$U_{\mu i}U_{\tau i}0,i3.$$
(2)
We write this condition in terms of an arbitrary angle $`\varphi `$ and an arbitrary phase $`\eta `$:
$$\mathrm{sin}\varphi U_{\mu i}=e^{i\eta }\mathrm{cos}\varphi U_{\tau i}0,i3.$$
(3)
Note that, in our formalism, we permit CP-violating phases. Introducing the quantity
$$A=\underset{i3}{}\left[|U_{\mu i}|^2+|U_{\tau i}|^2\right]$$
(4)
and using Eq. (3) along with the unitarity of the mixing matrix one can easily show that $`^\mathrm{?}`$
$$A=1,$$
(5)
$$U_{\mu 3}=\mathrm{sin}\varphi e^{i\delta }e^{i\eta },$$
(6)
$$U_{\tau 3}=\mathrm{cos}\varphi e^{i\delta },$$
(7)
where $`\delta `$ is a phase to be determined, and
$$U_{\alpha 3}=0,\alpha \mu ,\tau ,$$
(8)
Introducing the states
$$|\stackrel{~}{\nu }_\mu =\mathrm{cos}\varphi |\nu _\mu +\mathrm{sin}\varphi e^{i\eta }|\nu _\tau ,$$
(9)
and
$$|\stackrel{~}{\nu }_\tau =\mathrm{sin}\varphi e^{i\eta }|\nu _\mu +\mathrm{cos}\varphi |\nu _\tau ,$$
(10)
It follows that
$$|\stackrel{~}{\nu }_\mu =\frac{1}{\mathrm{cos}\varphi }\underset{i3}{}U_{\mu i}|\nu _i,$$
(11)
$$|\stackrel{~}{\nu }_\tau =e^{i\delta }|\nu _3.$$
(12)
and
$$|\nu _\alpha =\underset{i3}{}U_{\alpha i}|\nu _i,\alpha \mu ,\tau .$$
(13)
This is a remarkable result which simply follows from the assumption of Eq. (3). This assumption leads to a decoupling of all the other flavors from the chosen (the third in our choice) mass eigenstate in the neutrino mixing matrix.
### Three Active Flavors
For three active flavors we get
$$\left(\begin{array}{c}|\nu _e\\ |\stackrel{~}{\nu }_\mu \end{array}\right)=\left(\begin{array}{cc}U_{e1}& U_{e2}\\ U_{\mu 1}/\mathrm{cos}\varphi & U_{\mu 2}/\mathrm{cos}\varphi \end{array}\right)\left(\begin{array}{c}|\nu _1\\ |\nu _2\end{array}\right).$$
(14)
The solar neutrino data in this case could be explained by either the matter-enhanced or vacuum $`\nu _e\stackrel{~}{\nu }_\mu `$ oscillations.
In the special case of $`\varphi =\pi /4`$, the full mixing matrix is given by $`^\mathrm{?}`$
$$\left(\begin{array}{ccc}\mathrm{cos}\theta & \mathrm{sin}\theta & 0\\ \sqrt{2}\mathrm{sin}\theta & \sqrt{2}\mathrm{cos}\theta & \frac{1}{\sqrt{2}}e^{i\delta }\\ \sqrt{2}\mathrm{sin}\theta & \sqrt{2}\mathrm{cos}\theta & \frac{1}{\sqrt{2}}e^{i\delta }\end{array}\right).$$
(15)
The limiting case of $`\theta =\pi /4`$ and $`\delta =0`$ yields bi-maximal mixing of three active neutrinos $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$.
### An Arbitrary Number of Flavors
In general $`N`$ flavors mix with the fundamental representation of $`U(N)`$. An arbitrary $`U(N)`$ element can be written as a product of $`N(N1)/2`$ different non-commuting $`SU(2)`$ rotations and a diagonal matrix:
$$U_{i\alpha }^{}=R_{12}R_{13}R_{14}\mathrm{}R_{23}R_{24}\mathrm{}\left(\begin{array}{cccc}e^{i\delta _1}& 0& 0& .\\ 0& e^{i\delta _2}& 0& .\\ 0& 0& e^{i\delta _3}& .\\ .& .& .& .\end{array}\right)$$
(16)
where e.g.
$$R_{14}=\left(\begin{array}{ccccc}C_{14}& 0& 0& S_{14}^{}& .\\ 0& 1& 0& 0& .\\ 0& 0& 1& 0& .\\ S_{14}& 0& 0& C_{14}^{}& .\\ .& .& .& .& .\end{array}\right).$$
(17)
Our choice of parameters leads to
$`C_{\alpha 3}`$ $`=`$ $`1,\alpha 2`$
$`C_{23}`$ $`=`$ $`\mathrm{cos}\varphi `$
$`S_{23}`$ $`=`$ $`e^{i\eta }\mathrm{sin}\varphi ,`$ (18)
hence our choice reduces the number of parameters from $`N(N1)/2`$ to $`(N^23N+4)/2`$.
## 2 Specific Cases
Here we summarize implications of our scheme for three different experimental situations.
### 2.1 Atmospheric Neutrinos
If we have only active neutrinos with $`m_1m_2`$ we have the standard result:
$$P(\nu _\mu \nu _\tau )=\mathrm{sin}^22\varphi \mathrm{sin}^2\left[\frac{(m_3^2m_2^2)L}{4E}\right]$$
(19)
If we have only one sterile state in addition to the active neutrinos and there is the mass hierarchy $`m_4>m_3m_2>m_1`$ (where $`m_2^2m_1^2`$ is of order of the solar neutrino solution) we get the following result for the $`\nu _\mu \nu _\tau `$ conversion probability:
$`P`$ $`(\nu _\mu \nu _\tau )=\mathrm{sin}^22\varphi \mathrm{sin}^2\left[{\displaystyle \frac{(m_3^2m_2^2)L}{4E}}\right]+4|U_{\mu 4}|^4\mathrm{tan}^2\varphi \mathrm{sin}^2\left[{\displaystyle \frac{(m_4^2m_2^2)L}{4E}}\right]`$ (20)
$``$ $`8\mathrm{sin}^2\varphi |U_{\mu 4}|^2\mathrm{sin}\left[{\displaystyle \frac{(m_4^2m_2^2)L}{4E}}\right]\mathrm{sin}\left[{\displaystyle \frac{(m_3^2m_2^2)L}{4E}}\right]\mathrm{cos}\left[{\displaystyle \frac{(m_4^2m_3^2)L}{2E}}\right].`$
It will be instructive to do a fit to the SuperKamiokande atmospheric neutrino data with Eq. (20).
### 2.2 Reactor Neutrinos
In our scheme, if the value of $`(m_2^2m_1^2)`$ is determined from the solar neutrino data, for reactor neutrino experiments we can assume $`(m_2^2m_1^2)E/L`$. We then have
$$P(\nu _e\nu _e)=1\mathrm{sin}^22\theta _{\mathrm{eff}}\mathrm{sin}^2\left[\frac{(m_4^2m_1^2)L}{4E}\right],$$
(21)
where
$$\mathrm{sin}^22\theta _{\mathrm{eff}}=4|U_{e4}|^2(1|U_{e4}|^2).$$
(22)
For the large values of $`(m_4^2m_1^2)`$ that would help the r-process nucleosynthesis in the neutrino-driven wind models of supernova CHOOZ experiment gives a bound of $`|U_{e4}|^2<0.047`$ $`^\mathrm{?}`$. The best limit, $`|U_{e4}|^2<0.005`$ comes from the BUGEY experiment $`^\mathrm{?}`$ and is still consistent with the conversion into sterile neutrinos in supernovae $`^\mathrm{?}`$.
### 2.3 Neutrinoless Double Beta Decay
The current data indicates
$$_{ee}=\underset{i}{}m_i|U_{ei}|^2<0.5\mathrm{eV}.$$
(23)
In our scheme $`_{ee}=_{i3}m_i|U_{ei}|^2`$. Thus for three flavors $`_{ee}`$ depends only on $`m_1`$ and $`m_2`$, not on $`m_3`$.It is possible to enforce $`_{ee}0`$ for bi-maximal mixing $`^{\mathrm{?},\mathrm{?}}`$. When sterile neutrinos are included this puts a limit on $`m_4`$. One should emphasize that the uncertainties of the nuclear matrix elements could be rather large so it may not be necessary to impose $`_{ee}0`$.
## 3 Conclusions
We explored the implications of imposing the constraint that two neutrino flavors (which for definiteness we take to be $`\nu _\mu `$ and $`\nu _\tau `$) are similarly coupled to the mass basis in addition to the unitarity constraints. We allow three active and an arbitrary number of sterile neutrinos. We show that in this scheme one of the mass eigenstates decouples from the problem, reducing the dimension of the flavor space by one. This result allows significant simplification in the treatment of matter-enhanced neutrino transformation where multiple flavors and level crossings are involved.
When the constraint of Eq. (1) is imposed, which was motivated by the recent experimental results at Superkamiokande, the form of Eq. (16) indicates the existence of a coset structure of the neutrino mixing matrix. Recent related work discussed the existence of a an $`Sp(4)`$ symmetry in the neutrino mass sector $`^\mathrm{?}`$. It was shown that the most general neutrino mass Hamiltonian sits in the $`Sp(4)/SU(2)\times U(1)`$ coset space where U(1) is the chirality transformation and the $`SU(2)`$ generates the see-saw transformation. At the moment it is not clear what the relation, if any, between these two coset structures is.
## Acknowledgments
I thank G. Fuller and T. Weiler for discussions. This work was supported in part by the U.S. National Science Foundation Grant No. PHY-9605140 at the University of Wisconsin, in part by the University of Wisconsin Research Committee with funds granted by the Wisconsin Alumni Research Foundation, and in part by the Alexander von Humboldt-Stiftung, Germany. The very kind hospitality of Hans Weidenmüller at the Max-Planck-Institut für Kernphysik is much appreciated.
## References
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# Parity-violating electromagnetic interactions in QED₃ at finite temperature
## 1 Introduction:
It is well known by now that, in odd space-time dimensions, one can add a topological term to the Lagrangian density of a gauge field, in addition to the usual Maxwell term. Such a term is known as the Chern-Simons term and a theory with such a term is conventionally called a Chern-Simons theory . In $`2+1`$ dimensions, for example, the Chern-Simons (CS) action has the form
$$S_{\mathrm{CS}}=Md^3x\mathrm{tr}ϵ^{\mu \nu \lambda }A_\mu \left(_\nu A_\lambda +\frac{2g}{3}A_\nu A_\lambda \right).$$
(1)
Here “tr” denotes trace over the matrix indices of the gauge fields, $`g`$ the coupling constant while $`M`$ is an arbitrary parameter with the dimensions of mass.
The CS action has several interesting features. Unlike the standard Maxwell action for the gauge fields, it is a topological action. In a theory with a Maxwell term, the CS action generates a mass for the gauge fields . While it is invariant under small gauge transformations, the CS action, for a non-Abelian theory, is not invariant under topologically nontrivial large gauge transformations. Rather, its change is proportional to the winding number associated with the gauge transformation. Explicitly, under
$$A_\mu U^1A_\mu U+\frac{1}{g}U^1_\mu U,$$
(2)
the CS action transforms as
$$S_{\mathrm{CS}}S_{\mathrm{CS}}+\frac{8\pi ^2M}{g^2}W,$$
(3)
where
$$W=\frac{1}{24\pi ^2}d^3x\mathrm{Tr}ϵ^{\mu \nu \lambda }_\mu UU^1_\nu UU^1_\lambda UU^1$$
(4)
is a topological integer known as the winding number of the gauge transformation. For vanishing winding number, the gauge transformations are called small gauge transformations, while for any nontrivial value of the winding number, they are known as large gauge transformations. The CS action clearly is not invariant under large gauge transformations. However, the path integral and, therefore, the theory is, provided the coefficient of the CS term is quantized as
$$\frac{4\pi M}{g^2}=n,$$
(5)
with $`n`$ an integer.
The CS action, in $`2+1`$ dimensions, is known to violate discrete symmetries like $`P`$ and $`T`$. Furthermore, the mass term for a fermion (in the irreducible two component representation) is also known to violate these symmetries. Therefore, if we have massive fermions interacting with a background non-Abelian gauge field, one expects the radiative corrections due to fermions to generate a CS term in the effective action. In fact, it is known that radiative corrections, at zero temperature, shift the value of the tree level CS coefficient such that (assume, for simplicity, $`m>0`$)
$$MM\frac{g^2N_f}{8\pi },$$
(6)
where $`N_f`$ represents the number of fermion flavors. It is clear now that, even if we start with a consistent theory with tree level quantization given by Eq. (5), the radiative corrections change this coefficient and the effective theory will continue to be invariant under large gauge transformations only for an even number of fermion flavors. An even number of fermion flavors is also required to cancel a global anomaly in such theories and, therefore, we see that, in such a case, once the tree level CS coefficient is quantized, the quantum theory continues to have large gauge invariance at the one loop level. In such a theory, it is also known that there is no higher loop corrections to the CS coefficient at zero temperature so that the full quantum theory continues to be invariant under large gauge transformations.
In contrast, it was observed that, at finite temperature, the one loop radiative corrections due to fermions shift the tree level CS coefficient as (We would see later that this corresponds to a particular limit.)
$$MM\frac{g^2N_f}{8\pi }\mathrm{tanh}\frac{\beta m}{2},$$
(7)
where $`\beta =\frac{1}{T}`$ in units of the Boltzmann constant. This, of course, reduces to Eq. (6) when $`T0`$. However, for any nonzero temperature, this is a continuous function and, therefore, even when the tree level CS coefficient is quantized and the number of fermion flavors is even, it cannot take a discrete value as would be required for large gauge invariance to hold. It would appear, therefore, that large gauge invariance would be violated at finite temperature. On the other hand, this is rather strange since temperature is not expected to affect gauge invariance, small or large.
The possible understanding of this puzzle has led to a lot of interest in this topic and only recently, a mechanism for its resolution has been found within the context of the $`0+1`$ dimensional Abelian CS theory . Basically, the resolution of the puzzle in the $`0+1`$ dimensional model goes as follows. For $`N_f`$ flavors of fermions interacting with an Abelian gauge background, at zero temperature, the radiative corrections due to fermions generate only the CS term (namely, only the one point function). On the other hand, at finite temperature, the effective action due to fermions can be exactly evaluated and has the form
$$\mathrm{\Gamma }_f=iN_f\mathrm{log}\left(\mathrm{cos}\frac{a}{2}+i\mathrm{tanh}\frac{\beta m}{2}\mathrm{sin}\frac{a}{2}\right),$$
(8)
where
$$a=𝑑tA(t),$$
(9)
with $`A(t)`$ representing the gauge field. This shows that, unlike at zero temperature, all possible amplitudes are generated in the effective action at finite temperature. Second, all the terms in the effective action are non-extensive and, while every individual term in the effective action violates large gauge invariance, for an even number of fermion flavors, the full effective action is invariant under
$$aa+2\pi N,$$
(10)
which represents the large gauge transformation in this case.
By now, the $`0+1`$ dimensional models have been studied from various points of view . First of all, since we do not expect to be able to evaluate the effective action in closed form in the $`2+1`$ dimensional case, the $`0+1`$ dimensional theory has been studied exhaustively in the perturbative approach . This gives rise to many interesting features. Similarly, if we were to study the $`2+1`$ dimensional theory perturbatively, a signature of large gauge invariance may lie in the large gauge Ward identity. With this in mind, large gauge Ward identities have been derived for the $`0+1`$ dimensional theories which have quite distinctive features. To better understand whether the non-extensive structure is special to $`0+1`$ dimension, the effective action for a fermion, in $`1+1`$ dimensions, interacting with an Abelian gauge background has also been evaluated at finite temperature and it turns out that the effective action, in this case, is extensive although non-local and non-analytic as would be expected in a thermal background. The analysis of the $`0+1`$ dimensional model has also been generalized to $`2+1`$ dimensional models for a restrictive gauge background . Namely, it has been shown that for a single fermion interacting with an Abelian gauge background of the form $`A_0=A_0(t)`$ and $`\stackrel{}{A}=\stackrel{}{A}(\stackrel{}{x})`$, the effective action has the form
$$\mathrm{\Gamma }^{}=\frac{e}{2\pi }d^2x\mathrm{arctan}\left(\mathrm{tanh}\frac{\beta m}{2}\mathrm{tan}\frac{ea}{2}\right)B,$$
(11)
where the magnetic field is defined to be $`B=ϵ^{ij}_iA_j`$.
It is natural to believe that the effective action in Eq. (11) does not represent the complete effective action of the fermion theory in $`2+1`$ dimensions. In fact, it is quite clear that the gauge background is quite restrictive. And, more importantly, the effective action in Eq. (11) does not exhibit non-locality or non-analyticity as would be expected from a thermal effective action. On the other hand, it does represent an all order calculation, be it for a very specific gauge background. It is, of course, quite clear that an exact evaluation of the effective action in a general gauge background is impossible. The only way to go beyond the CS action is through perturbation theory and possibly through the use of the large gauge Ward identity. With this in mind, we have decided to evaluate the parity violating part of the box diagram for a fermion interacting with an arbitrary Abelian gauge background which may serve as a first step towards understanding the question of the effective action and, therefore, large gauge invariance in the $`2+1`$ dimensional theory. Even the calculation of the simple box diagram turns out to be extremely difficult and we had to make use of symbolic computer programs in the intermediate steps. However, the calculation does bring out some interesting features of the theory. The main results of our analysis were already reported in . In this paper, we describe the details of our calculation.
The paper is organized as follows. In section 2, we compile our notation as well as various identities in $`2+1`$ dimensions, which lead to the fact that all the odd point functions vanish in this theory. (This is really a consequence of $`C`$ invariance.) Consequently, one need to look at only even point functions. In section 3, we exhibit the small gauge invariance of the fermion loop at finite temperature. In section 4, we discuss the choice of a small gauge invariant tensor basis which simplifies the calculations. We obtain the parity violating part of the box diagram at zero temperature as well as the quartic effective action associated with this. We evaluate the finite temperature amplitude in two distinct limits, namely, the long wave and the static limits, which shows that the thermal amplitude is indeed non-analytic. We discuss various features of the result and construct the corresponding effective actions. We show that it is really in the static limit that the question of large gauge invariance comes up. In section 5, we derive a large gauge Ward identity and solve for the leading terms in the static limit, which coincides with the effective action, Eq. (11), obtained in the restrictive gauge background. In section 6, we present a brief conclusion along with future directions.
## 2 Notations and Conventions:
Let us consider a single flavor of fermion interacting with a background Abelian gauge field described by the Lagrangian density
$$=\overline{\psi }\left(\gamma ^\mu (i_\mu eA_\mu )m\right)\psi .$$
(12)
Here, $`e`$ represents the electromagnetic coupling strength and we use a diagonal metric with signatures $`(+,,)`$ as well as assume that $`m>0`$. The spinors are two component complex spinors and the Dirac matrices can be represented in terms of the Pauli matrices $`\stackrel{}{\sigma }`$ as follows
$$\gamma ^0=\sigma _2,\gamma ^1=i\sigma _1,\gamma ^2=i\sigma _3,$$
(13)
so that
$$(\gamma ^0)^{}=\gamma ^0,(\gamma ^1)^{}=\gamma ^1,(\gamma ^2)^{}=\gamma ^2,$$
(14)
and
$$(\gamma ^0)^2=1=(\gamma ^1)^2=(\gamma ^2)^2.$$
(15)
The $`2\times 2`$ gamma matrices satisfy some interesting relations such as
$$\gamma ^\mu \gamma ^\nu =\eta ^{\mu \nu }+iϵ^{\mu \nu \lambda }\gamma _\lambda ,$$
(16)
where $`ϵ^{\mu \nu \lambda }`$ represents the totally anti-symmetric Levi-Civita tensor with $`ϵ^{012}=1`$. Relation (16) shows that, unlike in four dimensions, in $`2+1`$ dimensions, we have
$$\mathrm{Tr}\gamma ^\mu \gamma ^\nu \gamma ^\lambda =2iϵ^{\mu \nu \lambda }.$$
(17)
It is worth noting here that the gamma matrices satisfy the relation
$$\mathrm{Tr}\gamma ^{\mu _1}\gamma ^{\mu _2}\mathrm{}\gamma ^{\mu _{2n+1}}=\mathrm{Tr}\gamma ^{\mu _{2n+1}}\gamma ^{\mu _{2n}}\mathrm{}\gamma ^{\mu _1},$$
(18)
which is quite useful in showing that all the odd point functions, in this theory, vanish which, in turn, is a reflection of charge conjugation invariance of the theory. (A word of caution here, namely, that this holds only in the Abelian theory. The presence of internal symmetry generators invalidates this for non-Abelian theories.) Similarly, for an even number of gamma matrices, we have
$$\mathrm{Tr}\gamma ^{\mu _1}\gamma ^{\mu _2}\mathrm{}\gamma ^{\mu _{2n}}=\mathrm{Tr}\gamma ^{\mu _{2n}}\gamma ^{\mu _{2n1}}\mathrm{}\gamma ^{\mu _1},$$
(19)
which helps simplify the calculation of even point functions. There is one other $`2+1`$ dimensional identity which is quite useful in simplifying the calculations, namely, for any arbitrary vector, $`A^\mu `$, we have
$$A^\mu ϵ^{\nu \lambda \sigma }+A^\nu ϵ^{\lambda \mu \sigma }+A^\lambda ϵ^{\mu \nu \sigma }A^\sigma ϵ^{\mu \nu \lambda }=0,$$
(20)
which is really a statement of the fact that in $`2+1`$ dimensions, we cannot have a fourth rank anti-symmetric tensor.
## 3 Gauge Invariance of the Fermion Loop:
In trying to evaluate the effective action due to the fermions, let us next show that, at finite temperature, the $`n`$-point amplitude generated by the fermion loop is gauge invariant, at least under small gauge transformations.
It is simpler to see the small gauge invariance in the real time formalism where there is a doubling of fields . (We will use the closed time path formalism .) In this case, the propagator acquires a $`2\times 2`$ matrix structure and an $`n`$-point amplitude can have both $`+`$ and $``$ type of external vertices. For simplicity, we will show gauge invariance only for amplitudes containing vertices of $`+`$ type (namely, the original vertices) although everything can be carried over to vertices of other kind.
For the type of amplitude that we are interested in (namely, ones with $`+`$ vertices), we need only one component of the fermion propagator, namely,
$$S_{++}(k)=(k/+m)[\frac{1}{k^2m^2+iϵ}+2i\pi n_F(|k^0|)\delta (k^2m^2)].$$
(21)
Here, $`n_F`$ represents the fermion distribution function. Defining $`q=k^{}k`$, we have
$`S_{++}(k^{})q/S_{++}(k)`$ $`=`$ $`\left[{\displaystyle \frac{1}{k^2m^2+iϵ}}+2i\pi n_F(|k^0|)\delta (k^2m^2)\right]`$ (22)
$`\times (k/^{}+m)((k/^{}m)(k/m))(k/+m)`$
$`\times \left[{\displaystyle \frac{1}{k^2m^2+iϵ}}+2i\pi n_F(|k^0|)\delta (k^2m^2)\right]`$
$`=`$ $`S_{++}(k)S_{++}(k^{}).`$
This relation is identical to the one at zero temperature.
Let us next consider a fermion loop with $`N`$ photons, carrying momenta $`p_1,p_2\mathrm{}p_N`$ and indices $`\mu _1,\mu _2,\mathrm{},\mu _N`$ in an ordered way and define
$$k_r=k+p_1+p_2+\mathrm{}+p_r,$$
where $`k`$ represents the momentum in the loop. Let us next attach an extra photon line with momentum $`q`$ and index $`\mu `$ (see Fig.1) between the photon lines carrying momenta $`p_r`$ and $`p_{r+1}`$ (all the lines are of $`+`$ type). Contracting this diagram with $`q_\mu `$, we obtain (we are going to neglect the coupling constants as well as an overall sign coming from the fermion loop)
$`g_r^{\mu _1\mathrm{}\mu _N}`$ $`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^3}\mathrm{Tr}S_{++}(k)\gamma ^{\mu _1}S_{++}(k_1)\gamma ^{\mu _2}\mathrm{}\gamma ^{\mu _r}}`$ (23)
$`\times S_{++}(k_r)q/S_{++}(k_r+q)\gamma ^{\mu _{r+1}}\mathrm{}S_{++}(k_{N1}+q)\gamma ^{\mu _N}`$
$`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^3}\mathrm{Tr}S_{++}(k)\gamma ^{\mu _1}S_{++}(k_1)\mathrm{}\gamma ^{\mu _r}}`$
$`\times \left(S_{++}(k_r)S_{++}(k_r+q)\right)\gamma ^{\mu _{r+1}}\mathrm{}S_{++}(k_{N1}+q)\gamma ^{\mu _N}`$
$`=`$ $`{\displaystyle }{\displaystyle \frac{d^3k}{(2\pi )^3}}\mathrm{Tr}[S_{++}(k)\gamma ^{\mu _1}\mathrm{}\gamma ^{\mu _r}S_{++}(k_r)\gamma ^{\mu _{r+1}}\mathrm{}S_{++}(k_{N1}+q)\gamma ^{\mu _N}`$
$`S_{++}(k)\gamma ^{\mu _1}\mathrm{}\gamma ^{\mu _r}S_{++}(k_r+q)\gamma ^{\mu _{r+1}}\mathrm{}S_{++}(k_{N1}+q)\gamma ^{\mu _N}].`$
Here, we have used the relation Eq. (22) in the intermediate steps.
If we now sum over all the possible insertions of the photon line with momentum $`q`$, terms cancel pairwise to give
$`{\displaystyle \underset{r=1}{\overset{N}{}}}g_r^{\mu _1\mathrm{}\mu _N}`$ $`=`$ $`{\displaystyle }{\displaystyle \frac{d^3k}{(2\pi )^3}}\mathrm{Tr}[S_{++}(k)\gamma ^{\mu _1}S_{++}(k_1)\mathrm{}S_{++}(k_{N1})\gamma ^{\mu _N}`$ (24)
$`S_{++}(k+q)\gamma ^{\mu _1}S_{++}(k_1+q)\mathrm{}S_{++}(k_{N1}+q)\gamma ^{\mu _N}]`$
$`=`$ $`0.`$
Here, we have shifted the variable of integration in the second term by $`kkq`$ to obtain the final result. Such a shift is, of course, meaningful if the integrand is well behaved. We note here that, at finite temperature, the temperature dependent terms are ultraviolet finite and, therefore, such a shift is allowed if the zero temperature part is well behaved, which we know to be true. This argument parallels the zero temperature argument and shows that the $`n`$-point amplitudes generated by the fermion loop are gauge invariant even at finite temperature. Note that this argument can be easily extended to any number of dimensions.
As we have said, the gauge invariance of amplitudes with mixed vertices can be shown in an analogous manner. We simply note here that with a matrix propagator of the form
$$S(k)=\left(\begin{array}{cc}S_{++}(k)& S_+(k)\\ S_+(k)& S_{}(k)\end{array}\right),$$
(25)
where,
$`S_+(k)`$ $`=`$ $`2i\pi (k/+m)(n_F(|k^0|)\theta (k^0))\delta (k^2m^2)`$
$`S_+(k)`$ $`=`$ $`2i\pi (k/+m)(n_F(|k^0|)\theta (k^0))\delta (k^2m^2)`$
$`S_{}(k)`$ $`=`$ $`(k/+m)[{\displaystyle \frac{1}{k^2m^2iϵ}}+2i\pi n_F(|k^0|)\delta (k^2m^2)],`$ (26)
and a matrix vertex of the form
$$\mathrm{\Gamma }^\mu =\left(\begin{array}{cc}\hfill \gamma ^\mu & \hfill 0\\ \hfill 0& \hfill \gamma ^\mu \end{array}\right),$$
(27)
it is easy to see that the generalization of Eq. (22) takes the form
$$S(k^{})(q_\mu \mathrm{\Gamma }^\mu )S(k)=S(k)S(k^{}).$$
(28)
Furthermore, along with the identities ($`q=k^{}k`$),
$$S_\pm (k^{})q/S_\pm (k)=0,$$
(29)
the gauge invariance of any mixed amplitude follows in a completely analogous manner.
## 4 Box Diagram:
Since all the odd point functions vanish in this theory and the two point function is already known (we will come back to the non-analyticity in the two point function later), the next meaningful amplitude to evaluate is the four point function. Furthermore, we are interested only in the parity violating part of this amplitude. This calculation, of course, is extremely cumbersome. However, as we have seen in the last section, the four point function has to be invariant under small gauge transformations and we expect this to be of help. While small gauge invariance alone predicts uniquely the form of the $`n`$-point amplitude in the $`0+1`$ dimensional theory, it is not so in $`2+1`$ dimensions. For example, we know that the $`n`$-point amplitude in $`0+1`$ dimensions has to be of the form
$$\mathrm{\Pi }_{(n)}=\alpha _n\delta (p_1)\delta (p_2)\mathrm{}\delta (p_{n1}),$$
if small gauge invariance has to hold. However, let us note that in the $`2+1`$ dimensional theory, even at the level of the parity violating four point amplitude, there are several possible structures that are compatible with small gauge invariance. Each of the structures below (and possibly more), for example,
$`\mathrm{\Pi }_{PV}^{\mu \nu \lambda \rho }`$ $``$ $`u^\mu u^\nu u^\lambda ϵ^{\rho \sigma \tau }u_\sigma p_{4,\tau }\delta (up_1)\delta (up_2)\delta (up_3)+perm.`$
$`\mathrm{\Pi }_{PV}^{\mu \nu \lambda \rho }`$ $``$ $`u^\mu u^\nu ϵ^{\lambda \rho \tau }p_{4,\tau }\delta (up_1)\delta (up_2)\delta ^3(p_3+p_4)+perm.`$
$`\mathrm{\Pi }_{PV}^{\mu \nu \lambda \rho }`$ $``$ $`u^\mu u^\nu ϵ^{\lambda \rho \sigma }u_\sigma (u(p_3p_4))\delta (up_1)\delta (up_2)\delta (p_3^{})\delta (p_4^{})+perm.`$
where $`u^\mu `$ denotes the velocity of the heat bath, is compatible with small gauge invariance. This is what makes the calculation hard. However, one can simplify the calculation somewhat by choosing a (small) gauge invariant tensor basis for this amplitude.
### 4.1 The Calculation:
The graphs which contribute to the four photon amplitude are shown in Fig. 2. There are three other contributions obtained by charge conjugation. To evaluate these diagrams, we use the analytically continued imaginary-time thermal perturbation theory . This approach can be formulated so as to express the thermal Greens function in terms of forward scattering amplitudes of an on-shell fermion in an external electromagnetic field, as depicted in Fig. 3. Each of these forward scattering amplitude diagrams is obtained by cutting one of the internal lines of the box diagrams in Fig. 2. This, therefore, generates a total of $`4\times 6=24`$ diagrams, which can be systematically obtained from the graph in Fig. 3, by permutations of the external momenta and polarizations. The contribution of the box diagrams, at finite temperature, can then be written in the form
$$\mathrm{\Pi }^{\mu \nu \lambda \rho }(p_1,p_2,p_3,p_4)=\frac{e^4}{(2\pi )^2}\frac{d^2\stackrel{}{k}}{2\omega _k}(n_F(\omega _k)\frac{1}{2})[\underset{ijkl}{}B_{(ijkl)}^{\mu \nu \lambda \rho }+(kk)].$$
(30)
Here $`\omega _k=\sqrt{k^2+m^2}`$, $`n_F(\omega _k)=(\mathrm{e}^{\omega _k/T}+1)^1`$, and the sum is over the permutations $`(ijkl)`$ of $`(1234)`$. Each $`B`$ has a numerator which involves a trace over the Dirac indices. For example, we have
$$\begin{array}{ccc}& & B_{(1234)}^{\mu \nu \lambda \rho }=\hfill \\ & & \\ & & \frac{\mathrm{tr}\left[(/k+m)\gamma ^\mu (/k+/p_1+m)\gamma ^\nu (/k+/p_{12}+m)\gamma ^\lambda (/k+/p_{123}+m)\gamma ^\rho \right]}{\left(2kp_1+p_1^2\right)\left(2kp_{12}+p_{12}^2\right)\left(2kp_{123}+p_{123}^2\right)}|_{k_0=\omega _k},\hfill \end{array}$$
(31)
where $`p_{12}=p_1+p_2`$, etc. Here, we are only interested in the contributions from the trace in Eq. (31) which contain odd powers of the mass, since these will lead to parity-breaking terms (remember that the fermion mass breaks parity).
Let us first study the zero temperature contribution coming from the box diagram, which is associated with the factor $`1/2`$ in the first bracket of Eq. (30), as $`n_F(\omega _k)`$ vanishes in this limit. The computation can be performed explicitly in the low momentum region, where $`|p_\mu |m`$. The result can then be expressed in terms of a series in powers of $`p/m`$, which begins with the leading contribution
$`\mathrm{\Pi }_{PV,T=0}^{\mu \nu \lambda \rho }`$ $`=`$ $`{\displaystyle \frac{ie^4}{16\pi m^6}}\left[ϵ^{\mu \nu \alpha }p_1^\alpha (p_2)^2+ϵ^{\mu \alpha \beta }p_1^\alpha p_2^\beta p_2^\nu \right]`$ (32)
$`\times `$ $`\left[\eta ^{\lambda \rho }p_3p_4p_3^\rho p_4^\lambda \right]+\mathrm{permutations}.`$
It is interesting to note that this result is consistent with the Coleman-Hill theorem which implies that, in the four point Greens function at zero temperature, the terms of order $`p`$ should be absent. In fact, the above structure shows that the parity-violating contributions, generated by the box diagram at $`T=0`$, begin only with terms of order $`(p/m)^5`$. In the configuration space, the low-energy effective action associated with Eq. (32) can be written in the form
$$\mathrm{\Gamma }_{PV,T=0}^4=\frac{e^4}{64\pi m^6}\mathrm{d}^3xϵ^{\mu \nu \lambda }F_{\mu \nu }\left(^\tau F_{\tau \lambda }\right)F^{\rho \sigma }F_{\rho \sigma },$$
(33)
which is manifestly Lorentz and gauge invariant (small and large).
It is worth pointing out that this is the unique, lowest order (in derivatives) parity violating quartic action that one can construct at zero temperature and can be thought of as the generalization of the result of Karplus and Neuman (to the parity violating amplitude in $`2+1`$ dimensions). Of course, one can naively write down other possible structures, for example, of the form
$$S_{PV,T=0}^4=d^3xϵ^{\mu \nu \lambda }_\nu F_{\lambda \tau }F_{\mu \rho }F^{\rho \sigma }F_\sigma ^\tau .$$
(34)
However, using identities such as in Eq. (20) as well as the Bianchi identity, it is straightforward to show that the two structures in Eqs. (33)-(34) are related by a simple multiplicative constant. One can, of course, also construct structures with three epsilon tensors, but they reduce to one of the two forms above. This shows that the lowest order, parity violating quartic action, at zero temperature, has a unique form given in Eq. (33). In fact, the identities, in $`2+1`$ dimensions, are so restrictive that the general form of the lowest order (in derivatives) parity violating effective action can be determined to have the form
$$\mathrm{\Gamma }_{PV}=\mathrm{\Gamma }_{CS}+\underset{n=1}{}a_nϵ^{\mu \nu \lambda }F_{\mu \nu }(^\tau F_{\tau \lambda })(F^{\rho \sigma }F_{\rho \sigma })^n,$$
(35)
with the coefficient $`a_n`$ to be determined perturbatively ($`a_1`$ is already determined in Eq. (33)).
The evaluation of the temperature dependent part of the box diagram, on the other hand, is extremely cumbersome and, as we have mentioned earlier, we would like to systematize the calculation by first selecting a gauge invariant basis which we do next.
### 4.2 Gauge Invariant Tensor Basis:
Let us next construct a set of gauge invariant tensor basis for the parity violating part of the four point amplitude. We note that the tensors in this basis must be linear in the Levi-Civita tensor (odd number of epsilons are, of course, allowed, but reduce to a single epsilon upon using various identities). Furthermore, the tensor basis should also reflect symmetry under exchange of external photon lines. At finite temperature, in addition to the usual tensor structures, we also have the velocity $`u^\mu `$ of the heat bath and, thus, there are, in general, many such structures that one can construct. However, it is practically impossible to carry out the calculation for a general configuration of momenta. For this reason, we have chosen to work with a special configuration of momenta, namely,
$$p_1=p_2=p_3=p=\frac{1}{3}p_4.$$
(36)
In this special configuration, the number of linearly independent, gauge invariant tensor structures is rather easy to determine. For example, for tensor structures where the Levi-Civita tensor has two free indices, there are only two linearly independent structures possible, namely,
$`T_1^{\mu \nu \lambda \rho }`$ $`=`$ $`ϵ^{\sigma \lambda \rho }p_\sigma \left(\eta ^{\mu \nu }{\displaystyle \frac{p^\mu p^\nu }{p^2}}\right)+perm.`$
$`T_2^{\mu \nu \lambda \rho }`$ $`=`$ $`ϵ^{\sigma \lambda \rho }p_\sigma \left(u^\mu {\displaystyle \frac{pu}{p^2}}p^\mu \right)\left(u^\nu {\displaystyle \frac{pu}{p^2}}p^\nu \right)+perm.`$ (37)
These two independent structures are, in fact, quite easy to understand intuitively. Let us recall that, for the parity conserving part of the two point function, there are two independent tensor structures at finite temperature (there are really three structures with a constraint) and the parity violating part of the self-energy has a unique structure with the epsilon tensor. The two structures above simply arise as products of the parity violating structure with the two independent parity conserving structures.
One can similarly look for tensor structures where two of the indices of the Levi-Civita tensor are contracted. There are again only two linearly independent, gauge invariant tensor structures of this kind that one can construct and they have the forms
$`T_3^{\mu \nu \lambda \rho }`$ $`=`$ $`ϵ^{\sigma \tau \rho }u_\sigma p_\tau \left(u^\lambda {\displaystyle \frac{pu}{p^2}}p^\lambda \right)\left(u^\mu {\displaystyle \frac{pu}{p^2}}p^\mu \right)\left(u^\nu {\displaystyle \frac{pu}{p^2}}p^\nu \right)+perm.`$
$`T_4^{\mu \nu \lambda \rho }`$ $`=`$ $`ϵ^{\sigma \tau \rho }u_\sigma p_\tau \left(u^\lambda {\displaystyle \frac{pu}{p^2}}p^\lambda \right)\left(\eta ^{\mu \nu }{\displaystyle \frac{p^\mu p^\nu }{p^2}}\right)+perm.`$ (38)
The four structures in Eqs. (37)-(38) represent a complete set of linearly independent, gauge invariant basis for the parity violating part of the box diagram in this special momentum configuration. (There are other structures possible, but they are not linearly independent.) Therefore, the parity violating part of the four point amplitude can be written as
$$\mathrm{\Pi }_{(4),PV}^{\mu \nu \lambda \rho }=\underset{i=1}{\overset{4}{}}C_iT_i^{\mu \nu \lambda \rho }.$$
(39)
where the coefficients $`C_i`$ are to be determined from the actual evaluation of the Feynman diagrams. Explicit calculation shows that $`C_3=C_4=0`$, so that the parity violating part of the four point amplitude can be expressed in terms of only the first two structures in Eq. (37).
### 4.3 Non-analyticity:
In evaluating the box amplitude at finite temperature, one faces yet another difficulty. Namely, thermal amplitudes are known to be non-analytic at the origin in the energy-momentum plane , which is understood as resulting from new branch cuts that develop at finite temperature due to the possibility of additional channels of reaction. This also translates to the fact that the temperature dependent effective action has a non-analytic structure . This is not of importance in $`0+1`$ dimension where there is no non-analyticity. But, this becomes quite crucial in higher dimensions. Thus, for example, the parity violating part of the self-energy, in $`2+1`$ dimensions, has the form
$$\mathrm{\Pi }_{PV}^{\mu \nu }(p)=\frac{ime^2}{(2\pi )^2}ϵ^{\mu \nu \lambda }p_\lambda \frac{d^2k}{\omega _k}\mathrm{tanh}\left(\frac{\omega _k}{2T}\right)(\frac{1}{p^2+2kp}+kk).$$
(40)
Although it is not widely appreciated, the integrand in Eq. (40) is non-analytic at $`p^\mu =0`$ and depending on how one evaluates the integral, the result would be different. For example, in the long wave limit (LW), the leading term, at high temperature, of the parity violating part of the self-energy takes the form
$$\mathrm{\Pi }_{PV}^{\mu \nu (LW)}(p^0,\stackrel{}{p}=0)=\frac{ie^2}{8\pi }\frac{m}{T}\mathrm{ln}\left(\frac{m}{T}\right)ϵ^{0\mu \nu }p^0+\mathrm{},$$
(41)
giving rise to a leading quadratic effective action of the form
$$\mathrm{\Gamma }_{CS}^{(LW)}=\frac{e^2m}{16\pi T}\mathrm{ln}\frac{m}{T}d^3xϵ^{0ij}A_iE_j.$$
(42)
In contrast, in the static limit (S), the leading behavior of the parity violating term in the self-energy has the form
$$\mathrm{\Pi }_{PV}^{\mu \nu (S)}(p^0=0,\stackrel{}{p})=\frac{ie^2}{4\pi }\mathrm{tanh}\left(\frac{m}{2T}\right)ϵ^{\mu \nu j}p_j+\mathrm{}$$
(43)
giving rise to a leading quadratic effective action of the form
$$\mathrm{\Gamma }_{CS}^{(S)}=\frac{e^2}{4\pi }\mathrm{tanh}\left(\frac{m}{2T}\right)d^3xA_0B.$$
(44)
There are several things to note from this analysis. First, the form of the effective actions, in the two different limits, are quite different. Their leading temperature dependence is also quite distinct. And, finally, if we think of finite temperature as compactifying the time direction and, thereby, inducing a large gauge transformation, then, the effective action in the long wave limit is invariant under such large gauge transformations while the problem of large gauge invariance really manifests in the static limit. Without going into detail, we would like to point out here that the leading contribution to the CS term vanishes if we approach the origin $`p^\mu =0`$ along a light-like direction. It is also worth pointing out here that the original calculation of the CS term corresponds to evaluating it in the static limit (the question of non-analyticity was not very much appreciated then).
It is clear, therefore, that in evaluating the box diagram, at finite temperature, we expect the amplitude to be non-analytic as well. In fact, as in the case of the self-energy, we are going to evaluate this amplitude only in the long wave and the static limits. Let us first concentrate on the long wave limit. In this limit ($`\stackrel{}{p}=0`$) we have the relation $`p^\mu =(pu)u^\mu `$ and, therefore, it is clear that three of the four basis tensor structures in Eqs. (37)-(38) identically vanish (as we have mentioned earlier, the last two structures do not contribute to the parity violating part of the four point amplitude at all), namely,
$$T_2^{\mu \nu \lambda \rho }=0=T_3^{\mu \nu \lambda \rho }=T_4^{\mu \nu \lambda \rho },$$
and the only non-vanishing basis tensor takes the simple form (with a multiplicative factor taken out)
$$T_1^{\mu \nu \lambda \rho }=ϵ^{\sigma \lambda \rho }u_\sigma (\eta ^{\mu \nu }u^\mu u^\nu )+perm.$$
(45)
In this case, the amplitude can be written as
$$\mathrm{\Pi }_{PV}^{\mu \nu \lambda \rho (LW)}=C_1ϵ^{\sigma \lambda \rho }u_\sigma (\eta ^{\mu \nu }u^\mu u^\nu )+perm.$$
(46)
We note here from the explicit structure in Eq. (46) that in the long wave limit, the amplitude is nontrivial only when all the external indices take spatial values. Furthermore, the coefficients $`C_1`$ which depend on the momenta, temperature etc, are to be evaluated from the Feynman diagram and have the form
$$C_1=\frac{6ie^4mp_0}{\pi }_0^{\mathrm{}}d|\stackrel{}{k}|\frac{|\stackrel{}{k}|}{\omega _k}\mathrm{tanh}\frac{\omega _k}{2T}\frac{(3\omega _k^25m^2+2p_0^2)}{(p_0^2\omega _k^2)(p_0^24\omega _k^2)(9p_0^24\omega _k^2)}.$$
(47)
For $`|p_0|m`$, we can expand this in a series of the form
$$C_1=\frac{ie^4mTp_0}{16}\underset{l=\mathrm{}}{\overset{\mathrm{}}{}}\left\{\left(\frac{5m^2}{\mathrm{\Delta }_l^6}+\frac{3}{\mathrm{\Delta }_l^4}\right)\mathrm{ln}\left(1+\frac{\mathrm{\Delta }_l^2}{m^2}\right)\frac{5}{\mathrm{\Delta }_l^4}\frac{1}{2m^2\mathrm{\Delta }_l^2}\right\},$$
(48)
where $`\mathrm{\Delta }_l(2l+1)\pi T`$. In the high temperature limit, the leading contribution comes from the last term in Eq. (48). Performing the summation over $`l`$, we then obtain that
$$C_1(Tm)=\frac{ie^4p_0}{128}\frac{1}{mT}.$$
(49)
Therefore, we see that the leading contribution, in the long wave limit, comes from an extensive effective action of the form
$$\stackrel{~}{\mathrm{\Gamma }}_{PV}^{4(LW)}=\frac{e^4}{512mT}\mathrm{d}^3xϵ^{0ij}E_i\left(_t^1E_j\right)\left(_t^1E_k\right)\left(_t^1E_k\right),$$
(50)
where $`\stackrel{}{E}`$ denotes the electric field. This action is non-local and manifestly gauge invariant (both under small and large gauge transformations) much like the quadratic effective action in the long wave limit. We would like to note here that, in the long wave limit, we have evaluated the amplitude for arbitrary values of the energies, but have chosen to present the results only for the special configuration of Eq. (36) for simplicity.
Next, let us turn to the discussion of the thermal behavior of the box diagram in the static limit, where $`p_i^0=0`$. In this case, due to the very complicated angular integrations, the calculations are extremely difficult, even when using computer algebra. As a result, we have restricted ourselves, in this calculation, to the special configuration of the external spatial momenta (of Eq. (36)), where $`\stackrel{}{p}_1=\stackrel{}{p}_2=\stackrel{}{p}_3=\stackrel{}{p}=\frac{1}{3}\stackrel{}{p}_4`$. In this case, we note that the tensor structure $`T_2^{\mu \nu \lambda \rho }`$ would give contributions only to $`\mathrm{\Pi }_{PV}^{000i}`$ whereas $`T_1^{\mu \nu \lambda \rho }`$ would give contributions to both $`\mathrm{\Pi }_{PV}^{000i}`$ as well as $`\mathrm{\Pi }_{PV}^{0ijk}`$ (in the rest frame of the heat bath). Thus, we see that, in the static limit, the amplitude can have only an odd number of temporal indices (unlike the long wave limit). Let us parameterize the two nontrivial amplitudes as
$`\mathrm{\Pi }_{PV}^{000i(S)}`$ $`=`$ $`2(C_1+C_2)ϵ^{0ij}{\displaystyle \frac{p_j}{|\stackrel{}{p}|}}={\displaystyle \frac{1}{4|\stackrel{}{p}|^2}}ϵ^{0ij}p_j\mathrm{\Pi }_1(\stackrel{}{p},T)`$ (51)
$`\mathrm{\Pi }_{PV}^{0ijk(S)}`$ $`=`$ $`2C_1ϵ^{0kl}{\displaystyle \frac{p_l}{|\stackrel{}{p}|}}\left({\displaystyle \frac{p^ip^j}{\stackrel{}{p}^2}}+\eta ^{ij}\right)`$ (52)
$`=`$ $`{\displaystyle \frac{1}{12|\stackrel{}{p}|^4}}ϵ^{0kl}p_l\left(|\stackrel{}{p}|^2\eta ^{ij}+p^ip^j\right)\mathrm{\Pi }_2(\stackrel{}{p},T),`$
where $`\mathrm{\Pi }_{1,2}(\stackrel{}{p},T)`$ are rather complicated functions of the momenta and the temperature. However, for small momenta, namely, $`|\stackrel{}{p}|m,T`$, we can expand these in a powers series in the momenta and each term in the series can be evaluated in a straightforward manner. Thus, for example, the leading term in $`\mathrm{\Pi }_1`$, in this domain, can be evaluated to have the form
$$\mathrm{\Pi }_1(\stackrel{}{p},T)=\frac{6ie^4}{4\pi }\left[\mathrm{tanh}\left(\frac{m}{2T}\right)\mathrm{tanh}^3\left(\frac{m}{2T}\right)\right]\frac{|\stackrel{}{p}|^2}{T^2}+O\left(\frac{|\stackrel{}{p}|^4}{m^2T^2}\right).$$
(53)
In the high temperature limit, this behaves as $`\frac{1}{T^3}`$, which is quite different from the leading $`\frac{1}{T}`$ behavior of the result (49) in the long wave limit. Let us note here that $`\mathrm{\Pi }_2`$ can also be evaluated in a similar fashion and has the leading high temperature behavior
$$\mathrm{\Pi }_2(\stackrel{}{p},T)=\frac{17ie^4}{16800\pi }\frac{m^3}{T^7}p^4.$$
(54)
It is interesting that terms with lower powers of momentum in $`\mathrm{\Pi }_2`$ identically vanish. As a result, we see that the leading term in $`\mathrm{\Pi }_{PV}^{0ijk(S)}`$ is strongly suppressed at high temperature compared with $`\mathrm{\Pi }_{PV}^{000i(S)}`$ which, in turn, is suppressed relative to $`\mathrm{\Pi }_{PV}^{ijkl(LW)}`$.
The leading contribution given in Eqs. (51) and (53) can be associated with the effective non-extensive action (remember that $`p_4=3p`$)
$$\stackrel{~}{\mathrm{\Gamma }}_{PV}^{4(S)}=\frac{e^4T}{48\pi }\left[\mathrm{tanh}\left(\frac{m}{2T}\right)\mathrm{tanh}^3\left(\frac{m}{2T}\right)\right]\mathrm{d}^3xa_0^3B,$$
(55)
where we have defined
$$a_0=_0^\beta dtA_0(t,\stackrel{}{x})$$
(56)
and $`B`$ is the magnetic field. This form, which may also hold in the quasi-static limit, is consistent with the result derived from the all-orders effective action noted earlier in the special gauge background (see Eq. (11)). We note here that the effective action that would give rise to the amplitude $`\mathrm{\Pi }_{PV}^{0ijk(S)}`$ in Eqs. (52) and (54) can also be determined in a similar manner, but is highly suppressed at high temperature and, unlike the non-extensive action in Eq. (55), would have an extensive, be it non-local structure characteristic of thermal actions. Let us note here that our calculations have been done in the small momentum approximation (which would correspond to a derivative expansion of the effective action). It is well known that , in such an expansion, it is impossible to pick out delta function structures characteristic of non-extensive actions unless one sums the series, which in the present case is simply impossible. In fact, even the evaluation of the leading term in the small momentum expansion already pushes us to the limit of our computational abilities (we really mean even with the use of computers). Therefore, in isolating delta function structures, we have been guided by our earlier experience from the studies in $`0+1`$ dimension , namely, that if an amplitude has a delta function structure, then, in the small momentum expansion, the amplitude vanishes if the variable has a nonzero value and is nonzero only when the variable assumes a vanishing value. This we have checked explicitly. It remains an open question as to whether one can find a better way of isolating delta function structures from a calculation of the leading term in the small momentum expansion.
To conclude this section, therefore, we have found that the temperature dependent part of the parity violating four point amplitude is non-analytic, much like the self-energy. The effective actions, in general, contain both extensive as well as non-extensive terms. In the long wave limit, the leading term at high temperature goes as $`\frac{1}{T}`$ and the effective action associated with this is extensive. Furthermore, this action is invariant under both small and large gauge transformations, much like the CS action in the long wave limit. In the static limit, the leading term in the effective action is non-extensive and behaves as $`\frac{1}{T^3}`$ at high temperature. Furthermore, while this action is invariant under small gauge transformations, it violates large gauge invariance. Thus, large gauge invariance seems to hold order by order in the long wave limit, while it is the static limit where large gauge invariance appears to be an issue at every order.
## 5 Large Gauge Ward Identity:
To a given order, the quasi-static perturbative contributions are not invariant under large gauge transformations generated by $`ea_0ea_0+2\pi N`$, where $`N`$ is a topological integer. But one can derive, in this case, a Ward identity for large gauge invariance, which relates the amplitudes obtained in perturbation theory. To this end, motivated by the structure of Eq. (55), let us write the corresponding all order effective action in the form
$$\stackrel{~}{\mathrm{\Gamma }}^{(S)}=\frac{eT}{2\pi }\mathrm{d}^3x\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{a})B,$$
(57)
where $`\stackrel{~}{a}=ea_0`$. It has been noted in that in the special background $`A_0=A_0(t)`$ and $`\stackrel{}{A}=\stackrel{}{A}(\stackrel{}{x})`$, $`\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{a})`$ corresponds to the real part of the effective action $`\mathrm{\Gamma }_f`$ in Eq. (8), with the identification $`a\stackrel{~}{a}`$. This action obeys, for a single fermion flavor, the large gauge Ward identity
$$\frac{^2\mathrm{\Gamma }^{(1)}}{\stackrel{~}{a}^2}=i\left[\frac{1}{4}\left(\frac{\mathrm{\Gamma }^{(1)}}{\stackrel{~}{a}}\right)^2\right],$$
(58)
where the one point function has the value
$$\frac{\mathrm{\Gamma }^{(1)}}{\stackrel{~}{a}}|_{\stackrel{~}{a}=0}=\frac{1}{2}\mathrm{tanh}\frac{\beta m}{2}.$$
(59)
In order to derive the large gauge Ward identity satisfied by $`\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{a})=\mathrm{}\left[\mathrm{\Gamma }^{(1)}(\stackrel{~}{a})\right]`$, we write
$$\mathrm{\Gamma }^{(1)}(\stackrel{~}{a})=\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{a})+iI(\stackrel{~}{a}),$$
(60)
where $`I`$ denotes the imaginary part of the action $`\mathrm{\Gamma }^{(1)}`$, and substitute this relation into the nonlinear equation (58). Equating to zero the resulting real and imaginary parts, we obtain the following system of coupled equations
$$\frac{^2\stackrel{~}{\mathrm{\Gamma }}}{\stackrel{~}{a}^2}=2\frac{\stackrel{~}{\mathrm{\Gamma }}}{\stackrel{~}{a}}\frac{I}{\stackrel{~}{a}};\frac{^2I}{\stackrel{~}{a}^2}=\frac{1}{4}+\left(\frac{I}{\stackrel{~}{a}}\right)^2\left(\frac{\stackrel{~}{\mathrm{\Gamma }}}{\stackrel{~}{a}}\right)^2.$$
(61)
We must now eliminate from the first equation $`I/\stackrel{~}{a}`$, so as to express $`^2\stackrel{~}{\mathrm{\Gamma }}/\stackrel{~}{a}^2`$ solely in terms of functionals of $`\stackrel{~}{\mathrm{\Gamma }}`$. After some analysis, it turns out that a consistent solution of the above set of equations requires $`I/\stackrel{~}{a}`$ to have the form
$$\frac{I}{\stackrel{~}{a}}=A\mathrm{sin}(\omega \stackrel{~}{\mathrm{\Gamma }})+B\mathrm{cos}(\omega \stackrel{~}{\mathrm{\Gamma }}),$$
(62)
where the coefficients $`A`$ and $`B`$, as well as the frequency $`\omega `$, must be determined from the boundary conditions. One of these conditions can be read directly from (59) and the fact that $`\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{a})=\mathrm{}\left[\mathrm{\Gamma }^{(1)}(\stackrel{~}{a})\right]`$. The other condition follows from the form (57) of the effective action $`\stackrel{~}{\mathrm{\Gamma }}^{(S)}`$ which, as a consequence of invariance under charge conjugation, is a functional involving only even powers of $`A_\mu `$. Consequently, $`\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{a})`$ must contain only odd powers of $`\stackrel{~}{a}`$ and therefore, in particular, $`^2\stackrel{~}{\mathrm{\Gamma }}/\stackrel{~}{a}^2`$ should vanish at $`\stackrel{~}{a}=0`$. These conditions, together with the set of Eqs. (61), determine uniquely $`A`$, $`B`$ and $`\omega `$ in Eq. (62), so that
$$\frac{I}{\stackrel{~}{a}}=\frac{1}{2\mathrm{sinh}\beta m}\mathrm{sin}(2\stackrel{~}{\mathrm{\Gamma }}).$$
(63)
Using this form, we find from the first relation in Eq. (61) that $`\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{a})`$ satisfies the large gauge Ward identity
$$\frac{^2\stackrel{~}{\mathrm{\Gamma }}}{\stackrel{~}{a}^2}=\frac{1}{\mathrm{sinh}\beta m}\frac{\stackrel{~}{\mathrm{\Gamma }}}{\stackrel{~}{a}}\mathrm{sin}\left(2\stackrel{~}{\mathrm{\Gamma }}\right).$$
(64)
This identity, which reflects the underlying large gauge invariance of the quasi-static $`\mathrm{QED}_3`$ theory, relates higher point Greens functions to lower ones. However, unlike the Ward identity for small gauge transformations, the relation (64) is nonlinear. In some sense, this is expected for large gauge transformations which are topologically nontrivial. The relation (64), in fact, allows us to check for large gauge invariance perturbatively. Note from Eqs. (44), (55) and (57) that
$$\frac{\stackrel{~}{\mathrm{\Gamma }}}{\stackrel{~}{a}}|_{\stackrel{~}{a}=0}=\frac{1}{2}\mathrm{tanh}\frac{\beta m}{2},\frac{^3\stackrel{~}{\mathrm{\Gamma }}}{\stackrel{~}{a}^3}|_{\stackrel{~}{a}=0}=\frac{1}{4}\left(\mathrm{tanh}\frac{\beta m}{2}\mathrm{tanh}^3\frac{\beta m}{2}\right).$$
(65)
The identity in Eq. (64) leads to (remember that $`\stackrel{~}{\mathrm{\Gamma }}`$ is odd in $`\stackrel{~}{a}`$ and hence vanishes for $`\stackrel{~}{a}=0`$),
$$\frac{^3\stackrel{~}{\mathrm{\Gamma }}}{\stackrel{~}{a}^3}|_{\stackrel{~}{a}=0}=\frac{2}{\mathrm{sinh}\beta m}\left(\frac{\stackrel{~}{\mathrm{\Gamma }}}{\stackrel{~}{a}}|_{\stackrel{~}{a}=0}\right)^2.$$
(66)
This can be easily seen to hold from the relations in Eq. (65). In fact, the solution of the Ward identity (64), subject to the above boundary conditions, is given by
$$\stackrel{~}{\mathrm{\Gamma }}(\stackrel{~}{a})=\mathrm{arctan}\left[\mathrm{tanh}\frac{\beta m}{2}\mathrm{tan}\left(\frac{\stackrel{~}{a}}{2}\right)\right].$$
(67)
Note that in this solution, which sums up the leading perturbative effects in this region, the tangent is invariant under the large gauge transformations $`\stackrel{~}{a}\stackrel{~}{a}+2\pi N`$. (Incidentally, large gauge invariance would also require quantization of the magnetic flux, which we do not get into here.) Substituting the form (67) in the expression (57), we obtain for $`\stackrel{~}{\mathrm{\Gamma }}(s)`$ a result which agrees, in the static limit of $`\mathrm{QED}_3`$, with the parity-breaking effective action previously discussed in the literature .
## 6 Conclusion:
In this paper, we have studied the radiatively generated parity violating part of the four point amplitude in a theory of a single fermion interacting with an arbitrary Abelian gauge background in $`2+1`$ dimensions at finite temperature. We have shown that the zero temperature part of the parity violating quartic action is unique and, in fact, so is the structure of the complete parity violating part of the effective action. In evaluating the temperature dependent contribution, we have pointed out various obstacles that one has to face and have systematically shown how one can handle these in a given calculation. Of importance is the non-analyticity of thermal amplitudes as well as of the thermal effective actions. We have discussed this in detail for the CS term (self-energy) as well as for the parity violating part of the four photon amplitude. In particular, we have shown that the behavior of the leading amplitudes and, therefore, the leading effective actions in the long wave and static limits are quite distinct at high temperature. Furthermore, while the leading term in the quartic effective action is extensive (but non-local) in the long wave limit, it is non-extensive in the static limit. We have found that, in the long wave limit, large gauge invariance is manifest order by order. In contrast, it appears to be violated order by order in the static limit.
These results can be understood intuitively from the following heuristic arguments. Note that, in $`2+1`$ dimensions, we can always write
$$A_0(t,\stackrel{}{x})=\frac{1}{\beta }_0^\beta 𝑑t^{}A_0(t^{},\stackrel{}{x})+_t\mathrm{\Omega }(t,\stackrel{}{x}),$$
(68)
$$A_i(t,\stackrel{}{x})=\frac{1}{^2}^jF_{ij}(0,\stackrel{}{x})+\left(_0^t\frac{t}{\beta }_0^\beta \right)dt^{}E_i(t^{},\stackrel{}{x})+_i\mathrm{\Omega }(t,\stackrel{}{x}),$$
(69)
where,
$$\mathrm{\Omega }(t,\stackrel{}{x})=\left(_0^t\frac{t}{\beta }_0^\beta \right)dt^{}A_0(t^{},\stackrel{}{x})\frac{1}{^2}\stackrel{}{}\stackrel{}{A}(0,\stackrel{}{x})$$
(70)
Thus, in a particular gauge, we can think of $`a_0(\stackrel{}{x})`$, $`B(0,\stackrel{}{x})`$ and $`\stackrel{}{E}(t,\stackrel{}{x})`$ as representing the physically meaningful variables. From this, it is clear that, in the long wave limit, the only meaningful variable is the electric field which is both small and large gauge invariant. Consequently, the effective action, in this limit, would be large gauge invariant order by order. In contrast, in the static limit, all of the three variables are meaningful implying that the leading term (in derivatives) would involve an odd number of $`a_0(\stackrel{}{x})`$ and a single $`B(0,\stackrel{}{x})`$. Of course, there can be other terms, but they will be higher order in the number of derivatives. Furthermore, order by order, the leading term would violate large gauge invariance.
We have written down a large gauge Ward identity that the leading order terms of the parity violating effective action in the static limit must satisfy for large gauge invariance to hold. This identity can be solved to obtain the leading, all order parity violating effective action which coincides with the action proposed earlier in a restrictive gauge background. However, it is worth remembering that this does not represent the full effective action – rather, it only represents the leading term of the full parity violating effective action.
This study has been carried out within the context of an Abelian gauge theory as a first step towards understanding the question of large gauge invariance at finite temperature. The main interest is, of course, the study of this issue within the context of a non-Abelian gauge theory, which is work in progress.
This work was supported in part by U.S. Dept. Energy Grant DE-FG 02-91ER40685, NSF-INT-9602559 as well as by CNPq, Brazil.
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# Elementary Theory of Line Broadening and Four-wave Mixing in Nonequilibrium Many-Particle Systems
## 1 Time-Dependent Perturbation Theory
The theory here is based on standard time-dependendent perturbation theory, e.g. as presented by Baym . We review here the basic results given in Chapter 12 of Baym.
We begin by assuming that the Hamiltonian of a quantum mechanical system is given by
$$H=H_0+V,$$
where $`H_0`$ is the single-particle Hamiltonian and $`V`$ is a scattering term which is small compared to $`H_0`$. The system of interest can be anything from free atoms or ions to carriers in a semiconductor or metal. We ignore all relativistic effects, however.
We assume that the system has been prepared at time $`t=0`$ in the quantum mechanical state $`|\psi _0`$. At a later time $`t`$, the state is written as $`|\psi _t`$.
The Schrödinger equation gives the time evolution of the system as
$$i\mathrm{}\frac{}{t}|\psi _t=(H_0+V)|\psi _t$$
In the interaction representation, we define a new state $`|\psi (t)`$ ($`t`$ in parenthesis rather than subscript) given by
$$|\psi (t)=e^{iH_0t/\mathrm{}}|\psi _t$$
and a new operator $`V(t)=e^{iH_0t/\mathrm{}}Ve^{iH_0t/\mathrm{}}`$. In this representation, the Schrödinger equation is rewritten as
$$i\mathrm{}\frac{}{t}|\psi (t)=V(t)|\psi (t),$$
which has the advantage of not depending on $`H_0`$.
Integrating this equation as a perturbation series, one can show that
$`|\psi (t)`$ $`=`$ $`e^{(i/\mathrm{}){\scriptscriptstyle V(t)𝑑t}}|\psi (0)`$
$`=`$ $`\left(1+(1/i\mathrm{}){\displaystyle _0^t}V(t^{})𝑑t^{}+(1/i\mathrm{})^2{\displaystyle _0^t}𝑑t^{}{\displaystyle _0^t^{}}𝑑t^{\prime \prime }V(t^{})V(t^{\prime \prime })+\mathrm{}\right)|\psi (0)`$
## 2 Connection to Time-Independent Perturbation Theory
In general, we are interested in knowing the amount of depletion of the initial state after some time $`t`$. Abbreviating $`|\psi (0)=|0`$, we write
$`0|\psi (t)`$ $`=`$ $`(1+(1/i\mathrm{})0|V|0{\displaystyle _0^t}dt^{}`$ (2)
$`+(1/i\mathrm{})^2{\displaystyle _0^t}𝑑t^{}{\displaystyle _0^t^{}}𝑑t^{\prime \prime }{\displaystyle \underset{m}{}}0|V|mm|V|0e^{(i/\mathrm{})(E_0E_m)t^{}}e^{(i/\mathrm{})(E_mE_0)t^{\prime \prime }}`$
$`+(1/i\mathrm{})^3{\displaystyle _0^t}𝑑t^{}{\displaystyle _0^t^{}}𝑑t^{\prime \prime }{\displaystyle _0^{t^{\prime \prime }}}𝑑t^{\prime \prime \prime }{\displaystyle \underset{m,n}{}}0|V|nn|V|mm|V|0`$
$`\times e^{(i/\mathrm{})(E_0E_n)t^{}}e^{(i/\mathrm{})(E_nE_m)t^{\prime \prime }}e^{(i/\mathrm{})(E_mE_0)t^{\prime \prime \prime }}+\mathrm{}).`$
Here we have inserted a sum over the complete set of states $`|mm|=1`$ which are assumed to be eigenstates of the Hamiltonian $`H_0`$. Then we have
$`0|\psi (t)`$ $`=`$ $`(1+(1/i\mathrm{})t0|V|0+{\displaystyle \frac{1}{2}}(1/i\mathrm{})^2t^2|0|V|0|^2`$ (3)
$`+(1/i\mathrm{}){\displaystyle \underset{m0}{}}{\displaystyle _0^t}𝑑t^{}{\displaystyle \frac{|m|V|0|^2}{E_0E_m+i\eta }}e^{(i/\mathrm{})(E_0E_m)t^{}}\left(e^{(i/\mathrm{})(E_mE_0)t^{}}1\right)`$
$`+{\displaystyle \frac{1}{3!}}(1/i\mathrm{})^3t^3|0|V|0|^3`$
$`+2(1/i\mathrm{})^20|V|0{\displaystyle \underset{m0}{}}{\displaystyle _0^t}𝑑t^{}{\displaystyle _0^t^{}}𝑑t^{\prime \prime }{\displaystyle \frac{|m|V|0|^2}{E_0E_m+i\eta }}e^{(i/\mathrm{})(E_0E_m)t^{}}`$
$`\times (e^{(i/\mathrm{})(E_mE_0)t^{}}1)+\mathrm{}).`$
In going from (2) to (3), we have used some assumptions. The upper and lower bounds of each time integral give terms of the form
$$\left(e^{(i/\mathrm{})(E_mE_0)t^{}}1\right).$$
The first, exponential term cancels the exponential term of the next time integral. The second term, -1, gives a vanishing term in the integration because when the last time integration over $`dt^{}`$ is performed, it gives an integral of the form
$$𝒟(E_m)𝑑E_m\left(e^{i(E_mE_0)t/\mathrm{}}1\right)\frac{|m|V|0|^2}{(E_0E_m)^2}.$$
Assuming $`𝒟(E_m)|m|V|0|^2`$ is continuous and slowly varying with $`E_m`$, this integral vanishes because the leading term is odd, i.e.
$$\left(e^{i(E_mE_0)t/\mathrm{}}1\right)i(E_mE_0)t/\mathrm{}$$
for $`(E_mE_0)0`$, and when $`(E_mE_0)0`$, the $`1/(E_0E_m)^2`$ term in the denominator and the fast oscillation of the $`e^{i(E_mE_0)t/\mathrm{}}`$ term kill this integral. This is the “random-phase approximation” (RPA); essentially, it means that we ignore memory of the past and keep only terms from the upper bound of the time integrals.
Note also that rigorously, to introduce the term $`i\eta `$ in the denominator, which allows us to treat the pole at $`E_m=E_0`$, we must assume that $`V=Ve^{\eta t}`$, where $`\eta 0`$. Then we will have
$$_0^te^{\eta t^{}}𝑑t^{}=\frac{e^{\eta t}1}{\eta }=\frac{1\eta t+\mathrm{}1}{\eta }=t$$
and
$`{\displaystyle _0^t}{\displaystyle _0^t^{}}e^{\eta t^{}}e^{\eta t^{\prime \prime }}𝑑t^{}𝑑t^{\prime \prime }={\displaystyle _0^t}e^{\eta t^{}}{\displaystyle \frac{e^{\eta t^{}}1}{\eta }}𝑑t^{}={\displaystyle \frac{e^{2\eta t}1}{2\eta ^2}}{\displaystyle \frac{e^{\eta t}1}{\eta ^2}}`$
$`={\displaystyle \frac{(12\eta t+2\eta t^2+\mathrm{}1)2(1\eta t+\frac{1}{2}t^2+\mathrm{}1)}{2\eta ^2}}={\displaystyle \frac{1}{2}}t^2,`$
etc., which gives us the same result as if we had simply done the integrals assuming $`V`$ is time-independent and then inserted $`i\eta `$ whereever it is needed to take care of a pole.
The series in (3) has the same form as the series expansion of an exponential. By induction, we can write
$$0|\psi (t)=\mathrm{exp}\left[(i/\mathrm{})\left(0|V|0+\underset{m0}{}\frac{|m|V|0|^2}{E_0E_m+i\eta }+\mathrm{}\right)t\right]$$
(4)
This result (4) is extremely useful. It is not an approximation, but is exact to all orders of $`t`$, within the limits of the RPA, for any time-independent $`V`$. The series inside the exponential is just the time-independent perturbation series for the energy correction due to the particle interactions (see, e.g. Baym , chapter 11), which is typically called “Rayleigh-Schrödinger” perturbation theory.
By the Dirac formula, the second-order term is
$`{\displaystyle \underset{m0}{}}{\displaystyle \frac{|m|V|0|^2}{E_0E_m+i\eta }}`$ $`=`$ $`P\left({\displaystyle \underset{m0}{}}{\displaystyle \frac{|m|V|0|^2}{E_0E_m}}\right)i\pi {\displaystyle \underset{m0}{}}|m|V|0|^2\delta (E_0E_m)`$ (5)
$`=`$ $`\mathrm{\Delta }^{(2)}i\mathrm{\Gamma }^{(2)}`$
where the second term is just the total scattering rate. $`\mathrm{\Delta }^{(1)}=0|V|0`$ is called the mean-field energy, $`\mathrm{\Delta }^{(2)}`$ is called the real self-energy and $`\mathrm{\Gamma }^{(2)}`$ is called the imaginary self-energy. From (4) we have therefore
$$0|\psi (t)=e^{(i/\mathrm{})(\mathrm{\Delta }^{(1)}+\mathrm{\Delta }^{(2)}i\mathrm{\Gamma }^{(2)})t}$$
(6)
The probability of being in the state $`|0`$, given by $`|0|\psi (t)|^2`$, decreases over time as $`e^{(2\mathrm{\Gamma }/\mathrm{})t}`$ due to out-scattering. Thus one can see the reason why the imaginary self-energy is associated with an out-scattering rate.
## 3 Absorption and Emission Line Shape in the Single-Particle Picture
We now imagine that we have a single particle prepared in a given state which can decay by two channels, namely by coupling to an external field (normally a photon field) and by scattering to other internal states. We write
$$H=H_0+V+V^{}$$
where $`V^{}`$ is a smaller perturbation than $`V`$, which gives the interband photon-electron interaction. We will only be concerned about the first-order contributions of this term, i.e. we will not worry about renormalization of the electron states due to electron-photon interaction. We will allow renormalization due to $`V`$, however. Therefore we write
$$|\psi ^{}(t)=|\psi (0)+\frac{1}{i\mathrm{}}_0^t𝑑t^{}V^{}(t^{})|\psi (t^{}),$$
where $`|\psi (t^{})`$ is assumed to include the evolution due to $`V`$ given in (1).
We are interested in the rate of emission of a photon with frequency $`\omega `$ and momentum $`𝐤`$. We assume that the eigenstates of $`H_0`$ are single-particle states, each with a given momentum. Therefore by momentum conservation, only one state with momentum $`𝐤`$, which we will identify as $`|0`$, couples to the photon state $`|\omega ,𝐤`$. We write the rate of emission as the probability of being in the photon state $`|\omega ,𝐤`$ at $`t=\mathrm{}`$, after starting in the initial state $`|0`$, divided by the total time spent in the initial state:
$$\frac{1}{\tau (\omega ,𝐤)}=\frac{\left|\omega ,𝐤|\psi ^{}(\mathrm{})\right|^2}{_0^{\mathrm{}}𝑑t|0|\psi (t)|^2}.$$
By the assumptions above, by momentum conservation we write $`\omega ,𝐤|V^{}|m=V^{}\delta _{0,m}`$. Then we have
$`\left|\omega ,𝐤|\psi ^{}(\mathrm{})\right|^2`$ $`=`$ $`\left|{\displaystyle \frac{1}{i\mathrm{}}}{\displaystyle _0^{\mathrm{}}}𝑑t^{}{\displaystyle \underset{m}{}}\omega ,𝐤|e^{(i/\mathrm{})H_0t^{}}V^{}e^{(i/\mathrm{})H_0t^{}}|mm|\psi (t^{})\right|^2`$
$`=`$ $`\left|{\displaystyle \frac{1}{i\mathrm{}}}{\displaystyle _0^{\mathrm{}}}𝑑t^{}\omega ,𝐤|V^{}|0e^{i\omega t^{}}e^{(i/\mathrm{})E_0t^{}}0|\psi (t^{})\right|^2`$
$`=`$ $`|V^{}|^2\left|{\displaystyle \frac{1}{i\mathrm{}}}{\displaystyle _0^{\mathrm{}}}𝑑te^{i\omega t^{}}e^{(i/\mathrm{})E_0t^{}}e^{(i/\mathrm{})(\mathrm{\Delta }^{(1)}+\mathrm{\Delta }^{(2)}i\mathrm{\Gamma }^{(2)})t^{}}\right|^2`$
$`=`$ $`|V^{}|^2\left|{\displaystyle \frac{1}{\mathrm{}\omega Ei\mathrm{\Gamma }}}\right|^2`$
$`=`$ $`|V^{}|^2{\displaystyle \frac{1}{(\mathrm{}\omega E)^2+\mathrm{\Gamma }^2}}`$
where we have abbreviated $`E=E_0+\mathrm{\Delta }^{(1)}+\mathrm{\Delta }^{(2)}`$ and $`\mathrm{\Gamma }=\mathrm{\Gamma }^{(2)}`$ from (6). The normalization factor is
$`{\displaystyle _0^{\mathrm{}}}𝑑t\left|0|\psi (t)\right|^2`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑t\left|e^{(i/\mathrm{})(\mathrm{\Delta }^{(1)}+\mathrm{\Delta }^{(2)}i\mathrm{\Gamma }^{(2)})t}\right|^2`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑te^{2\mathrm{\Gamma }t/\mathrm{}}`$
$`=`$ $`{\displaystyle \frac{\mathrm{}}{2\mathrm{\Gamma }}}`$
This gives the total rate as
$$\frac{1}{\tau (\omega ,𝐤)}=|V^{}|^2\frac{2\mathrm{\Gamma }/\mathrm{}}{(\mathrm{}\omega E)^2+\mathrm{\Gamma }^2}$$
(7)
which when $`\mathrm{\Gamma }0`$, is
$$\frac{1}{\tau (\omega ,𝐤)}=\frac{2\pi }{\mathrm{}}|V^{}|^2\delta (\mathrm{}\omega E).$$
The interpretation of this result is that while the photon has definite momentum and energy, the particle emitting the photon has definite momentum but indefinite energy. The energy uncertainty comes from the Heisenberg uncertainty relation $`\mathrm{\Delta }E\mathrm{\Delta }t\mathrm{}`$, that is, the shorter the time spent in state $`|0`$, the greater the energy uncertainty.
## 4 Connection to the Boltzmann Equation in Many-Particle Theory
We have so far been vague about the nature of the interaction $`V`$. In the previous section we assumed that we had a set of single-particle states with a time-independent out-scattering term $`V`$. In general, however, the particle of interest will be one of many particles in a system, and $`V`$ will be an interaction term between the particles.
The theory of Sections 1 and 2 is completely general; that is, the state $`|0`$ can be taken as a many-particle state instead of a single-particle state. If it is taken as a many-particle state, however, then the result (6) is not as useful, since it involves the total self-energy of the system, while we are typically interested in the single-particle self-energy.
We define the instantaneous many-particle state using creation and destruction operators $`a_𝐤`$ and $`a_𝐤^{}`$ (Baym , chapter 19) as
$$|0=\underset{𝐤}{}\frac{\left(a_𝐤^{}\right)^{n_𝐤}}{\sqrt{n_𝐤!}}|\text{vac},$$
(8)
where $`n_𝐤`$ gives the instantaneous occupation number of each state (this is called a “Fock” state). The “vacuum” state $`|\text{vac}`$ is the zero-particle state, which in the case of a solid means the ground state of the system. Note that we do not need to assume an equilibrium distribution of particles; we can use an instantaneous nonequilibrium distribution if we have that information.
A typical interaction term is written in terms of the same creation and destruction operators, e.g. a two-body, number-conserving term,
$$V=\frac{1}{2}\underset{𝐤_1,𝐤_2,𝐤_3}{}U_{𝐤_1,𝐤_2,𝐤_3,𝐤_4}a_{𝐤_4}^{}a_{𝐤_3}^{}a_{𝐤_2}a_{𝐤_1}$$
(9)
where the summation is not over $`𝐤_4`$ because it is implicitly assumed that momentum is conserved so that $`𝐤_4=𝐤_1+𝐤_2𝐤_3`$. The interaction energy $`U`$ is assumed to be symmetric on exchange of $`𝐤_1`$ with $`𝐤_3`$ or $`𝐤_2`$ with $`𝐤_4`$.
We are concerned about the evolution of the single-particle state $`𝐤`$, which is to say, the number of particles in state $`𝐤`$ as a function of time. The change in the number of particles in a time $`t`$ is given by
$`dn_𝐤`$ $`=`$ $`\psi _t|n_𝐤|\psi _t0|n_𝐤|0`$ (10)
$`=`$ $`\psi (t)|e^{iH_0t/\mathrm{}}n_𝐤e^{iH_0t/\mathrm{}}|\psi (t)0|n_𝐤|0`$
$`=`$ $`0|e^{(i/\mathrm{}){\scriptscriptstyle V(t)𝑑t}}n_𝐤e^{(i/\mathrm{}){\scriptscriptstyle V(t)𝑑t}}|00|n_𝐤|0`$
$`=`$ $`0|e^{(i/\mathrm{}){\scriptscriptstyle V(t)𝑑t}}[n_𝐤,e^{(i/\mathrm{}){\scriptscriptstyle V(t)𝑑t}}]|0.`$
The operator $`n_𝐤`$ commutes with $`H_0`$, by definition. If it commutes with $`V`$, then there is no change in $`n_𝐤`$ over time. We can resolve the commutator in (10) by using the relations
$`[n_𝐤,a_𝐤^{}]`$ $`=`$ $`a_𝐤\delta _{𝐤,𝐤^{}}`$
$`[n_𝐤,a_𝐤^{}^{}]`$ $`=`$ $`a_𝐤^{}\delta _{𝐤,𝐤^{}}`$ (11)
which are valid, surprisingly, for both boson and fermion creation and destruction operators. For a four-operator term in the interaction (9), we have
$`n_𝐤a_{𝐤_4}^{}a_{𝐤_3}^{}a_{𝐤_2}a_{𝐤_1}`$ $`=`$ $`a_{𝐤_4}^{}a_{𝐤_3}^{}a_{𝐤_2}a_{𝐤_1}\delta _{𝐤,𝐤_4}+a_{𝐤_4}^{}a_{𝐤_3}^{}a_{𝐤_2}a_{𝐤_1}\delta _{𝐤,𝐤_3}`$
$`a_{𝐤_4}^{}a_{𝐤_3}^{}a_{𝐤_2}a_{𝐤_1}\delta _{𝐤,𝐤_2}a_{𝐤_4}^{}a_{𝐤_3}^{}a_{𝐤_2}a_{𝐤_1}\delta _{𝐤,𝐤_1}`$
$`+a_{𝐤_4}^{}a_{𝐤_3}^{}a_{𝐤_2}a_{𝐤_1}n_𝐤`$
Thus
$`[n_𝐤,V]`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{𝐤_1,𝐤_2}{}}\left(U_{0321}a_𝐤^{}a_{𝐤_3}^{}a_{𝐤_2}a_{𝐤_1}+U_{3021}a_{𝐤_\mathrm{𝟑}}^{}a_𝐤^{}a_{𝐤_2}a_{𝐤_1}U_{3201}a_{𝐤_3}^{}a_{𝐤_2}^{}a_𝐤a_{𝐤_1}U_{3210}a_{𝐤_3}^{}a_{𝐤_2}^{}a_{𝐤_1}a_𝐤\right)`$ (12)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{𝐤_1,𝐤_2}{}}(U_D\pm U_E)\left(a_𝐤^{}a_{𝐤_1}^{}a_{𝐤_2}a_{𝐤_3}a_{𝐤_3}^{}a_{𝐤_2}^{}a_{𝐤_1}a_𝐤\right)`$
where $`U_D`$ refers to the direct term and $`U_E`$ to the exchange term, and the + sign is for bosons and the - sign is for fermions. (For hard-sphere scattering, i.e. s-wave scattering, $`U`$ is a constant, which gives a factor of 4 enhancement of the scattering cross section for bosons and is forbidden for fermions.)
We first write out the series expansion,
$`dn_𝐤`$ $`=`$ $`0|\left(1(1/i\mathrm{}){\displaystyle _0^t}V(t^{})𝑑t^{}+\mathrm{}\right)`$
$`\times \left((1/i\mathrm{}){\displaystyle _0^t}𝑑t^{}[n_𝐤,V(t^{})]+(1/i\mathrm{})^2{\displaystyle _0^t}𝑑t^{}{\displaystyle _0^t^{}}𝑑t^{\prime \prime }[n_𝐤,V(t^{})V(t^{\prime \prime })]+\mathrm{}\right)|0`$
Any terms in the right-hand series multiplied by the leading “1” in the left-hand series vanish, since $`0|[n_𝐤,A]|0`$ vanishes for any operator $`A`$. Note that $`V^{}(t)=V(t)`$. The leading-order term is therefore
$`dn_𝐤`$ $`=`$ $`0|(1/\mathrm{}^2)\left({\displaystyle _0^t}𝑑t^{}V(t^{})\right)\left({\displaystyle _0^t}𝑑t^{\prime \prime }[n_𝐤,V(t^{\prime \prime })]\right)|0`$
$`=`$ $`{\displaystyle \underset{m}{}}(1/\mathrm{}^2)\left({\displaystyle _0^t}𝑑t^{}e^{(i/\mathrm{})(E_0E_m)t^{}}\right)\left({\displaystyle _0^t}𝑑t^{\prime \prime }e^{(i/\mathrm{})(E_mE_0)t^{\prime \prime }}\right)0|V|mm|[n_𝐤,V]|0`$
$`=`$ $`{\displaystyle \underset{m}{}}\left({\displaystyle \frac{e^{(i/\mathrm{})(E_0E_m)t}1}{E_0E_m}}\right)\left({\displaystyle \frac{e^{(i/\mathrm{})(E_0E_m)t}1}{E_0E_m}}\right)0|V|mm|[n_𝐤,V]|0`$
where the sum over states $`m`$ is over all possible Fock states. The time-dependent factors are resolved using the identities
$`\underset{t\mathrm{}}{lim}{\displaystyle \frac{\left(e^{ixt}1\right)\left(e^{ixt}1\right)}{x^2}}`$ $`=`$ $`\underset{t\mathrm{}}{lim}{\displaystyle \frac{\mathrm{sin}^2(xt/2)}{x^2}}`$
$`=`$ $`\delta (x)2\pi t`$
We therefore have
$`dn_𝐤`$ $`=`$ $`{\displaystyle \underset{m}{}}0|V|mm|[n_𝐤,V]|0{\displaystyle \frac{2\pi t}{\mathrm{}}}\delta (E_0E_m).`$ (13)
Since the different Fock states are assumed orthonormal, the summation in $`V`$ is eliminated because only the terms which couple $`|0`$ to $`|m`$ survive. Since a destruction operator $`a_𝐤`$ acting to the right on a state with $`n_𝐤`$ particles gives a factor $`\sqrt{n_𝐤}`$, and a creation operator $`a_𝐤^{}`$ gives a factor $`\sqrt{1\pm n_𝐤}`$, where the + sign is for bosons (stimulated emission) and the - sign is for fermions (Pauli exclusion), using (12), (13) therefore becomes
$`{\displaystyle \frac{dn_𝐤}{dt}}`$ $`=`$ $`{\displaystyle \frac{2\pi }{\mathrm{}}}{\displaystyle \underset{𝐤_1,𝐤_2}{}}|U_D\pm U_E|^2\left[n_{𝐤_3}n_{𝐤_2}(1\pm n_{𝐤_1})(1\pm n_𝐤)n_𝐤n_{𝐤_1}(1\pm n_{𝐤_2})(1\pm n_{𝐤_3})\right]`$ (14)
$`\times \delta (E_𝐤+E_{𝐤_1}E_{𝐤_3}E_{𝐤_3})`$
where we have assumed $`t`$ is a very small quantity $`dt`$. This is the quantum Boltzmann equation, which gives the total rate of out-scattering for a state $`𝐤`$. It is the same energy we would have written down if we had simply written the total scattering rate as the sum of the Fermi’s Golden Rule rate,
$$\frac{2\pi }{\mathrm{}}|m|V|0|^2\delta (E_0E_m),$$
over all allowed processes of single-particle states $`|m`$. It is correct to order $`V^2`$, but the limit $`t\mathrm{}`$ in (4) implies that the time step $`dt`$ must be long compared to the oscillation time $`\mathrm{}/(E_0E_m)`$, i.e. we are not concerned about behavior on time scales so short that the equivalent energy uncertainty is comparable to the typical collision energies. This is the random-phase approximation again.
We can make one simplification of the Boltzmann equation by treating $`n_𝐤`$ as a continuous variable, in which $`n_𝐤`$ is equal to its average value in an element of phase space $`d^3k`$, and converting the summation to an integral, to get
$`{\displaystyle \frac{d\overline{n}_𝐤}{dt}}`$ $`=`$ $`{\displaystyle \frac{2\pi }{\mathrm{}}}\left({\displaystyle \frac{L^3}{(2\pi )^3}}\right)^2{\displaystyle d^3k_1d^3k_2|U_D\pm U_E|^2\left[n_{𝐤_3}n_{𝐤_2}(1\pm n_{𝐤_1})(1\pm n_𝐤)n_𝐤n_{𝐤_1}(1\pm n_{𝐤_2})(1\pm n_{𝐤_3})\right]}`$ (15)
$`\times \delta (E_𝐤+E_{𝐤_1}E_{𝐤_3}E_{𝐤_4}).`$
We can make an additional simplification if we assume that the system is isotropic, i.e. the distribution function $`n_𝐤`$ depends only on the magnitude of $`𝐤`$ and not on the direction. Then we can integrate analytically over all the angles to reduce this integral to just a double integral which can then be solved numerically for nonequilibrium isotropic distributions. This has been used to produce predictions for various nonequilibrium systems .
These different rates clarify the difference between $`T_1`$, $`T_2`$ etc. in experiments. The rate (14) gives the rate of depletion of a single quantum state, sometimes called the “dephasing” rate in optics. The time constant for this decay is called “$`T_2`$” in NMR terminology. The rate (15) gives the rate of depletion of states with the same macroscopic characteristics as state $`𝐤`$, and the time constant is called “$`T_1`$” in NMR terminology.
We now have calculated the many-body equivalent of the formula (6) for the decay rate. We can now use this in the calculation of line broadening.
## 5 Line broadening in Many-Particle theory
As in Section 3, we are interested in the rate of emission of a photon with momentum $`𝐤`$ and frequency $`\omega `$ from a given initial state, via a weak particle-photon interaction in the presence of a particle-particle interaction. In the many-particle theory we will define the creation and destruction operators $`c_𝐤^{}`$ and $`c_𝐤`$ for photons; the total many-particle state is a product of Fock states of both photons and the interacting particles. As before, we write $`H=H_0+V+V^{}`$, where $`V^{}`$ is the particle-photon interaction term, which we assume has the form
$$V^{}=\underset{𝐤_1}{}P_{𝐤_1}c_{𝐤_1}^{}a_{𝐤_1}.$$
Although this form is strictly unphysical, it is a good model for radiative transitions. Electrons are, of course, conserved, but if we restrict our attention to only one band, then radiative transitions which cause band-to-band transitions involve the destruction of an electron in one band. On the other hand, excitons are not conserved, so this notation allows us to treat electrons and excitons on equal footing as in the previous section. Also, we drop the hermitian conjugate term because we neglect the possibility of photon absorption; we assume there are no photons at $`t=0`$.
The number of photons emitted in a time $`t`$ is then
$`N_𝐤`$ $`=`$ $`\psi (t)|c_𝐤^{}c_𝐤|\psi (t)`$ (16)
$`=`$ $`0|e^{(i/\mathrm{}){\scriptscriptstyle (V(t)+V^{}(t))𝑑t}}[c_𝐤^{}c_𝐤,e^{(i/\mathrm{}){\scriptscriptstyle (V(t)+V^{}(t))𝑑t}}]|0`$
As in the previous section, we assume $`|0`$ is an exact Fock state which is not an eigenstate of the Hamiltonian. This will complicate the analysis (preventing us from using standard many-body formalism) but allows us to treat the case of a highly nonequilibrium system. As before, we write out the series,
$`N_𝐤`$ $`=`$ $`0|(1(1/i\mathrm{}){\displaystyle _0^t}(V(t^{})+V^{}(t^{}))dt^{}`$
$`+(1/i\mathrm{})^2{\displaystyle _0^t}dt^{}{\displaystyle _0^t^{}}dt^{\prime \prime }(V(t^{\prime \prime })+V^{}(t^{\prime \prime }))(V(t^{})+V^{}(t^{}))\mathrm{})`$
$`\times ((1/i\mathrm{}){\displaystyle _0^t}dt^{}[c_𝐤^{}c_𝐤,(V(t^{})+V^{}(t^{}))]`$
$`+(1/i\mathrm{})^2{\displaystyle _0^t}dt^{}{\displaystyle _0^t^{}}dt^{\prime \prime }[c_𝐤^{}c_𝐤,(V(t^{})+V^{}(t^{}))(V(t^{\prime \prime })+V^{}(t^{\prime \prime }))]+\mathrm{})|0`$
Since $`V^{}`$ is not number-conserving, only terms with products of the pair $`V^{}V^{}`$ survive. We restrict ourselves to only terms which are first order in $`V^{}V^{}`$. The photon number $`c_𝐤^{}c_𝐤`$ commutes with $`V(t)`$, while the relation (11) implies that $`[c_𝐤^{}c_𝐤,V^{}]=P_𝐤c_𝐤^{}a_𝐤`$, which in turn implies that only terms with $`V^{}=P_𝐤^{}a_𝐤^{}c_𝐤`$ survive.
At this point we can make an additional assumption in order to greatly simplify the calculation. This is to assume the “dilute” limit, which is that the occupation number of state $`𝐤`$ in $`|0`$ is not greater than one, and in-scattering to any particular state $`𝐤`$ is negligible compared to out-scattering. One can see that this will be the case because for any out-scattering term $`a_𝐤^{}a_{𝐤_3}^{}a_{𝐤_2}a_{𝐤_\mathrm{𝟏}}`$ for which the four momenta satisfy momentum conservation exactly, the probability of having both of states $`𝐤_1`$ and $`𝐤_2`$ occupied simultaneously is very low. Note that in equilibrium, the average net scattering into and out of any state over a long time equals zero, but in this case the many-body state $`|0`$ has been prepared with $`n_𝐤=1`$ at $`t=0`$, and therefore the out-scattering at that time will greatly exceed the in-scattering.
In terms like $`VV^{}`$ there are combinations like $`a_{𝐤_4}^{}a_{𝐤_3}^{}a_{𝐤_2}a_{𝐤_1}a_𝐤`$. No double occupancy means that terms with $`a_𝐤a_𝐤`$ vanish. On the other hand, terms in $`V`$ for which no $`𝐤_i=𝐤`$ commute with $`a_𝐤`$ and one can show that all such terms vanish. We will use $`V_𝐤(t)`$ to indicate the subset of terms in the summation of $`V(t)`$ which do not commute with $`a_𝐤`$. We then get
$`N_𝐤`$ $`=`$ $`0||P_𝐤|^2((1/i\mathrm{}){\displaystyle _0^t}a_𝐤^{}c_𝐤e^{(i/\mathrm{})(\mathrm{}\omega E_𝐤)t^{}}dt^{}+(1/i\mathrm{})^2{\displaystyle _0^t}dt^{}{\displaystyle _0^t^{}}dt^{\prime \prime }V_𝐤(t^{\prime \prime })a_𝐤^{}c_𝐤e^{(i/\mathrm{})(\mathrm{}\omega E_𝐤)t^{}}`$ (17)
$`(1/i\mathrm{})^3{\displaystyle _0^t}dt^{}{\displaystyle _0^t^{}}dt^{\prime \prime }{\displaystyle _0^{t^{\prime \prime }}}dt^{\prime \prime \prime }V_𝐤(t^{\prime \prime \prime })V_𝐤(t^{\prime \prime })a_𝐤^{}c_𝐤e^{(i/\mathrm{})(\mathrm{}\omega E_𝐤)t^{}}+\mathrm{})`$
$`\times ((1/i\mathrm{}){\displaystyle _0^t}dt^{}c_𝐤^{}a_𝐤e^{(i/\mathrm{})(\mathrm{}\omega E_𝐤)t^{}}+(1/i\mathrm{})^2{\displaystyle _0^t}dt^{}{\displaystyle _0^t^{}}dt^{\prime \prime }c_𝐤^{}a_𝐤e^{(i/\mathrm{})(\mathrm{}\omega E_𝐤)t^{}}V_𝐤(t^{\prime \prime })`$
$`+(1/i\mathrm{})^3{\displaystyle _0^t}dt^{}{\displaystyle _0^t^{}}dt^{\prime \prime }{\displaystyle _0^{t^{\prime \prime }}}dt^{\prime \prime \prime }c_𝐤^{}a_𝐤e^{(i/\mathrm{})(\mathrm{}\omega E_𝐤)t^{}}V_𝐤(t^{\prime \prime })V_𝐤(t^{\prime \prime \prime })+\mathrm{})|0`$
$`=`$ $`|P_𝐤|^20|\left((1/i\mathrm{}){\displaystyle _0^t}e^{(i/\mathrm{})(\mathrm{}\omega E_𝐤)t^{}}𝑑t^{}\right)`$
$`\times (1(1/i\mathrm{}){\displaystyle _0^t^{}}V_𝐤(t^{\prime \prime })dt^{\prime \prime }+(1/i\mathrm{})^2{\displaystyle _0^t^{}}{\displaystyle _0^{t^{\prime \prime }}}V_𝐤(t^{\prime \prime })V_𝐤(t^{\prime \prime \prime })dt^{\prime \prime }dt^{\prime \prime \prime }+\mathrm{})`$
$`\times n_𝐤\left((1/i\mathrm{}){\displaystyle _0^t}e^{(1/i\mathrm{})(\mathrm{}\omega E_𝐤)t^{}}𝑑t^{}\right)`$
$`\times \left(1+(1/i\mathrm{}){\displaystyle _0^t^{}}V_𝐤(t^{\prime \prime })𝑑t^{\prime \prime }+(1/i\mathrm{})^2{\displaystyle _0^t^{}}{\displaystyle _0^{t^{\prime \prime }}}V_𝐤(t^{\prime \prime })V_𝐤(t^{\prime \prime \prime })𝑑t^{\prime \prime }𝑑t^{\prime \prime \prime }+\mathrm{}\right)|0`$
$`=`$ $`{\displaystyle \frac{|P_𝐤|^2}{\mathrm{}^2}}0|\left({\displaystyle _0^t}e^{(i/\mathrm{})(\mathrm{}\omega E_𝐤)t^{}}𝑑t^{}\right)e^{(i/\mathrm{})_0^t^{}V_𝐤(t^{\prime \prime })𝑑t^{\prime \prime }}`$
$`\times n_𝐤\left({\displaystyle _0^t}e^{(1/i\mathrm{})(\mathrm{}\omega E_𝐤)t^{}}𝑑t^{}\right)e^{(i/\mathrm{})_0^t^{}V_𝐤(t^{\prime \prime })𝑑t^{\prime \prime }}|0`$
We now break this into two parts, inserting a complete set of states $`|mm|`$ and separating the term with $`|00|`$ from the rest.
$`N_𝐤`$ $`=`$ $`{\displaystyle \frac{|P_𝐤|^2}{\mathrm{}^2}}[\left({\displaystyle _0^t}e^{(i/\mathrm{})(\mathrm{}\omega E_𝐤)t^{}}dt^{}0|e^{(i/\mathrm{})_0^t^{}V_𝐤(t^{\prime \prime })𝑑t^{\prime \prime }}|0\right)`$ (18)
$`\times n_𝐤\left({\displaystyle _0^t}e^{(1/i\mathrm{})(\mathrm{}\omega E_𝐤)t^{}}𝑑t^{}0|e^{(i/\mathrm{})_0^t^{}V_𝐤(t^{\prime \prime })𝑑t^{\prime \prime }}|0\right)`$
$`+{\displaystyle \underset{m0}{}}(1/\mathrm{}^2)|{\displaystyle _0^t}dt^{}e^{(i/\mathrm{})(\mathrm{}\omega E_𝐤)t^{}}{\displaystyle _0^t^{}}dt^{\prime \prime }0|V_𝐤(t^{\prime \prime })|m|^2+\mathrm{}]`$
The second, cross term vanishes, having terms like
$$\left|_0^t𝑑t^{}e^{i\omega _1t^{}}\frac{e^{i\omega _2t^{}}1}{\omega _2}\right|^2=\left|\frac{e^{i(\omega _1+\omega _2)t}1}{(\omega _1+\omega _2)\omega _2}\frac{e^{i\omega _1t}}{\omega _1\omega _2}\right|^2$$
which vanish in the limit of the random-phase approximation, while the other factors are given by (4) – (6). In the limit $`t\mathrm{}`$, we then have
$`N_𝐤`$ $`=`$ $`{\displaystyle \frac{|P_𝐤|^2}{\mathrm{}^2}}n_𝐤\left|{\displaystyle _0^{\mathrm{}}}𝑑t^{}e^{(i/\mathrm{})(\mathrm{}\omega E_𝐤)t^{}}e^{(i/\mathrm{})(\mathrm{\Delta }^{(1)}+\mathrm{\Delta }^{(2)}i\mathrm{\Gamma }^{(2)})}\right|^2`$ (19)
$`=`$ $`|P_𝐤|^2n_𝐤{\displaystyle \frac{1}{(\mathrm{}\omega E)^2+\mathrm{\Gamma }^2}},`$
where we have used the abbreviated notation for $`E`$ and $`\mathrm{\Gamma }`$ of Section 3, and
$`\mathrm{\Gamma }^{(2)}`$ $`=`$ $`\pi {\displaystyle \underset{m0}{}}|m|V_𝐤|0|^2\delta (E_0E_m)`$ (20)
$`=`$ $`\pi {\displaystyle \underset{m0}{}}\left|m|{\displaystyle \frac{1}{2}}{\displaystyle \underset{𝐤_1,𝐤_2}{}}(U_D\pm U_E)a_{𝐤_3}^{}a_{𝐤_2}^{}a_{𝐤_1}a_𝐤|0\right|^2\delta (E_0E_m)`$
$`=`$ $`\pi {\displaystyle \underset{𝐤_1,𝐤_2}{}}|U_D\pm U_E|^2n_𝐤n_{𝐤_1}(1\pm n_{𝐤_2})(1\pm n_{𝐤_3})\delta (E_𝐤+E_{𝐤_1}E_{𝐤_3}E_{𝐤_3})`$
Note that the integral in (20) is the same as in (14), neglecting the in-flow terms, which are negligible in the dilute limit for a particular state $`𝐤`$, even though the average in-flow to a region of phase space $`d^3k`$ around state $`𝐤`$, found in (15), equals the out-flow from the same region in equilibrium.
We can rewrite (19) as
$`{\displaystyle \frac{dN_𝐤}{dt}}`$ $`=`$ $`|P_𝐤|^2{\displaystyle \frac{dn_𝐤}{dt}}{\displaystyle \frac{1}{(\mathrm{}\omega E)^2+\mathrm{\Gamma }^2}}`$ (21)
$`=`$ $`|P_𝐤|^2{\displaystyle \frac{2\mathrm{\Gamma }/\mathrm{}}{(\mathrm{}\omega E)^2+\mathrm{\Gamma }^2}}`$
in which (14) is used to obtain $`dn_𝐤/dt2\mathrm{\Gamma }/\mathrm{}`$.
This all has been essentially a long justification of a very simple formula, which we already obtained in (7) using the single-particle picture. This calculation shows us the limits of validity of this formula, however. In particular, it assumes the dilute limit. In the gas of a high-density boson gas, the situation is much more complex, as shown by Shi, Verachaka and Griffin for case of a quasiequilibrium boson gas. When state $`𝐤`$ is nondegenerate, however, this calculation shows that we can use the the interaction term $`V^{}=|U_D\pm U_E|\sqrt{n_1}\sqrt{1\pm n_2}\sqrt{1\pm n_3}`$ with generality in the single-particle picture for a particle in interacting with a bath that is in equilibrium or far from equilibrium. This type of interaction potential has been used to predict the Lorentzian line broadening of luminescence from various semiconductors , known as “homogeneous” broadening (as opposed to “inhomogeneous” broadening, due to random fluctuations of the band gap.)
An even simpler method, which is justified by this calculation, is to simply write $`H=H_0^{}+V^{}`$ where $`H_0^{}`$ includes the self-energy correction to the single-particle energy due to scattering by $`V`$. In this case we would write $`e^{(i/\mathrm{})H_0^{}t}|0=e^{(i/\mathrm{})(E_𝐤+\mathrm{\Delta }_𝐤^{(1)}+\mathrm{\Delta }_𝐤^{(2)}i\mathrm{\Gamma }^{(2)})t}`$ and then treat the interaction $`V^{}`$ in first-order perturbation theory.
An important implication of the calculations here is that even in the case when there is substantial line broadening, the energy conserving delta-function in (14) should not be replaced by a Lorentzian. One can show that if one were to do this, total energy conservation of the system would be violated.
## 6 The Connection to $`\chi ^{(3)}`$ in Four-Wave Mixing
We have seen that the imaginary self energy gives the line broadening of optical transitions. Alternatively, as seen above, it also corresponds to the out-scattering rate. This rate can often be measured directly in time-resolved measurements. The two methods are complementary. If the rate is extremely fast, it may be too fast for time-resolved methods, but then it gives signficant line broadening. On the other hand, if the line broadening is too small for the experimental spectral resolution to pick up, it corresponds to slow decay rate in time-resolved measurements.
The nonlinear susciptibility is defined in terms of Maxwell’s wave equation in the presence of a polarization current $`J`$,
$$\frac{^2E}{x^2}=\mu _0ϵ_0\frac{^2E}{t^2}+\mu _0\frac{J}{t}.$$
(22)
The polarization current depends on the electric field in a complex way, which is represented to third order as
$$J=ϵ_0\frac{}{t}\left(\chi E+\chi ^{(2)}E^2+\chi ^{(3)}E^3\right).$$
(23)
To calculate $`\chi ^{(3)}`$ from first principles, then, we need to calculate the current due to the oscillating dipole moment of the medium. If we recall that the J = eV = eP/m, we see that fundamentally we want to calculate the momentum $`P`$ of the oscillator as a function of electric field. The value of $`\chi ^{(3)}`$ can then be identified as the part of this which depends to the third order on the electric field $`E`$.
For simplicity we will assume that the oscillator has two states which couple only weakly to other states. These two states nominally correspond to the valence and conduction electron bands in a semiconductor. We will label the ground state $`|v`$ and the excited state $`|c`$. We will also assume that the light field couples only the valence band to the excited band and has no effect on coupling within bands.
Having justified its use in the previous section, we use the simple approach of replacing each state’s energy with the appropriate renormalized energy including imaginary self-energy due to intraband interactions, subject to the same constraint of no multiple occupancy of states, and treat the optical transitions in first-order perturbation theory. Let the time-dependent electron state be $`|e_t`$. In the interaction representation we use $`|e(t)=e^{iH_0t/\mathrm{}}|e_t`$. Then the oscillator momentum of interest is
$`P(t)`$ $`=`$ $`e_t|P|e_t=e(t)|e^{iH_0t/\mathrm{}}Pe^{iH_0t/\mathrm{}}|e(t)`$
$`=`$ $`e(t)|e^{iH_0t/\mathrm{}}{\displaystyle \underset{i}{}}|ii|P{\displaystyle \underset{j}{}}|jj|e^{iH_0t/\mathrm{}}|e(t)`$
$`=`$ $`e(t)|ce^{i(E_c+i\mathrm{\Gamma }_c)t/\mathrm{}}c|P|vv|e(t)e^{i(E_vi\mathrm{\Gamma }_v)t/\mathrm{}}`$
$`+e(t)|ve^{i(E_v+i\mathrm{\Gamma }_v)t/\mathrm{}}v|P|cc|e(t)e^{i(E_ci\mathrm{\Gamma }_c)t/\mathrm{}}`$
$`=`$ $`2\mathrm{}v|Pce(t)|vc|e(t)e^{i(E_vE_c)t/\mathrm{}}e^{(\mathrm{\Gamma }_v+\mathrm{\Gamma }_c)t/\mathrm{}},`$
where $`v|P|c`$ is the standard oscillator strength of the transition.
The time dependence of the electronic state in response to the electric field is determined by the Schrödinger equation
$$i\mathrm{}\frac{}{t}|e_t=H|e_t$$
(25)
$`H`$ $`=`$ $`H_0+V^{}`$ (26)
$`V^{}`$ $`=`$ $`{\displaystyle \frac{e}{2mc}}PA`$ (27)
$`E`$ $`=`$ $`{\displaystyle \frac{1}{c}}{\displaystyle \frac{A}{t}}.`$ (28)
$`V^{}`$ is the radiation term which connects the valence and conduction band states, and $`H_0`$ is the total Hamiltonian of the electron states. $`A`$ and $`E`$ are assumed to be classical fields. Let $`A(t)=A_0e^{i\omega t}`$. Then
$`E(t)`$ $`=`$ $`(i/c)A_0\omega e^{i\omega t}=E_0e^{i\omega t}`$ (29)
which implies
$$V^{}=\frac{ePE_0}{2m\omega }e^{i\omega t}$$
(30)
In the interaction representation,
$$V^{}(t)=e^{iH_0t/\mathrm{}}\frac{ePE_0}{2m\omega }e^{i\omega t}e^{iH_0t/\mathrm{}}.$$
(31)
Third-order time-dependent perturbation theory then gives
$`|e(t)`$ $`=`$ $`|e(t_0)+{\displaystyle \frac{1}{i\mathrm{}}}{\displaystyle _{t_0}^t}𝑑t^{}V^{}(t^{})|e(t_0)`$
$`+{\displaystyle \frac{1}{(i\mathrm{})^2}}{\displaystyle _{t_0}^t}𝑑t^{}{\displaystyle _{t_0}^t^{}}𝑑t^{\prime \prime }V^{}(t^{})V^{}(t^{\prime \prime })|e(t_0)`$
$`+{\displaystyle \frac{1}{(i\mathrm{})^3}}{\displaystyle _{t_0}^t}𝑑t^{}{\displaystyle _{t_0}^t^{}}𝑑t^{\prime \prime }{\displaystyle _{t_0}^{t^{\prime \prime }}}𝑑t^{\prime \prime \prime }V^{}(t^{})V^{}(t^{\prime \prime })V^{}(t^{\prime \prime \prime })|e(t_0)`$
where $`t_0\mathrm{}`$.
Let $`|e(t_0)=|v`$. Then since we assume that $`V^{}`$ does not couple intraband states, so that only $`c|V^{}|v`$ terms survive, we have
$`c|e(t)`$ $`=`$ $`{\displaystyle \frac{1}{i\mathrm{}}}{\displaystyle _{t_0}^t}𝑑t^{}c|V^{}(t^{})|v`$
$`+{\displaystyle \frac{1}{(i\mathrm{})^3}}{\displaystyle _{t_0}^t}𝑑t^{}{\displaystyle _{t_0}^t^{}}𝑑t^{\prime \prime }{\displaystyle _{t_0}^{t^{\prime \prime }}}𝑑t^{\prime \prime \prime }c|V^{}(t^{})|vv|V^{}(t^{\prime \prime })|cc|V^{}(t^{\prime \prime \prime })|v`$
$`v|e(t)`$ $`=`$ $`1+{\displaystyle \frac{1}{(i\mathrm{})^2}}{\displaystyle _{t_0}^t}𝑑t^{}{\displaystyle _{t_0}^t^{}}𝑑t^{\prime \prime }v|V^{}(t^{})|cc|V^{}(t^{\prime \prime })|v`$
Therefore
$`P(t)`$ $`=`$ $`2\mathrm{}v|P|ce^{i(E_vE_c)t/\mathrm{}}e^{(\mathrm{\Gamma }_v+\mathrm{\Gamma }_c)t/\mathrm{}}[{\displaystyle \frac{1}{i\mathrm{}}}{\displaystyle _{t_0}^t}dt^{}c|V^{}(t^{})|v`$
$`+{\displaystyle \frac{1}{(i\mathrm{})^3}}{\displaystyle _{t_0}^t}𝑑t^{}{\displaystyle _{t_0}^t^{}}𝑑t^{\prime \prime }{\displaystyle _{t_0}^{t^{\prime \prime }}}𝑑t^{\prime \prime \prime }c|V^{}(t^{})|vv|V^{}(t^{\prime \prime })|cc|V^{}(t^{\prime \prime \prime })|v`$
$`+{\displaystyle \frac{1}{(i\mathrm{})^3}}\left({\displaystyle _{t_0}^t}𝑑t^{}c|V^{}(t^{})|v\right)\left({\displaystyle _{t_0}^t}𝑑t^{}{\displaystyle _{t_0}^t^{}}𝑑t^{\prime \prime }c|V^{}(t^{})|vv|V^{}(t^{\prime \prime })|c\right)`$
$`+.𝒪(V^5)]`$
There is no second-order $`\chi ^{(2)}`$ by symmetry. We will concentrate on the first third-order term; the calculation of the second term gives similar results.
Substituting in for $`V^{}(t)`$, and performing the time integrals explicitly, we get
$`P(t)`$ $`=`$ $`2\mathrm{}[|v|P|c|^2e^{i(E_vE_c)t/\mathrm{}}e^{(\mathrm{\Gamma }_v+\mathrm{\Gamma }_c)t/\mathrm{}}{\displaystyle \frac{eE_0}{2m\omega }}{\displaystyle \frac{1}{i\mathrm{}}}{\displaystyle _{t_0}^t}dt^{}e^{i(E_ci\mathrm{\Gamma }_c)t^{}/\mathrm{}}e^{i\omega t^{}}e^{i(E_v+i\mathrm{\Gamma }_v)t^{}/\mathrm{}}`$
$`+|v|P|c|^4e^{i(E_vE_c)t/\mathrm{}}e^{(\mathrm{\Gamma }_v+\mathrm{\Gamma }_c)t/\mathrm{}}\left({\displaystyle \frac{eE_0}{2m\omega }}\right)^3{\displaystyle \frac{1}{(i\mathrm{})^3}}{\displaystyle _{t_0}^t}𝑑t^{}{\displaystyle _{t_0}^t^{}}𝑑t^{\prime \prime }{\displaystyle _{t_0}^{t^{\prime \prime }}}𝑑t^{\prime \prime \prime }`$
$`\times \left(e^{i(E_ci\mathrm{\Gamma }_c)t^{}/\mathrm{}}e^{i\omega t^{}}e^{i(E_v+i\mathrm{\Gamma }_v)t^{}/\mathrm{}}\right)\left(e^{i(E_vi\mathrm{\Gamma }_v)t^{\prime \prime }/\mathrm{}}e^{i\omega t^{\prime \prime }}e^{i(E_c+i\mathrm{\Gamma }_c)t^{\prime \prime }/\mathrm{}}\right)`$
. $`\times \left(e^{i(E_ci\mathrm{\Gamma }_c)t^{\prime \prime \prime }/\mathrm{}}e^{i\omega t^{\prime \prime \prime }}e^{i(E_v+i\mathrm{\Gamma }_v)t^{\prime \prime \prime }/\mathrm{}}\right)]`$
$`=`$ $`2\mathrm{}[|v|P|c|^2e^{i(E_vE_c)t/\mathrm{}}e^{(\mathrm{\Gamma }_v+\mathrm{\Gamma }_c)t/\mathrm{}}{\displaystyle \frac{eE_0}{2m\omega }}{\displaystyle \frac{e^{i(E_c\mathrm{}\omega E_v)t/\mathrm{}}}{(E_cE_v\mathrm{}\omega +i(\mathrm{\Gamma }_c+\mathrm{\Gamma }_v))}}e^{(\mathrm{\Gamma }_v+\mathrm{\Gamma }_c)t/\mathrm{}}.`$
$`+|v|P|c|^4e^{i(E_vE_c)t/\mathrm{}}e^{(\mathrm{\Gamma }_v+\mathrm{\Gamma }_c)t/\mathrm{}}\left({\displaystyle \frac{eE_0}{2m\omega }}\right)^3{\displaystyle \frac{1}{(i\mathrm{})^2}}{\displaystyle _{t_0}^t}𝑑t^{}{\displaystyle _{t_0}^t^{}}𝑑t^{\prime \prime }`$
$`\times \left(e^{i(E_ci\mathrm{\Gamma }_c)t^{}/\mathrm{}}e^{i\omega t^{}}e^{i(E_v+i\mathrm{\Gamma }_v)t^{}/\mathrm{}}\right)\left(e^{i(E_vi\mathrm{\Gamma }_v)t^{\prime \prime }/\mathrm{}}e^{i\omega t^{\prime \prime }}e^{i(E_c+i\mathrm{\Gamma }_c)t^{\prime \prime }/\mathrm{}}\right)`$
. $`\times {\displaystyle \frac{e^{i(E_c\mathrm{}\omega E_v)t^{\prime \prime }/\mathrm{}}}{(E_cE_v\mathrm{}\omega +i(\mathrm{\Gamma }_c+\mathrm{\Gamma }_v))}}e^{(\mathrm{\Gamma }_v+\mathrm{\Gamma }_c)t^{\prime \prime }/\mathrm{}}]`$ (37)
$`=`$ $`2\mathrm{}[|v|P|c|^2{\displaystyle \frac{eE_0e^{i\omega t}}{2m\omega }}{\displaystyle \frac{1}{(E_cE_v\mathrm{}\omega +i(\mathrm{\Gamma }_c+\mathrm{\Gamma }_v))}}.`$
$`+|v|P|c|^4\left({\displaystyle \frac{eE_0e^{i\omega t}}{2m\omega }}\right)^3{\displaystyle \frac{1}{(E_cE_v3\mathrm{}\omega +i3(\mathrm{\Gamma }_c+\mathrm{\Gamma }_v))}}\left({\displaystyle \frac{1}{2\omega }}\right)`$
$`\times .{\displaystyle \frac{1}{(E_cE_v\mathrm{}\omega +i(\mathrm{\Gamma }_c+\mathrm{\Gamma }_v))}}e^{2(\mathrm{\Gamma }_c+\mathrm{\Gamma }_v)t/\mathrm{}}]`$
The first term of the final result is the linear $`\chi `$ which goes into the index of refraction, and the second term is proportional to $`E^3`$, i.e. it gives $`\chi ^{(3)}`$. Here we have deduced only a frequency-tripling nonlinear effect with an ingoing and outgoing resonance, because we assumed only one input frequency. Of course, if we had written $`A(t)=A_1(e^{i\omega _1t}+e^{i\omega _1t})+A_2(e^{i\omega _2t}+e^{i\omega _2t})+A_3(e^{i\omega _3t}+e^{i\omega _3t})`$ instead of $`A(t)=A_0e^{i\omega t}`$ as we did, then we would get all the frequency mixing terms associated with third-order optics, e.g. four-wave mixing (FWM).
As seen in (6), the prefactor $`e^{i(E_vE_c)t/\mathrm{}}e^{(\mathrm{\Gamma }_v+\mathrm{\Gamma }_c)t/\mathrm{}}`$ in each of the terms is exactly canceled out by a term with opposite sign in the exponent from the time integrals, so that the oscillation of the polarization is only at the driving frequency and the third-order harmonics.
Suppose that instead of continuing on forever, the electric field is shut off at $`t=0`$, i.e.
$$A(t)=\{\begin{array}{cc}A_0e^{i\omega t},& t<0\\ 0,& t>0\end{array}$$
(38)
Then for times $`t>0`$, the integrals in (6) are time-independent. The only time dependence in $`P(t)`$ comes from the $`e^{i(E_vE_c)t/\mathrm{}}e^{(\mathrm{\Gamma }_v+\mathrm{\Gamma }_c)t/\mathrm{}}`$ prefactor. The polarization continues to oscillate at frequency $`(E_vE_c)/\mathrm{}`$, which is resonant at both $`\omega `$ and $`3\omega `$ (the sum frequency). If there is no scattering (damping) it will continue ringing forever. The imaginary part of the self energy gives the rate at which the phase-coherent polarization dies. This is why it can be called a “dephasing” rate. There is also the feature, seen in (37), that the polarization grows exponentially in time at twice the decay rate, which implies the counterintuitive prediction that a FWM signal will have a longer rise time if there is less dephasing. This has been seen experimentally .
In the case of delayed FWM, a third, “probe” wave follows two “pump” waves by some time delay. The strength of the fourth wave then gives a measure of the $`T_2`$ time. This is because the first two waves each create a polarization wave in the medium due to the linear term in (6), which decays according to the prefactor $`e^{(\mathrm{\Gamma }_v+\mathrm{\Gamma }_c)t/\mathrm{}}`$. These two waves then become contributions to $`E`$ at later times, which mix with the probe wave via the the third-order term. Since the strengths of the pump polarization waves decay, so will the delayed FWM signal, with exactly the same time constant.
We have seen then that the FWM dephasing rate and the line broadening are both controlled by the same imaginary self-energy, i.e. the $`T_2`$ time. This leads one to expect that the line width and the FWM dephasing rate should be inversely proportional. This is generally true, but can be incorrect when a continuum of states is excited, (e.g. by an ultrafast pulse with considerable Heisenberg energy uncertainty), which leads to interference between different oscillator frequencies, as seen in Ref. .
## 7 Scaling Laws in Recent Experiments
The electron-electron Coulomb scattering process represents a special case of the two-body scattering discussed in Section 4, because the interaction cross section diverges in the case of zero momentum transfer, and this divergence must be removed self-consistently by a screening length which depends on the density and the instantaneous energy distribution of the electrons. A previous publication presented the results of scaling laws for electron-electron scattering based on a Boltzmann integral calculation taking into account the dependence of the self-consistent screening length on the density, using a method valid both for nonequilibrium and equilibrium electron plasmas. The main conclusions of that work are reviewed here.
First, the proper integral for the electron-electron scattering depends on the type of experiment. There are three different integrals which relate to the experiments. The first is the “total” scattering rate, i.e. the imaginary self-energy calculated in (20) above. This integral determines the rate of decay of $`\chi ^{(3)}`$, as shown above. A second integral is given by (15) weighted by the momentum exchanged by the electrons in a given scattering event. This typically controls the electron scattering rate determined in transport measurements. A third integral is given by (15) weighted by the energy exchanged by the electrons in a given scattering event. This determines the evolution of the carrier distribution function, which is typically recorded in time-resolved luminescence experiments. These three rates are approximately the same for hard-sphere scattering but are very different for Coulomb scattering. This means that the “relaxation time approximation,” in which all scattering processes are assumed to be characterized by a single “relaxation time,” is justified in the case of short-range interactions but breaks down completely for long-range Coulomb scattering.
Second, the scaling laws as a function of temperature and density for these various integrals were determined for both two dimensions and three dimensions. At that time, several experiments had been done which gave scaling laws consistent with the results of the calculations; recent experiments also fulfill the predictions of that theory. Table 1 gives a summary of the predicted density dependences and the experiments which have observed them. An important result is that this theory predicts that the dephasing rate is independent of density at low density, in basic agreement with recent experiments (Ref. ; also seen indirectly in Ref. .) The “high-density regime” is defined as the regime in which the classical screening length becomes comparable to the interparticle spacing. A full treatment of the scattering rate in this regime would require accounting for the quantum wavefunctions; however, the proper scaling law for the dephasing can be found by the simple assumption that the screening length is pinned at the interparticle spacing. This assumption also correctly gives the crossover from the high-density to low-density scaling regimes.
A third main conclusion of Ref. was that the screening length of the electrons scales with density in same way even for a highly nonequilibrium distribution. Quantum memory effects are important in determining the exact evolution of polarization, but the scaling laws should remain the same even when quantum memory effects are taken into account.
Recent experiments showing a near-constant dephasing rate at low density have been interpreted in terms of optical phonon scattering . A full study of dephasing as a function of the excition photon energy and the temperature should distinguish between these two interpretations, since Ref. also gave predictions for the temperature dependence of the scaling laws.
## 8 Conclusions
The iostropic quantum Boltzmann equation has been used to produce quantitative predictions for numerous experimental systems far from equilibrium . These calculations show that the integrals which enter into these calculations are the same as the integrals which are used for optical line broadening and four-wave mixing, for the cases of a Fermi or nondegenerate Bose gas. In the case of electron Coulomb scattering, great care must be taken to use the properly weighted integral for different experiments.
Acknowledgements. This work has been supported by the National Science Foundation as part of Early Career award DMR-97-22239. One of the authors (D.S.) is a Cottrell Scholar of the Research Corporation.
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# 1 Introduction
## 1 Introduction
Recently there has been considerable interest in the field theories with large extra spacetime dimensions. In comparison to the standard Kaluza-Klein theory these extra dimensions may be restricted only to the gravity sector of theory while the Standard Model (SM) fields are assumed to be localized on the 4-dimensional spacetime (Antoniadis, Arkani-Hamed, Dimopoulos and Dvali (1998); Antoniadis, Arkani-Hamed, Dimopoulos and Dvali (1998); (Arkani-Hamed Dimopoulos and Dvali (1999)). It is promising scenario from the phenomenological point of view because it shifts the energy scale of unification from $`10^{19}`$GeV o $`10100`$TeV.
The gauge field theory is extended by inclusion of the dilatonic field in such of theories. Such fields appear also in a natural way in Kaluza-Klein theories (Appelquist, Chodos and Freund, (1987)), superstring inspired theories (Witten (1985); Ferrata, Lüst and Teisen (1989)) and in theories based on the noncommutative geometry approach (Chamseddine and Fröhlich (1993) ).
As previous studies have already shown the inclusion of a dilaton in a pure Yang - Mills theory has consequences already at the classical level. In particular the dilaton Yang - Mills theories possess ’particle - like’ solutions with finite energy which are absent in the pure Yang - Mills case. Analogous equations have recently been obtained for the ’t Hooft - Polyakov monopole model coupled to the dilatonic field (Lavreashvili and Maison (1992)).
## 2 The dilatonic gauge field theory
Dilatons appear in the higher dimensional theory after the process of the spacetime compactification. The main idea of the theory with large extra dimensions is that the gravity is realized in the more dimensional spacetime (the bulk) while the matter is confined in the four-dimemsional spacetime (the brane). To be clear and simple, we will consider the six-dimensional gravity. Let us now consider the action integral of Einstein-Yang-Mills-Higgs theory in the six-dimensional spacetime:
$$𝒮=d^6x\sqrt{g_6}L,$$
(1)
where $`g_6=det\left(g_{MN}\right)`$ and $`M=\{\mu ,i\}`$, $`N=\{\nu ,j\}`$ with $`x^M=\{x^\mu ,y^i\},i=1,2`$. The metrical tensor in the six-dimensional spacetime can be written:
$$g_{MN}=\left(\begin{array}{cc}e^{2\xi \left(x\right)/f_0}\overline{g}_{\mu \nu }& 0\\ 0& r_2^2\delta _{ij}e^{+2\xi \left(x\right)/f_0}\end{array}\right),$$
(2)
According to above definition we can write:
$$\sqrt{g_6}=\sqrt{\overline{g}}r_2^2e^{2\xi \left(x\right)/f_0}.$$
(3)
In the equation (2)
$$g_{\mu \nu }=e^{2\xi \left(x\right)/f_0}\overline{g}_{\mu \nu }$$
(4)
represents the four-dimensional metric in the Jordan frame while $`\overline{g}_{\mu \nu }`$ in the Einstein frame. We consider here the Lagrangian of the Einstein-Yang-Mills-Higgs field as follows:
$`L=L_g+L_{YMH}\delta \left(y\right),`$ (5)
$`L_g={\displaystyle \frac{1}{2\kappa _6}}\left(R2\mathrm{\Lambda }\right),`$ (6)
$`L_{YMH}={\displaystyle \frac{1}{4}}F_{\mu \nu }^aF^{a\mu \nu }+{\displaystyle \frac{1}{2}}D_\mu \mathrm{\Phi }^aD^\mu \mathrm{\Phi }^aU\left(\mathrm{\Phi }\right),`$ (7)
where $`\kappa _6`$ is the six-dimensional gravitational coupling. $`L`$, $`L_g`$, $`L_{YMH}`$ \- describe the total Lagrange function, the gravity in six-dimensional spacetime and the Yang-Mills-Higgs field parts on the brane emdeded in six-dimensional space, respectively. In general non-vanishing cosmological constant ($`\mathrm{\Lambda }0`$) is possible. This case leads to the interesting monopole solution (Lugo and Shaposhnik (1999, 2000)). In our paper we shall focus our attention on the $`\mathrm{\Lambda }=0`$ case. All calculations should include metric $`g_{MN}`$ (3), so for example $`D^\mu =g^{\mu \nu }D_\nu `$. Let us compactify the six-dimensional spacetime to the four-dimensional Minkowski one on the torus $`\left(_6_4\times 𝒮^1\times 𝒮^1\right)`$. In this paper we assume that the extra dimensions are compactified to two-dimensional torus with a single radius $`r_2`$. The six-dimensional action may be rewritten as:
$$𝒮=d^4xd^2y\sqrt{g_6}L=d^4x\sqrt{\overline{g}},$$
(8)
where $`d^2y=\left(2\pi r_2\right)^2`$ and $``$ is the effective Lagrange function in four-dimensional spacetime. The six-dimensional gravitational coupling $`\kappa _6=8\pi G_6`$ is convenient to define as
$$G_6^1=\frac{1}{\left(2\pi \right)^2}M^4,$$
where $`M`$ \- is the energy scale of the compactification ($`10100TeV`$). Compactification of the six-dimensional gravity on the torus gives the Lagrangian (8) for the four-dimensional gravity
$$L=\frac{1}{2\kappa }R_{\left(4\right)}$$
(9)
in the Einstein frame, where
$$\frac{1}{\kappa }=\frac{\left(2\pi r_2\right)^2}{\kappa _6}$$
(10)
is the four-dimensional coupling constant or $`\kappa =8\pi G_N=8\pi M_{Pl}^2`$ . From the above equation (10) we get
$$M_{Pl}^2=4\pi M^4r_2^2.$$
(11)
Cosmological consideration (Hall and Smith (1999)) gives the bound $`M100TeV`$ which corresponds $`r_25.1\times 10^5`$ mm from equation (11). Compactification of gravity on the five-dimensional spacetime is rather unpysical ($`r_210km`$), however the nice five-dimensional spacetime compactivitation was propsed recently (Randal and Sundram (1999a,b)). The Planck mass $`M_{Pl}`$ (11) is not longer a fundamental constant, it may change during the Universe evolution (Flanagen Tye and Wasserman (1999)). For four-dimensional Minkowski spacetime ($`\overline{g}_{\mu \nu }=\eta _{\mu \nu }`$)
$$R_{\left(4\right)}=\frac{4}{f_0^2}e^{2\xi \left(x\right)/f_0}\left\{_\mu \xi ^\mu \xi +f_0_\mu ^\mu \xi \right\}.$$
(12)
The last term in the equation (12) can by transform into the first one by differentiating by parts. The parameter $`f_0`$ (or re-scaling of the $`\xi \left(x\right)`$ field) is determined by the Planck mass (at present time) as:
$$f_0=\frac{1}{\sqrt{2\pi }}M_{Pl}\mathrm{4.87\hspace{0.17em}10}^{18}GeV/c^2$$
(13)
to produce the $`1/2`$ term in for the dilaton field in (16). The $`f_0`$ parameter determines the dilaton scale $`f_0`$. At the present time $`f_0`$ is rather high, so the interaction with dilatons can be neglected. However, in the early universe when the Planck mass $`M_{Pl}`$ was smaller (for details see (Flanagen Tye and Wasserman (1999))) also the value of the $`f_0`$ was smaller, too.
As a result of compactification of the six-dimensional Lagrangian we get the Lagrange function for the Yang-Mills-Higgs fields. Fluctuations around the four-dimensional Minkowski
$$\overline{g}_{\mu \nu }=\eta _{\mu \nu }+h_{\mu \nu }(x,y)$$
(14)
will produce interaction with Kaluza-Klein dilatons
$$h_{\mu \nu }(x,y)=\underset{𝐧}{}h_{\mu \nu }^𝐧\left(x\right)e^{i\frac{2\pi n^iy^i}{r_2}}$$
(15)
with the typical mass scale $`M`$ (for $`n^i0`$).
In this paper we shall apply this approach to the simplest $`SO\left(3\right)`$ gauge field theory. The $`SO\left(3\right)`$ gauge field theory has nice monopole solutions (’t Hooft (1976)) which have produced trouble in the cosmology and have been the reason to introduce the idea of inflation. The main idea of this paper is to examine how the monopole solution looks like in the dilatonic gauge field theory inspired by the Kaluza-Klein gravity with the TeV scale.
The dilatonic gauge field theory may described by the Lagrangian (defined by equation (8), in the first approximation when $`\overline{g}_{\mu \nu }=\eta _{\mu \nu }`$ in the Minkowski spacetime) we get
$`={\displaystyle \frac{1}{2}}_\mu \xi \left(x\right)^\mu \xi \left(x\right){\displaystyle \frac{1}{4}}e^{2\xi \left(x\right)/f_0}F_{\mu \nu }^aF^{a\mu \nu }+`$ (16)
$`{\displaystyle \frac{1}{2}}\left(D_\mu \mathrm{\Phi }^a\right)D^\mu \mathrm{\Phi }^ae^{2\xi \left(x\right)/f_0}U\left(\mathrm{\Phi }\right),`$
where for $`SO\left(3\right)`$ theory we have:
$$U\left(\mathrm{\Phi }\right)=\frac{\lambda }{4}\left(\mathrm{\Phi }^a\mathrm{\Phi }^av^2\right)^2$$
(17)
with the $`SO\left(3\right)`$ field strength tensor $`F_{\mu \nu }^a=_\mu W_\nu ^a_\nu W_\mu ^a+gϵ_{abc}W_\mu ^bW_\mu ^c`$. The $`SO\left(3\right)`$ gauge symmetry rotates the Higgs field $`\mathrm{\Phi }^a`$. The covariant derivative is given by $`D_\mu \mathrm{\Phi }^a=_\mu \mathrm{\Phi }^agϵ_{abc}W_\mu ^b\mathrm{\Phi }^c`$. Now we have $`D^\mu =\eta ^{\mu \nu }D_\nu `$. The Higgs potential has degenerate true vacua forming the sphere $`S^2`$ $`\left(\mathrm{\Phi }^a\mathrm{\Phi }^a=v^2\right)`$. The Euler-Lagrange equations for the Lagrangian (9) are scale - invariant:
$$x^\mu x^{}{}_{}{}^{\mu }=e^{\frac{u}{f_0}}x^\mu ,$$
(18)
$$\xi \xi ^{}=\xi +u,$$
(19)
$$\mathrm{\Phi }\mathrm{\Phi }^{}=\mathrm{\Phi },$$
(20)
$$W_\mu ^aW_\mu ^a=e^{\frac{u}{f_0}}W_\mu ^a.$$
(21)
These transformations change the Lagrange function in the following way:
$$^{}=e^{\frac{2u}{f_0}}.$$
(22)
This symmetry can be formulated equivalently as a scaling symmetry on the coordinates, and the dilaton is often denoted as a Goldstone boson for dilatations. The origin of the symmetry of the equations of motion is easy understood from the Kaluza-Klein origin of the action. The scale transformations are equivalent to a rescaling of the internal dimensions.
## 3 The dilatonic monopole.
The monopole scalar field configuration:
$`\mathrm{\Phi }^a`$ $`=`$ $`vh\left(r\right)n^a=vH\left(r\right){\displaystyle \frac{n^a}{gr}},`$ (23)
where
$`n^a`$ $`=`$ $`{\displaystyle \frac{x^a}{r}}`$ (24)
describes the ’hedgehog’ structure $`n^a`$ and scalar spherically symmetric field $`H\left(r\right)`$. The $`SO\left(3\right)`$ gauge field is described by the $`K\left(r\right)`$ field:
$$A_i^a=\epsilon _{aij}\frac{1}{gr}n^j\left(1K\left(r\right)\right).$$
(25)
The dilaton field is described by the $`S\left(x\right)`$ function:
$$\xi \left(x\right)=f_0S\left(x\right)$$
(26)
(if we introduce the dimensionless variable $`x=gvr`$). The field Euler-Lagrange equations generated by Lagrange function (16) for $`H\left(x\right)`$ and $`K\left(x\right)`$ and $`S\left(x\right)`$ are:
$$H^{\prime \prime }\left(x\right)\frac{2}{x^2}H\left(x\right)K\left(x\right)^2+\frac{\epsilon }{x^2}e^{2S\left(r\right)}\left(x^2H^2\left(x\right)\right)H\left(x\right)=0,$$
(27)
$$K^{\prime \prime }\left(x\right)+2K^{}\left(x\right)S^{}\left(x\right)\frac{1}{x^2}e^{2S\left(x\right)}H\left(x\right)^2K\left(x\right)\frac{K\left(x\right)}{x^2}\left(K\left(x\right)^21\right)=0,$$
(28)
$`S^{\prime \prime }\left(x\right)`$ $`+`$ $`{\displaystyle \frac{2}{x}}S^{}\left(x\right)\alpha ^1{\displaystyle \frac{1}{x^4}}e^{2S\left(x\right)}\left\{\left(1K\left(x\right)^2\right)^2+2x^2K^{}\left(x\right)^2\right\}`$
$`+`$ $`{\displaystyle \frac{1}{2}}\epsilon \alpha ^1e^{2S\left(r\right)}\left(1{\displaystyle \frac{H^2\left(x\right)}{x^2}}\right)^2=0.`$
In the dilatonic monopole we have two independent dimensionless constants:
$$\epsilon =\frac{\lambda }{g^2},$$
(30)
$$\alpha =\left(\frac{f_0}{v}\right)^2.$$
(31)
The mass (or the lowest energy) of the monopole in the rest frame is:
$$M_{mon}=\frac{4\pi v}{g}\rho \left(x\right)x^2𝑑x,$$
(32)
with the energy density given by:
$`\rho \left(x\right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\alpha S^{}\left(x\right)^2+{\displaystyle \frac{1}{2x^4}}e^{2S\left(x\right)}\left\{\left(1K\left(x\right)^2\right)^2+2x^2K^{}\left(x\right)^2\right\}`$
$`+`$ $`{\displaystyle \frac{1}{2x^4}}\left\{2H\left(x\right)^2K\left(x\right)^2+\left(xH^{}\left(x\right)H\left(x\right)\right)^2\right\}`$
$`+`$ $`{\displaystyle \frac{1}{4}}\epsilon e^{2S\left(x\right)}\left(1{\displaystyle \frac{H^2\left(x\right)}{x^2}}\right)^2.`$
Inside the monopole (according to the equations (27), (28) and (3)) the asymptotical behaviour when $`x0`$ is given as:
$$K\left(x\right)=1tx^2+O\left(3\right),$$
(34)
$$H\left(x\right)=ux^2+O\left(3\right),$$
(35)
$$S\left(x\right)=a+bx^2+O\left(3\right),$$
(36)
where u, t, a are local parameters and $`b`$ must be determined as:
$$b=\frac{24e^{2a}t^2ϵe^{2a}}{12\alpha }.$$
(37)
Far from the monopole core, if $`r\mathrm{}\left(x\mathrm{}\right)`$, both functions H(x) and K(x) should describe the normal vacuum ($`\mathrm{\Phi }^a\mathrm{\Phi }^a=v^2`$) with $`Hx`$ and $`K0`$ according to the $`H,K`$ function definition (23,25) remembering that $`x=gvr`$. In this limit the energy density (3) has a simple limit
$$\rho \left(x\right)=\frac{1}{2}\alpha S^{}\left(x\right)^2+\frac{1}{2x^4}e^{2S\left(x\right)}.$$
The monopole mass in this limit may be rewritten as:
$`M_{mon}={\displaystyle \frac{4\pi v}{g}}{\displaystyle \rho \left(x\right)x^2𝑑x}=`$
$`{\displaystyle \frac{4\pi v}{g}}{\displaystyle \left(\sqrt{\alpha }xS^{}+\frac{1}{x}e^S\right)^2𝑑x}+{\displaystyle \frac{4\pi v}{g}}\sqrt{\alpha }\left(e^{S\left(0\right)}e^{S\left(\mathrm{}\right)}\right).`$
The first term vanishes if the dilaton field obeys the Bogomolny equation
$$\sqrt{\alpha }xS^{}+\frac{1}{x}e^S=0.$$
This equation has a nice solution in the uniform normal vacuum (Bizon (1993))
$$S_b=ln\left(\left(\sqrt{\alpha }e^{S\left(\mathrm{}\right)}1/x\right)/\sqrt{\alpha }\right).$$
(38)
When $`r\mathrm{}`$ dilaton field should disappear in the true vacuum. This demand gives $`S\left(\mathrm{}\right)=0`$. So, When $`r\mathrm{}`$ we have the asymptotic behaviour of the solutions
$$H\left(x\right)=x\left(1we^{\sqrt{2\epsilon }x}\right)+O\left(1/x\right),$$
(39)
$$K\left(x\right)=ze^x+O\left(1/x\right),$$
(40)
$$S\left(x\right)=ln\left(\sqrt{\alpha }e^{S\left(\mathrm{}\right)}\frac{1}{\sqrt{\alpha }x}\right)+O\left(1/x\right).$$
(41)
Even when $`S\left(\mathrm{}\right)0`$ it may be removed by the dilaton transformation (19). We may solve the differential equations (27283) by the iteration method expanding them with respect to $`\epsilon `$:
$$H\left(x\right)=\underset{n=0}{\overset{\mathrm{}}{}}\epsilon ^nH_n\left(x\right),$$
(42)
$$K\left(x\right)=\underset{n=0}{\overset{\mathrm{}}{}}\epsilon ^nK_n\left(x\right),$$
(43)
$$S\left(x\right)=\alpha ^1\underset{n=0}{\overset{\mathrm{}}{}}\epsilon ^nS_n\left(x\right).$$
(44)
In the first step $`\left(n=0\right)`$ we obtain the equations:
$$H_0^{\prime \prime }\left(x\right)\frac{2}{x^2}H_0\left(x\right)K_0\left(x\right)^2=0,$$
(45)
$`K_0^{\prime \prime }\left(x\right)+{\displaystyle \frac{2}{\alpha }}K_0^{^{}}\left(x\right)S_0^{^{}}\left(x\right){\displaystyle \frac{1}{x^2}}e^{2S_0\left(x\right)/\alpha }H_0\left(x\right)^2K_0\left(x\right)`$ (46)
$`{\displaystyle \frac{K_0\left(x\right)}{x^2}}\left(K_0\left(x\right)^21\right)=0,`$
$$S_0\left(x\right)+\frac{2}{x}S_0^{}\left(x\right)\frac{e^{2S_0\left(x\right)/\alpha }}{x^4}\{(1K_0\left(x\right)^2)^2+2x^2K_0^{}\left(x\right)^2\}=0,$$
(47)
leading to the Prasad-Sommerfield solution (without the dilaton field $`S\left(x\right)`$). When $`\alpha \mathrm{}`$ we have the Prasad-Sommerfield solution (Prasad and Sommerfield (1975)) and we can find for $`H_0\left(x\right)`$ and $`K_0\left(x\right)`$ easily as :
$$H_0\left(x\right)=x/tanh\left(x\right)1,$$
(48)
$$K_0\left(x\right)=x/sinh\left(x\right).$$
(49)
Finding a nice analytical solution for the dilaton field (Fig.3, dotted line)
$`S_0\left(x\right)=a+{\displaystyle \frac{Q_D}{x}}+{\displaystyle \frac{1}{4sinh\left(x\right)^2x^2}}(1+2x^2+2xQ_D+cosh\left(2x\right)`$ (50)
$`2xQ_Dcosh\left(2x\right)2xsinh\left(2x\right))`$
was a crucial point of this paper. The leading term for the dilaton field at infinity will be a Coulomb one
$$S_0\left(x\right)=\frac{Q_D}{x}+O\left(1/x\right)$$
$`Q_D`$ is the dilatonic charge which originates from the global scale transformation (18-21). The similarity is striking, but we should remember that an electric charge comes from the $`\mathrm{exp}\left(ie\alpha \right)U\left(1\right)`$ gauge symmetry. The global scale transformation (18-21) is generated by the exponetrial transformation $`\mathrm{exp}\left(u/f_0\right)`$. However, the asymptotic bevaviour at $`x0`$ (36) admits $`Q_D=0`$.
We present the numerical solutions of the coupled set of differential equations (27, 28, 3) in the next section.
## 4 Numerical solutions
To solve the monopole equations numerically, we need the starting point (Press, Teukolsky and Vetterlino (1992)). To find the starting conditions we can use the solutions found from the variational procedures or from the Prasad-Sommerfield approximation (48, 49, 50). The trial solutions depending on the variational parameters must be postulated in such a way to fulfill the boundary conditions close to the center (34,35,36) of the monopole and at far outside (39,40,41). We postulate the trial solutions:
$$H_v\left(x\right)=x\frac{\left(ux+x^2\left(1e^{\sqrt{2\epsilon }x}\right)\right)}{\left(1+x^2\right)}ux^2+O\left(3\right),$$
(51)
$$K_v\left(x\right)=\frac{\left(1tx^2+zx^4\right)}{\left(1+x^4e^x\right)}1tx^2+O\left(3\right)$$
(52)
and
$$S_v\left(x\right)=\frac{\left(a+bx^2+Q_Dx^5\right)}{\left(1+x^6\right)}a+bx^2+O\left(3\right),$$
(53)
where $`O\left(3\right)`$ are corrections of the third order. For these functions the monopole mass was calculated and the trial functions with minimal energy was found. For the monopole without dilatons we get the $`M_{mon}`$ mass of monopole (if $`S\left(x\right)=0`$):
$$M_{mon}=15.886\times 10^{15}GeV.$$
The trial functions give the dilaton configuration close to the monopole case without dilatons. The minimal variational configurations for such a dilatonic monopole for $`Q_D=0`$ and $`v=10^{16}GeV`$ are presented on the Table 1. This shows that the mass of the dilatonic monopole dependences on the parameter $`\alpha `$ ($`f_0`$) and reaches the local minimum (for $`\alpha 1`$) lower than for monopole without dilatons.
For the dilatonic monopole the numerical method was independently verified using the Chebyshev polynomial expansion (Michaila (1999)). The monopole solutions for $`H\left(x\right)`$ and $`K\left(x\right)`$ are known very well so our attention is focused on the dilaton solution $`S\left(x\right)`$ especially. The behavior of the $`H\left(x\right)`$ and $`K\left(x\right)`$ solutions determined by the boundary conditions is the same as it is in the presence of the dilaton field.
The Chebyshev method allows to calculate the exact solution of the differential equations for the discrete set of points. The trial function provides the starting data for the numerical solution of the ordinary differential equation (ODE) (the shooting method (Press, Teukolsky and Vetterlino (1992) )) or the Chebyshev functions method. After this preliminary numerical calculation the method based on the Chebyshev polynomial was used.
Clenshaw and Curtis have proposed almost forty years ago an integration method based on the Chebyshev polynomials of the first kind of degree $`j`$,
$$t_j\left(x\right)=cos\left(jarccos\left(x\right)\right).$$
(54)
Since then, these methods have become standard. Since the Chebyshev polynomials are orthogonal and allows to rewrite the function $`f\left(x\right)`$ as
$$f\left(x\right)=\underset{j=0}{}\alpha _jt_j\left(x\right),$$
(55)
where (for $`j=0`$)
$$\alpha _0=\frac{1}{n}\underset{k=1}{\overset{n}{}}f\left(x_k\right)t_0\left(x_k\right),$$
(56)
and (for $`j0`$)
$$\alpha _j=\frac{2}{n}\underset{k=1}{\overset{n}{}}f\left(x_k\right)t_j\left(x_k\right).$$
(57)
The grid of $`n`$ points $`x_k`$ are zeros of the Chebyshev polynomial $`t_j\left(x\right)`$. This decomposition allows us to present the derivative of the function $`f\left(x\right)`$ as
$$f^{}\left(x_i\right)=\underset{k}{}D_{ik}f\left(x_k\right),$$
(58)
where the matrix
$$D_{ik}=\underset{j=0}{\overset{n1}{}}\frac{1}{c_j}t_j\left(x_k\right)t_j^{}\left(x_i\right),$$
(59)
and $`c_0=n`$, $`c_j=n/2`$ (at $`j0`$). This fact transforms the ordinary differentional equation:
$$\frac{d^2f}{dx^2}+p\left(x\right)\frac{df}{dx}+q\left(x\right)=r\left(x\right)$$
(60)
into an appropriate linear equation:
$$\underset{k}{}A_{ik}f\left(x_k\right)=u_i.$$
(61)
So, we can have exact solution for discrete number of points. This method may be also used to nonlinear equation
$$m\frac{d^2f}{dx^2}+p\left(x\right)\frac{df}{dx}+F(f,x)=0.$$
(62)
If we have a starting function $`f_0`$ then we expand $`f`$ around $`f_0`$
$$f\left(x\right)=f_o\left(x\right)+\epsilon \left(x\right),$$
(63)
and approximate the eq.(62) with the eq.(60) and then solve numerically. The solution may be treat now as a starting function for the next iteration, and so on. The iteration may last so long as an arbitrary precision is reached.
The perturbation around $`\epsilon `$ produces the series of differential equations
$$f_{a,b}^{\prime \prime }+\underset{b}{}p_{ab,n}\left(x\right)f_{b,n}^{}\left(x\right)+\underset{b}{}q_{ab,n}\left(x\right)f_{b,n}\left(x\right)=r_{a,n}\left(x\right),$$
(64)
where the vector
$$f_n=\{H_n\left(x\right),K_n\left(x\right),S_n\left(x\right)\}.$$
(65)
For example when $`n=1`$ the first equation ($`a=1`$) corresponds to $`p_{1b}=0`$, $`q_{11}=2K_0^2\left(x\right)/x^2`$, $`q_{12}=2H_0\left(x\right)K_0\left(x\right)/x^2`$, $`r_1=exp\left(2S_0\left(x\right)/\alpha \right)\left(H_0^2\left(x\right)x^2\right)H_0\left(x\right)/x^2`$, and so on.
In the monopole case the starting function are these obtained by the variational method. After expanding around trial functions (51,52,53) we obtain system of the differential equations of the type (60). After that the numerical solution may be obtained on the grid of the $`x_k`$.
The numerical solution for the dilaton field found by the Chebyshev numerical method is presented on the Fig.3 (the solid line).
## 5 Conclusions
The aim of this paper was to present a numerical study of the classical monopole solutions of the SO(3) theory coupled to the dilaton fields. We have shown that a monopole is surrounded by the dilaton cloud $`S\left(x\right)`$. In the field theories with large extra spacetime dimensions the Planck mass is not longer a fundamental constant and may change itself during the evolution of the universe. As a consequence the parameter $`f_0`$ changes, too. We have shown that the dilatonic monopole reaches the minimal mass when $`f_0v`$ with the mass lower a bit than for monopole without dilatons.
There is analytical solution in the Prasad-Sommerfield limit.
The spherically symmetric dilaton solutions coupled to the gauge field or gravity are interesting in their own and may moreover influence the monopole catalysis.
However in the theory inspired by the Kaluza-Klein theory with large extra dimensions also new interaction with massive ($`M`$) Kaluza-Klein gravitons takes place. In the four dimensional spacetime the monopole solutions is stable due to the monopole topological charge. Now the interaction with Kaluza-Klein gravitons $`h_{\mu \nu }^𝐧\left(x\right)`$ may cause disintegration of the monopole.
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# Non-unitarity of CKM matrix from vector singlet quark mixing and neutron electric dipole moment
## Abstract
In the standard model (SM) the lowest order contribution to the quark electric dipole moment (EDM) occurs at the three loop level. We show that the non-unitarity of the CKM matrix in models with an extended quark sector typically gives rise to a quark EDM at the two loop level which has no GIM-like suppression factors except the external quark mass. The induced neutron EDM is of order $`10^{29}`$ $`\mathrm{e}\mathrm{cm}`$ and can be well within the reach of the next generation of experiments if it is further enhanced by long distance physics as happens in the SM.
Keywords: electric dipole moment, CKM matrix, singlet quark
PACS numbers: 11.30.Er 12.15.Ff 12.60.-i 13.40.Em
An important target of particle physics is the determination of the Cabibbo-Kobayashi-Maskawa (CKM) matrix, which parameterizes the charged current interactions of quarks. In the standard model (SM) with three generations of quarks the CKM matrix is a $`3\times 3`$ unitary matrix, and the $`CP`$ violation is due to the presence of a nonzero phase in this matrix. The unitarity of the CKM matrix is essential in suppressing flavour changing neutral current (FCNC) processes. Using the CKM unitarity relation the $`CP`$ violating phase information can be elegantly displayed in terms of unitary triangles. The most interesting challenge of the $`B`$-factories now entering into operation, and of future collider $`B`$ experiments is to try to pin down the angles in the unitary triangles. If the measured angles violate either of the ‘triangle conditions’, or correspond to a point $`(\rho ,\eta )`$ which is outside the allowed region, then we will have evidence for new physics.
Violation of the CKM unitarity can appear in models with an extended quark sector. Various constraints on the possibility that exotic quarks mix with the ordinary SM quarks have been derived from low energy charged and neutral current phenomenology, $`Z`$ physics, FCNC processes and $`CP`$ violation in neutral meson systems. In this letter we propose to examine the unitarity of the CKM matrix in a different setting, namely, by investigating possible unitarity violating effects on the neutron electric dipole moment (EDM). We shall find that information from the neutron EDM is complementary to that from FCNC processes in $`B`$ physics and serves as a self-consistency check of the relevant theory.
In the minimal SM quark EDM’s vanish at the one loop level because the relevant amplitudes do not change the quark flavor and each CKM matrix element is accompanied by its complex conjugate so that no $`T`$-violating complex phase can arise. At the two loop level individual diagrams can have a complex phase, but it has been shown by Shabalin that their sum vanishes strictly. The null result was confirmed afterwards by several groups. It is thus thought that in the SM the lowest order contribution to quark EDM’s occurs at the three loop level. A recent calculation shows they are of order $`10^{35}`$ to $`10^{34}`$ e cm for $`u`$ and $`d`$ quarks. The extreme smallness of the quark and neutron EDM’s in the SM makes them particularly suited for searching for new physics. The current experimental upper bound, $`|d(n)|<6.3\times 10^{26}\mathrm{e}\mathrm{cm}`$, has put very strigent constraints on extensions of the SM, such as additional Higgs fields, right-handed currents, or supersymmetric partners. We reanalyzed the problem in Refs. and and found that the complete two loop cancellation can be attributed to two special features in the SM: the purely left-handed structure of the charged current and the unitarity of the $`3\times 3`$ CKM matrix. The cancellation introduced by the unitarity is the flavour diagonal analog of the GIM suppression in the FCNC processes, with the mere difference being that the cancellation is complete in the EDM case. Thus it is expected that in models with an extended quark sector the contributions to quark EDM’s at the two loop level are no longer cancelled thoroughly because of violation of the CKM unitarity, and that potentially large EDM’s for quarks and the neutron can then be induced. Ignoring possible logarithmic factors, they are of order, $`d(q)eg^4\pi ^4\stackrel{~}{\delta }m_q/m_W^2`$, where $`g`$ is the semi-weak coupling, $`\stackrel{~}{\delta }`$ is the rephasing invariant measure of $`CP`$ violation. Note that there are no GIM-like suppression factors except the external light quark mass which is required by chirality flip. Numerically they are of order $`10^{29}`$ $`\mathrm{e}\mathrm{cm}`$, well within the reach of the next generation of experiments if they are further enhanced by long distance physics as happens in the SM.
We consider here a model with one extra singlet down-type quark in a vector-like representation of the SM gauge group, $`SU(3)_C\times SU(2)_L\times U(1)_Y`$. In addition to the three quark generations, each consisting of the three representations$`(i=1,2,3)`$
$$Q_L^i(3,2)_{+1/6},u_R^i(3,1)_{+2/3},d_R^i(3,1)_{1/3},$$
(1)
we have the following vector-like representation:
$$d_4(3,1)_{1/3}+\overline{d}_4(\overline{3},1)_{+1/3}$$
(2)
Such a quark representation appears, for example, in $`E_6`$ GUTs . The model can be considered as a minimal extension of the SM in the sense that there is no other change in the gauge and scalar sectors. In particular, the charged currents remain purely left-handed and the $`CP`$ violation is still encoded in the CKM matrix.
After spontaneous symmetry breaking, the down-type singlet quark ($`d_4`$) mixes with the ordinary three down-type quarks so that the weak and mass eigenstates are related by a $`4\times 4`$ unitary matrix,
$$\left(\begin{array}{c}d^{}\\ s^{}\\ b^{}\\ d_4^{}\end{array}\right)_L=\left(\begin{array}{cccc}V_{ud}& V_{us}& V_{ub}& V_{u4}\\ V_{cd}& V_{cs}& V_{cb}& V_{c4}\\ V_{td}& V_{ts}& V_{tb}& V_{t4}\\ V_{0d}& V_{0s}& V_{0b}& V_{04}\end{array}\right)\left(\begin{array}{c}d\\ s\\ b\\ d_4\end{array}\right)_L.$$
(3)
In the basis where the up-type quarks are diagonalized, the submatrix consisting of the first three rows in the above matrix appears in $`SU(2)_L`$ $`W`$ and $`Z`$ couplings, and it is the generalized CKM matrix in this model. Note that there still exist unitarity relations among the up-type quarks although the matrix is no longer unitary. The non-unitarity of the matrix leads to FCNC $`Z`$ couplings amongst down-type quarks which have important repercussions on FCNC and $`CP`$ violating processes. In this work we shall focus on the purely charged current contributions to the EDM’s. We shall also discuss briefly the more complicated contributions involving FCNC interactions.
Let us consider the quark EDM. The effective Lagrangian for the EDM interaction is defined as $`_{\mathrm{eff}}=id/2\overline{\psi }\gamma _5\sigma _{\mu \nu }\psi F^{\mu \nu }`$, where $`F^{\mu \nu }`$ is the electromagnetic tensor and $`d`$ is the EDM of the fermion $`\psi `$. Again, no $`T`$-violating complex phase can arise at the one loop level. The contributing Feynman diagram at the two loop level is shown in Fig. 1. Following our previous works, we use the background field- ( or the nonlinear ) $`R_\xi `$ gauge with $`\xi =1`$. In this gauge, there is no $`W^\pm G^{}A`$ coupling and the $`W^+W^{}A`$ coupling is also very simple.\[ $`A`$ is the background electromagnetic field and $`G^{}`$ are would-be Goldstone fields. \] We use Greek and Latin letters to denote up- and down-type quarks respectively, and the external quark is denoted as $`u_e`$ or $`d_e`$. There are four groups of contributions, denoted as $`WW`$, $`WG`$, $`GW`$ and $`GG`$, where the first and second letters refer to the bosons exchanged in the outer and inner loops respectively.
Consider for example the $`u_e`$ quark EDM. We found as in the SM that the contribution from $`WW`$ is automatically cancelled without using any unitarity conditions while other contributions have the following separate structure:
$$V_{ek}V_{\alpha k}^{}V_{\alpha j}V_{ej}^{}[H_\alpha (k)H_\alpha (j)],$$
(4)
where $`H_\alpha (k)`$ is a function of masses of $`u_\alpha `$, $`d_k`$, $`u_e`$ and $`W^\pm `$. The crucial point is that the dependence on $`d_j`$ and $`d_k`$ masses splits up. This is mainly responsible for the thorough or partial cancellation occurring in the SM and beyond. Summing up the pair of mirror-reflected diagrams (i.e., $`jk`$) doubles the imaginary part of the product of CKM matrix elements while removing its real part. The summation over six pairs of $`(jk)`$, since we have four down-type quarks $`d_i`$, $`d_j`$, $`d_k`$ and $`d_l`$, completely cancels contributions amongst themselves due to the unitarity of the $`4\times 4`$ matrix $`V`$. For example, the $`H_\alpha (j)`$ term is
$$\begin{array}{cc}& 2i\mathrm{Im}[V_{ej}V_{\alpha j}^{}(V_{ek}^{}V_{\alpha k}+V_{el}^{}V_{\alpha l}+V_{ei}^{}V_{\alpha i})]\hfill \\ \hfill =& 2i\mathrm{Im}[V_{ej}V_{\alpha j}^{}(\delta _{e\alpha }V_{ej}^{}V_{\alpha j})]=0.\hfill \end{array}$$
(5)
In other words, the up-type quark EDM’s vanish strictly in the model with an extra down-type singlet quark as in the SM. However, for the down-type quark EDM’s in the model, we have one less up-type quark in virtual loops to complete the unitarity cancellation so that the cancellation is not thorough. For example, summing over the pairs of $`(\alpha \beta )`$ in the contribution to the $`d_e`$ quark EDM (see Fig. 1 for notations) and using the unitarity of $`V`$, the $`H_i(\alpha )`$ term is
$$\begin{array}{cc}& 2i\mathrm{Im}[V_{\alpha e}^{}V_{\alpha i}(V_{\beta i}^{}V_{\beta e}+V_{\gamma i}^{}V_{\gamma e})]\hfill \\ \hfill =& 2i\mathrm{Im}[V_{\alpha e}V_{\alpha i}^{}V_{0i}V_{0e}^{}],\hfill \end{array}$$
(6)
which is generally non-vanishing. Here $`u_\alpha `$, $`u_\beta `$ and $`u_\gamma `$ are the three up-type quarks. Therefore, the $`d_e`$ quark EDM is proportional to,
$$2i\underset{i,\alpha }{}\mathrm{Im}[V_{\alpha e}V_{\alpha i}^{}V_{0i}V_{0e}^{}]H_i(\alpha ).$$
(7)
Note that to avoid complete cancellation due to $`_i(V_{\alpha i}^{}V_{0i})=\delta _{\alpha 0}=0`$, $`H_i(\alpha )`$ must involve the $`d_i`$ mass.
Now let us evaluate analytically the $`d_e`$ quark EDM, $`d(d_e)`$. This is facilitated by the mass hierarchy in the SM, $`m_tm_Wm_q`$, where $`q`$ stands for other five quarks, and the assumption that $`m_{d_4}m_t`$. We should discriminate two kinds of down-type quarks with $`d_4`$ heavy and others light, as well as two kinds of up-type quarks with top heavy and others light. We want to retain only the terms that are least suppressed by light quark masses. Since the charged current is purely left-handed, the chirality flip needed for the EDM operator has to be made by the external quark mass. We found that $`H_i(\alpha )`$ is proportional to $`m_{u_\alpha }^2`$ for all of $`WG`$, $`GW`$ and $`GG`$ contributions when $`u_\alpha `$ is light. Therefore we only need to keep the top quark in up-type quarks. The leading terms involving the heavy $`d_4`$ quark come from the $`WG`$ and $`GG`$ contributions:
$$\begin{array}{ccc}\hfill d(d_e)_{\mathrm{heavy}}& =& em_{d_e}G_F^2m_W^2(4\pi )^4\mathrm{Im}[V_{te}V_{t4}^{}V_{04}V_{0e}^{}]\hfill \\ \multicolumn{3}{c}{[Q_u(\frac{23}{9}\frac{8}{9}\mu _t+\frac{4}{3}\mu _t\mathrm{ln}\frac{\mu _4}{\mu _t}\frac{16}{3}\mathrm{ln}\mu _t)}\\ \multicolumn{3}{c}{+Q_d(\frac{59}{9}\frac{1}{3}\mu _t2\mu _t\mathrm{ln}\frac{\mu _4}{\mu _t}+10\mathrm{ln}\mu _t)],}\end{array}$$
(8)
where $`\mu _t=m_t^2/m_W^2`$ and $`\mu _4=m_{d_4}^2/m_W^2`$. The leading terms involving the light $`d_i`$ quarks are independent of their masses so that we may use $`_iV_{\alpha i}^{}V_{0i}=V_{\alpha 4}^{}V_{04}`$ to sum up their contributions and obtain,
$$\begin{array}{ccc}\hfill d(d_e)_{\mathrm{light}}& =& em_{d_e}G_F^2m_W^2(4\pi )^4\mathrm{Im}[V_{te}V_{t4}^{}V_{04}V_{0e}^{}]\hfill \\ \multicolumn{3}{c}{[Q_u(8+\frac{16\pi ^2}{3}12\mathrm{ln}\mu _t+8\mathrm{ln}^2\mu _t)}\\ \multicolumn{3}{c}{+Q_d(4\frac{8\pi ^2}{3}+4\mathrm{ln}\mu _t8\mathrm{ln}^2\mu _t)],}\end{array}$$
(9)
which comes from the $`GW`$ contribution. We should emphasize that the above results are obtained in the limit of large mass hierarchy so that the cancellation of the EDM in the degeneracy limit of down-type quarks cannot be explicit therein.
We note that the $`d_e`$ quark EDM, $`d(d_e)`$, occurs at order $`g^4(m_{d_e}/m_W^2)`$ and there is no further suppression due to GIM mechanism. This is typical for models without CKM unitarity. Some terms are further enhanced by the heavy top mass, while the absence of the heaviest $`m_{d_4}^2`$ enhancement is consistent with general arguments based on gauge invariance and naive dimensional analysis.
For numerical analysis, we use the following parameters: $`G_F=1.2\times 10^5\mathrm{GeV}^2`$, $`m_W=80`$ GeV and $`m_t=175`$ GeV. Then we have for the $`d`$ quark,
$$\begin{array}{ccc}\hfill d(d)_{\mathrm{heavy}}& =& \mathrm{Im}[V_{td}V_{t4}^{}V_{04}V_{0d}^{}]\frac{m_d}{10\mathrm{MeV}}\hfill \\ \multicolumn{3}{c}{(5.3,0.86,+2.3)\times 10^{26}\mathrm{e}\mathrm{cm},}\\ \multicolumn{3}{c}{\mathrm{for}m_{d_4}=(200,300,400)\mathrm{GeV},}\\ \hfill d(d)_{\mathrm{light}}& =& \mathrm{Im}[V_{td}V_{t4}^{}V_{04}V_{0d}^{}]\frac{m_d}{10\mathrm{MeV}}\hfill \\ \multicolumn{3}{c}{3.3\times 10^{25}\mathrm{e}\mathrm{cm}.}\end{array}$$
(10)
Note that the same combination of CKM elements is involved in the two contributions. Since the ‘heavy’ part is generally smaller by one order of magnitude we retain below only the ‘light’ part which is independent of the $`d_4`$ mass. The product of CKM elements can be expressed in terms of the $`3\times 3`$ submatrix elements, e.g.,
$$\mathrm{Im}[V_{td}V_{t4}^{}V_{04}V_{0d}^{}]=\mathrm{Im}[V_{td}V_{tb}^{}Z_{bd}]\mathrm{Im}[V_{ts}V_{td}^{}Z_{ds}],$$
(11)
where $`Z_{ij}=V_{ui}V_{uj}^{}+V_{ci}V_{cj}^{}+V_{ti}V_{tj}^{},i,j=d,s,b`$ are precisely the couplings appearing in the FCNC $`Z`$ interactions amongst ordinary down-type quarks. The most stringent bounds on them come from the neutral meson mixing and FCNC decays. Here we adopt the bounds obtained by requiring that the new tree level FCNC effects do not exceed the experimental values. Some bounds may be relaxed if destructive interference occurs between them and other contributions, e.g., the box diagrams. We take $`|Z_{ds}|3\times 10^4,|Z_{sb}||Z_{bd}|10^3`$. These bounds are actually interrelated with the extraction of other elements like $`V_{td}`$ and $`V_{ts}`$. We do not attempt here a global analysis which is beyond the main interest of the present work, but simply adopt the following values for numerical estimate: $`|V_{td}|0.02,|V_{ts}||V_{cb}|0.04.`$ Then, $`\mathrm{Im}[V_{td}V_{t4}^{}V_{04}V_{0d}^{}]2\times 10^5`$, and
$$\begin{array}{c}|d(d)_{\mathrm{light}}|6.5\times 10^{30}\frac{m_d}{10\mathrm{MeV}}\mathrm{e}\mathrm{cm}.\hfill \end{array}$$
(12)
Using the $`SU(6)`$ relation for the neutron, we obtain
$$|d(n)|0.8\times 10^{29}\frac{m_d}{10\mathrm{MeV}}\mathrm{e}\mathrm{cm}.$$
(13)
A few comments are in order.
(1) The non-unitarity of the CKM matrix generally induces FCNC interactions of the $`Z`$ and Higgs bosons with down-type quarks. We found that mixed exchanges of $`W^\pm `$ and $`Z`$ (or Higgs) can contribute to the EDM. Since the Higgs boson decouples for heavy enough mass, we consider here the FCNC $`Z`$ couplings. There are two types of diagrams. The first one is similar to Fig. 1 and its leading term for the $`d(d)`$ contains the same matrix elements as in Eq.(10) and decreases the above result by about $`1/3`$ for $`m_{d_4}=300\mathrm{G}\mathrm{e}\mathrm{V}`$, while the $`d(u)`$ is severely suppressed. The second type has a different topology and involves the complicated issue of the renormalization of the non-unitary CKM matrix. The details of all this will be reserved for a future publication. The point we want to make here is that the purely charged current result in Eq.(13) gives us the correct order of magnitude of the neutron EDM since there is no plausible reason to expect strong cancellation between the charged current and FCNC contributions which generally involve several different Jarlskog rephasing invariants of $`CP`$ violation.
(2) The above discussion can be easily generalized to other models with exotic quarks. At the two loop level both up- and down-type quark EDM’s vanish strictly in the model with a sequential fourth generation. In the model with an extra up-type singlet quark, the down-type quark EDM’s vanish identically at two loop order. Since all down-type quarks are light in this case, the leading terms for the $`u_e`$ EDM must be proportional to $`m_{u_e}m_{d_i}^2`$ and are thus very small compared to the case considered above.
(3) We have also studied the $`P`$ and $`T`$ violating purely gluonic operators, e.g., the dimension-$`6`$ Weinberg operator. We found that they are severely suppressed by light quark masses in the current case as in the SM . Their contribution to the neutron EDM can be ignored. To obtain the quark chromoelectric dipole moment(CEDM), one merely replaces $`eQ_{u,d}`$ by $`g_s`$ in Eqns.$`(8)`$ and $`(9)`$ with $`g_s`$ being the QCD coupling. It is then clear that the quark CEDM is much smaller than the quark EDM so that the latter remains to be the dominant contribution in the neutron EDM.
(4)In the static limit the fermionic part of the EDM is identical to the spin operator. Since a significant amount of the proton spin is derived from the polarized strange quark sea, it seems reasonable that the strange quark also contributes to the neutron EDM. For the strange quark EDM, we have enhancement factors from masses and CKM elements. For the latter the dominant one is $`\mathrm{Im}[V_{tb}V_{ts}^{}Z_{sb}]`$. Suppose $`|\eta |10\%`$ of the proton spin is accounted for by the strange sea, then we have $`|d(n)_{\mathrm{strange}}|0.8\times 10^{29}(m_s/m_d)(|V_{ts}|/|V_{td}|)|\eta |\mathrm{e}\mathrm{cm}3.2\times 10^{29}\mathrm{e}\mathrm{cm}`$ for the quoted parameters and $`m_s=200`$ MeV.
(5) Another possible enhancement for the neutron EDM originates from long distance physics. Although there are still controversies concerning this, it seems reliable to get an enhancement factor of two orders of magnitude in the SM. If this persists in the case considered here, improvement of the upper bound on the neutron EDM in the near future will provide an interesting test of the unitarity of the CKM matrix which will be complementary to or even competitive with the bound from $`B`$ physics.
The work of X. Li was supported in part by the China National Natural Science Foundation under grant Numbers 19835060 and 19875072.
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# Global Bounds for the Lyapunov Exponent and the Integrated Density of States of Random Schrödinger Operators in One Dimension
## 1. Introduction
In this article we will consider random Schrödinger operators $`H(\omega )`$ in $`L^2()`$ of the form
(1)
$$H(\omega )=H_0+V_\omega ,H_0=\frac{d^2}{dx^2},V_\omega =\underset{j}{}\alpha _j(\omega )f(j),$$
where $`\{\alpha _j(\omega )\}_j`$ is a sequence of i.i.d. (independent, identically distributed) variables on a complete probability space $`(\mathrm{\Omega },,)`$ having a common distribution measure $`\kappa `$ (i.e. $`\{\alpha _j\mathrm{\Delta }\}=\kappa (\mathrm{\Delta })`$ for any Borel set $`\mathrm{\Delta }`$). In what follows we always suppose that $`\kappa `$ is supported on a compact interval and the single-site potential $`f`$ is integrable with support in the interval \[-1/2,1/2\]. Moreover, the random variables are assumed to form a stationary, metrically transitive random field, i.e. there are measure preserving ergodic transformations $`\{T_j\}_j`$ such that $`\alpha _j(T_k\omega )=\alpha _{jk}(\omega )`$ for all $`\omega \mathrm{\Omega }`$. The spectral properties of the operator (1) were studied in detail in . The results are most complete for the case when $`f`$ is the point interaction (see ).
The integrated density of states $`N(E)`$ and the Lyapunov exponent $`\gamma (E)`$ are important quantities associated with operators of the form (1) (see e.g. ). In particular, according to Ishii-Pastur-Kotani theorem the set $`\{E:\gamma (E)=0\}`$ is the essential support of the absolute continuous part of the spectral measure for $`H(\omega )`$.
The main idea of our approach is to approximate the operator (1) by means of the sequence
$$H^{(n)}(\omega )=H_0+\underset{j=n}{\overset{n}{}}\alpha _j(\omega )f(j)$$
with unchanged $`H_0`$, which converges to $`H(\omega )`$ in the strong resolvent sense. This differs from the usual approach where one puts the whole system in a box, which then tends to infinity (see e.g. ). In (see also ) we used this approximation to invoke scattering theory for the study the spectral properties of the limiting operator (1). Some other applications of scattering theory to the study of spectral properties of such type Schrödinger operators in one dimension can be found in and .
One of the important ingredients of our approach developed in is the Lifshitz-Krein spectral shift function. The spectral shift function naturally replaces the eigenvalue counting function usually used to construct the density of states for the operator (1). The celebrated Birman-Krein theorem (see e.g. ) relates the spectral shift function to scattering theory. In fact, up to a factor $`\pi ^1`$ it may be identified with the scattering phase for the pair ($`H^{(n)}(\omega )`$, $`H_0`$), i.e. $`\xi ^{(n)}(E;\omega )=\pi ^1\delta ^{(n)}(E;\omega )`$ when $`E>0`$,
$$\delta ^{(n)}(E;\omega )=\frac{1}{2i}\mathrm{log}detS^{(n)}(E;\omega )=\frac{1}{2i}\mathrm{log}det\left(\begin{array}{cc}T_\omega ^{(n)}(E)\hfill & \hfill R_\omega ^{(n)}(E)\\ L_\omega ^{(n)}(E)\hfill & \hfill T_\omega ^{(n)}(E)\end{array}\right).$$
Here $`|T^{(n)}(E)|^2`$ and $`|R^{(n)}(E)|^2=|L^{(n)}(E)|^2`$ have the meaning of transmission and reflection coefficients, respectively, such that $`|T^{(n)}(E)|^2+|R^{(n)}(E)|^2=1`$. For $`E<0`$ the quantity $`\xi ^{(n)}(E;\omega )`$ equals minus the counting function for $`H^{(n)}(\omega )`$.
In particular in we proved the almost sure existence of the limit
(2)
$$\xi (E)=\underset{n\mathrm{}}{lim}\frac{\xi ^{(n)}(E;\omega )}{2n+1},$$
which we called the spectral shift density. Also we proved the equality $`\xi (E)=N_0(E)N(E)`$, where $`N(E)`$ and $`N_0(E)=\pi ^1[\mathrm{max}(0,E)]^{1/2}`$ are the integrated density of states of the Hamiltonians $`H(\omega )`$ and $`H_0`$ respectively. This result also extends to higher dimension in the continuous and discrete cases. Also we showed that almost surely the Lyapunov exponent $`\gamma (E)`$ at energy $`E>0`$ is given as
(3)
$$\gamma (E)=\underset{n\mathrm{}}{lim}\frac{\mathrm{log}|T^{(n)}(E;\omega )|}{2n+1},$$
where $`T^{(n)}(E,\omega )`$ is the transmission amplitude for the pair of Hamiltonians ($`H^{(n)}(\omega )`$, $`H_0`$) at energy $`E`$. We recall that $`\gamma (E)`$ is defined as the upper Lyapunov exponent for the fundamental matrix at energy $`E`$ of the Schrödinger operator $`H(\omega )`$. The connection between the Lyapunov exponent and the transmission coefficient $`|T_\omega ^{(n)}(E)|`$ was recognized long ago . A complete proof has appeared in .
We note that the theory of the spectral shift function was also recently used to show that the integrated density of states is independent of the choice of boundary conditions on the sides of a large box, in which the system is put.
The conditions on the random variables $`\alpha _j`$ and the single-site potential $`f`$ stated above are slightly weaker than those in . However the results of which will be used below remain valid also in this more general case.
The aim of the present paper is to prove *global* bounds for the Lyapunov exponent and the integrated density of states, i.e. bounds which hold for all $`E>0`$ and describe the correct asymptotic behavior in the limit $`E\mathrm{}`$. These results are formulated as Theorems 1 and 2 below. To the best of our knowledge the first article to look for the asymptotic behavior of $`\gamma (E)`$ and $`N(E)`$ in the limit $`E\mathrm{}`$ is . The best known estimate for the integrated density of states is due to Kirsch and Martinelli \[10, Corollary 3.1\]. This bound however does not reproduce the correct asymptotic behavior of $`N(E)`$ in the large energy limit. Another estimate, which is due to Pastur and Figotin (see \[20, Sec. V.11.B\]), is valid for an $``$-metrically transitive random field. Since our potential $`V_\omega (x)`$ is a $``$-metrically transitive field this estimate does not apply directly to the present situation. Our two-sided estimate leads to the bound (23) below which is very close to that of Pastur and Figotin.
In what follows $`C`$ will denote a finite positive generic constant varying with the context, but which depends only on $`f`$ and $`\kappa `$.
We are indebted to Leonid Pastur for reading the preliminary version of this article.
## 2. The Lyapunov exponent
We recall that the scattering matrix $`S(E)`$ for a pair of Hamiltonians ($`H`$, $`H_0`$) on $`L^2()`$ at fixed energy $`E0`$ is a $`2\times 2`$ unitary matrix
(4)
$$S(E)=\left(\begin{array}{cc}T(E)& R(E)\\ L(E)& T(E)\end{array}\right),$$
where $`L(E)`$ and $`R(E)`$ denote the left and right reflection amplitudes respectively. The transmission amplitude $`T(E)`$ can vanish only for $`E=0`$ (see ). To any S-matrix (4) we associate the unimodular matrix
$$\mathrm{\Lambda }(E)=\left(\begin{array}{cc}\frac{1}{T(E)}& \frac{R(E)}{T(E)}\\ \frac{L(E)}{T(E)}& \frac{1}{\overline{T(E)}}\end{array}\right).$$
Let $`T_\alpha (E)`$, $`R_\alpha (E)`$, $`L_\alpha (E)`$ be the elements of the S-matrix at energy $`E`$ for the pair of operators ($`H_0+\alpha f`$, $`H_0`$) and $`\mathrm{\Lambda }_\alpha (E)`$ the corresponding $`\mathrm{\Lambda }`$-matrix. Also let $`\stackrel{~}{\mathrm{\Lambda }}_\alpha (E)=U_E^{1/2}\mathrm{\Lambda }_\alpha (E)U_E^{1/2}`$ with
$$U_E=\left(\begin{array}{cc}e^{i\sqrt{E}}& 0\\ 0& e^{i\sqrt{E}}\end{array}\right).$$
Explicitly we have
$$\stackrel{~}{\mathrm{\Lambda }}_\alpha (E)=\left(\begin{array}{cc}\frac{e^{i\sqrt{E}}}{T_\alpha (E)}& \frac{R_\alpha (E)}{T_\alpha (E)}\\ \frac{L_\alpha (E)}{T_\alpha (E)}& \frac{e^{i\sqrt{E}}}{\overline{T(E)}}\end{array}\right).$$
Consider the matrix
$$A(E)=𝔼\left\{\stackrel{~}{\mathrm{\Lambda }}_{\alpha (\omega )}(E)^{}\stackrel{~}{\mathrm{\Lambda }}_{\alpha (\omega )}(E)\right\}=\stackrel{~}{\mathrm{\Lambda }}_\alpha (E)^{}\stackrel{~}{\mathrm{\Lambda }}_\alpha (E)𝑑\kappa (\alpha )0,$$
where for brevity we write $`\alpha (\omega )`$ instead of $`\alpha _j(\omega )`$ with some $`j`$. Let $`\beta _+(E)`$ be the largest eigenvalue of $`A(E)`$ and $`\beta _{}(E)`$ the smallest. It will turn out below that $`\beta _+(E)1`$. Set $`\stackrel{~}{\gamma }(E)=(\mathrm{log}\beta _+(E))/20`$.
The first main result of the present article is
###### Theorem 1.
Given the Hamiltonian (1) and the distribution $`\kappa `$ for the coupling constant $`\alpha `$, for all $`E>0`$ the resulting Lyapunov exponent satisfies the upper bound
(5)
$$\gamma (E)\stackrel{~}{\gamma }(E).$$
In particular $`\gamma (E)`$ decays at least like $`1/\sqrt{E}`$ at infinity.
###### Proof.
Let $`\mathrm{\Lambda }^{(n)}(E;\omega )`$ denote the $`\mathrm{\Lambda }`$-matrix for the pair ($`H^{(n)}(\omega )`$, $`H_0`$), which by the factorization property can be represented in the form
(6)
$$\mathrm{\Lambda }^{(n)}(E;\omega )=U_E^{n1/2}\underset{j=n}{\overset{n}{}}\stackrel{~}{\mathrm{\Lambda }}_{\alpha _j(\omega )}(E)U_E^{n1/2}.$$
In fact, this factorization property is a consequence of the multiplicativity property of the fundamental matrix (see for a proof and for references to earlier work). A short calculation gives
(7)
$$|T^{(n)}(E;\omega )|^2=\frac{1}{4}\mathrm{tr}\left(\mathrm{\Lambda }^{(n)}(E;\omega )^{}\mathrm{\Lambda }^{(n)}(E;\omega )\right)+\frac{1}{2}.$$
With $`𝔼`$ denoting the expectation with respect to the measure $``$, by Jensen’s inequality and (7) we therefore have the estimate
(8) $`e^{2𝔼\{\mathrm{log}|T^{(n)}(E;\omega )|\}}𝔼\left\{|T^{(n)}(E;\omega )|^2\right\}`$
$`=`$ $`{\displaystyle \frac{1}{4}}𝔼\left\{\mathrm{tr}\left(\mathrm{\Lambda }^{(n)}(E;\omega )^{}\mathrm{\Lambda }^{(n)}(E;\omega )\right)\right\}+{\displaystyle \frac{1}{2}}.`$
From the factorization property (6) it follows that
(9)
$$\mathrm{tr}\left(\mathrm{\Lambda }^{(n)}(E;\omega )^{}\mathrm{\Lambda }^{(n)}(E;\omega )\right)=\mathrm{tr}\left(\underset{j=n}{\overset{n}{}}\stackrel{~}{\mathrm{\Lambda }}_{\alpha _j(\omega )}(E)^{}\underset{j=n}{\overset{n}{}}\stackrel{~}{\mathrm{\Lambda }}_{\alpha _j(\omega )}(E)\right).$$
We will now make use of the fact that the $`\alpha _k(\omega )`$ are i.i.d. random variables. For this purpose define the $`2\times 2`$ matrices $`A_j(E)0`$ recursively by $`A_0=I`$ and
(10)
$$A_j(E)=\stackrel{~}{\mathrm{\Lambda }}_\alpha (E)^{}A_{j1}(E)\stackrel{~}{\mathrm{\Lambda }}_\alpha (E)𝑑\kappa (\alpha ),$$
such that in particular $`A(E)=A_1(E)`$. Now it is easy to see that
(11) $`𝔼\left\{\mathrm{tr}\left(\mathrm{\Lambda }^{(n)}(E;\omega )^{}\mathrm{\Lambda }^{(n)}(E;\omega )\right)\right\}`$
$`=`$ $`\mathrm{tr}\left(𝔼\left(\mathrm{\Lambda }^{(n)}(E;\omega )^{}\mathrm{\Lambda }^{(n)}(E;\omega )\right)\right)=A_{2n+1}(E).`$
We now use the fact that the operator inequality $`0AA^{}`$ implies $`0\mathrm{tr}A\mathrm{tr}A^{}`$ and $`B^{}ABB^{}A^{}B`$ for all $`B`$. In particular we have $`A(E)\beta _+(E)I`$ from which we obtain the recursive estimates $`A_j(E)\beta _+(E)A_{j1}(E)\mathrm{}\beta _+(E)^jI`$ and hence
(12)
$$𝔼\left\{\mathrm{tr}\left(\mathrm{\Lambda }^{(n)}(E;\omega )^{}\mathrm{\Lambda }^{(n)}(E;\omega )\right)\right\}2\beta _+(E)^{2n+1}.$$
We remark that with the same arguments one proves the lower bound
$$2\beta _{}(E)^{2n+1}𝔼\left(\mathrm{tr}\left(\mathrm{\Lambda }^{(n)}(E;\omega )^{}\mathrm{\Lambda }^{(n)}(E;\omega )\right)\right).$$
The relation (3), the estimate (12) combined with (8) and Fatou’s lemma imply now
$`\gamma (E)`$ $``$ $`{\displaystyle \frac{1}{2}}\underset{n\mathrm{}}{lim}{\displaystyle \frac{\mathrm{log}𝔼\left\{|T^{(n)}(E;\omega )|^2\right\}}{2n+1}}`$
$``$ $`{\displaystyle \frac{1}{2}}\underset{n\mathrm{}}{lim}{\displaystyle \frac{\mathrm{log}\left(\beta _+(E)^{2n+1}/2+1/2\right)}{2n+1}}={\displaystyle \frac{1}{2}}\mathrm{log}\beta _+(E),`$
which proves the claim (5).
To establish the last claim of the theorem we recall the following well known estimates (see e.g. )
(13)
$$|T_\alpha (E)1|+|R_\alpha (E)|C\frac{1}{\sqrt{E}}$$
valid for all large $`E>0`$ uniformly for all $`\alpha `$ in the (compact) support of $`\kappa `$ for fixed $`f`$. Using the estimate (13) in (14) gives the estimate $`\beta _+(E)1+C/\sqrt{E}`$ for all large $`E`$. Since $`\stackrel{~}{\gamma }(E)=(\mathrm{log}\beta _+(E))/2`$, this concludes the proof of the theorem. ∎
Since $`\gamma (E)0`$, we obviously have the inequality $`\beta _+(E)1`$ for almost all $`E`$. We will give now a direct independent proof of this fact and simultaneously obtain an expression for $`\beta _+(E)`$. The matrix $`A(E)`$ may be written in the form
$$A(E)=\left(\begin{array}{cc}a(E)& b(E)\\ \overline{b(E)}& a(E)\end{array}\right)$$
with
(14) $`a(E)`$ $`=`$ $`{\displaystyle \left(\frac{2}{|T_\alpha (E)|^2}1\right)𝑑\kappa (\alpha )}`$
(15) $`b(E)`$ $`=`$ $`e^{i\sqrt{E}}{\displaystyle \frac{R_\alpha (E)}{T_\alpha (E)^2}𝑑\kappa (\alpha )}.`$
This gives the two eigenvalues of $`A(E)`$ in the form
(16)
$$\beta _\pm (E)=a(E)\pm |b(E)|$$
Obviously $`a(E)1`$ and hence $`\beta _+(E)1`$. In fact, $`a(E)=1`$ is possible if and only if $`R_\alpha (E)=0`$ for almost all $`\alpha `$ in the support of $`\kappa `$. Then also $`b(E)=0`$ and $`\beta _+(E)=1`$. Actually (if $`\mathrm{supp}\kappa `$ has at least one non-isolated point) we do not believe there are nontrivial $`f`$ and $`E`$ for which this holds but in any case for such $`E`$’s the Lyapunov exponent vanishes as is easily verified (see also ), so this is a trivial confirmation of estimate (5) in this case. In the remaining case we trivially have $`\beta _+(E)>1`$.
As an example we consider the random Kronig - Penney model which is formally obtained from $`H(\omega )`$ by replacing $`f`$ with the Dirac $`\delta `$-function at the origin. Then we have (correcting for a misprint on page 232 of )
(17) $`T_\alpha (E)`$ $`=`$ $`\left(1+i{\displaystyle \frac{\alpha }{2\sqrt{E}}}\right)^1`$
(18) $`R_\alpha (E)`$ $`=`$ $`i{\displaystyle \frac{\alpha }{2\sqrt{E}}}\left(1+i{\displaystyle \frac{\alpha }{2\sqrt{E}}}\right)^1`$
and our method still applies. This gives
(19) $`a(E)`$ $`=`$ $`1+{\displaystyle \frac{\alpha ^2}{4E}}`$
(20) $`b(E)`$ $`=`$ $`i{\displaystyle \frac{\alpha }{2\sqrt{E}}}{\displaystyle \frac{\alpha ^2}{4E}}.`$
Here for brevity by $``$ we denote the mean with respect to the probability measure $``$ such that
$$\alpha =𝔼\{\alpha (\omega )\}=\alpha 𝑑\kappa (\alpha ),\alpha ^2=𝔼\{\alpha (\omega )^2\}=\alpha ^2𝑑\kappa (\alpha ).$$
In particular (19) gives
(21)
$$\beta _+(E)=1+\frac{\alpha ^2}{4E}+\frac{1}{2\sqrt{E}}\left(\frac{\alpha ^2}{4E}+\alpha ^2\right)^{1/2}.$$
So also in this case $`\gamma (E)`$ decays at least like $`1/\sqrt{E}`$ as $`E\mathrm{}`$ and at least like $`1/E`$ if the mean $`\alpha `$ of $`\alpha `$ vanishes, i.e. if on average the coupling constant is zero.
## 3. The integrated density of states
We denote by $`\xi _\alpha (E)`$ the spectral shift function for the pair $`(H_0+\alpha f,H_0)`$. The second main result of this article is given by
###### Theorem 2.
For all $`E>0`$ the spectral shift density $`\xi (E)`$ for the operator (1) satisfies the following two-side bound
(22)
$$𝔼\{\xi _{\alpha (\omega )}(E)\}r(E)\xi (E)𝔼\{\xi _{\alpha (\omega )}(E)\}+r(E),$$
where
$$r(E)=\mathrm{min}\{\frac{1}{2},\frac{1}{\pi }𝔼\left\{\frac{|R_{\alpha (\omega )}(E)}{1|R_{\alpha (\omega )}(E)|}\right\}\}.$$
In particular $`𝔼\{\xi _{\alpha (\omega )}(E)\}`$ and $`r(E)`$ decays at least like $`1/\sqrt{E}`$ at infinity.
Remarks. 1. One can easily prove the following estimate
$$𝔼\{\xi _{\alpha (\omega )}(E)\}1\xi (E)𝔼\{\xi _{\alpha (\omega )}(E)\}+1,$$
which is valid for all $`E`$.
2. By the monotonicity of the spectral shift function with respect to perturbation $`\xi (E)0`$ if $`\mathrm{supp}\kappa _+`$ and $`\xi (E)0`$ if $`\mathrm{supp}\kappa _{}`$ for almost all $`E>0`$.
3. For large $`E>0`$ by (13)
$$r(E)=\mathrm{min}\{\frac{1}{2},\frac{1}{\pi }𝔼\left\{\frac{|R_{\alpha (\omega )}(E)}{1|R_{\alpha (\omega )}(E)|}\right\}\}=\frac{1}{\pi }𝔼\left\{\frac{|R_{\alpha (\omega )}(E)}{1|R_{\alpha (\omega )}(E)|}\right\}\frac{C}{\sqrt{E}}.$$
4. In we proved the relation $`\xi (E)=N_0(E)N(E)=\sqrt{E}/\pi N(E)`$, where $`N_0(E)`$ is the integrated density of states for the free operator $`H_0`$. Theorem 2 then gives the following two-sided bound for the integrated density of states
(23)
$$\frac{\sqrt{E}}{\pi }𝔼\{\xi _{\alpha (\omega )}(E)\}r(E)N(E)\frac{\sqrt{E}}{\pi }𝔼\{\xi _{\alpha (\omega )}(E)\}+r(E),E>0.$$
There are some other upper bounds on the integrated density of states. A well-known result is a one-sided bound due to Kirsch and Martinelli \[10, Corollary 3.1\],
$$N(E)\frac{C}{\sqrt{\eta }}𝔼\left\{_{1/2}^{1/2}(E+\eta V_\omega (x))_+𝑑x\right\}$$
for any $`\eta >0`$ and all $`E`$. This bound however does not reproduce the correct asymptotic behavior of $`N(E)`$ in the large energy limit.
5. The bounds (5) and (22) are of interest in the context of the Thouless formula (see e.g. )
(24)
$$\gamma (E)\gamma _0(E)=_{}\mathrm{log}|EE^{}|d\xi (E^{}),E,$$
where $`\gamma _0(E)=[\mathrm{max}(0,E)]^{1/2}`$ is the Lyapunov exponent for $`H_0`$. The Thouless formula in the form (24) can be viewed as a subtracted dispersion relation (see e.g. ).
###### Proof.
In we proved (see Theorem 3.3 there and its proof) that for any two potentials $`V_1`$ and $`V_2`$ with (compact) disjoint supports one has
$$\xi (E;H_0+V_1+V_2,H_0)=\xi (E;H_0+V_1,H_0)+\xi (E;H_0+V_2,H_0)+\xi _{12}(E)$$
with
$$\xi _{12}(E)=\frac{1}{2\pi i}\mathrm{log}\frac{1R_1(E)L_2(E)}{1\overline{R_1(E)}\overline{L_2(E)}},$$
where $`R_k(E)`$ and $`L_k(E)`$ are the right and left reflection coefficients for the Schrödinger equation with the potential $`V_k`$, $`k=1,2`$. Actually Theorem 3.3 in states that $`|\xi _{12}(E)|1/2`$ for all $`E0`$. Now we improve on this estimate. As in we set
$$L_k(E)=a_k(E)e^{i\delta _k^{(L)}},R_k(E)=a_k(E)e^{i\delta _k^{(R)}},k=1,2$$
with $`0a_k(E)1`$. Moreover $`a_k(E)=1`$ only when $`T_k(E)=0`$, which we recall can happen only if $`E=0`$. Therefore
$`\mathrm{log}{\displaystyle \frac{1R_1(E)L_2(E)}{1\overline{R_1(E)}\overline{L_2(E)}}}`$
$`=\mathrm{log}{\displaystyle \frac{1a_1(E)a_2(E)e^{i(\delta _1^{(R)}+\delta _2^{(L)})}}{1a_1(E)a_2(E)e^{i(\delta _1^{(R)}+\delta _2^{(L)})}}}`$
$`=2i\mathrm{arctan}{\displaystyle \frac{a_1(E)a_2(E)\mathrm{sin}(\delta _1^{(R)}+\delta _2^{(L)})}{1a_1(E)a_2(E)\mathrm{cos}(\delta _1^{(R)}+\delta _2^{(L)})}}.`$
By means of the inequality $`|\mathrm{arctan}x||x|`$ we immediately obtain
(25)
$$|\xi _{12}(E)|\mathrm{min}\{\frac{1}{2},\frac{1}{\pi }\frac{a_1(E)a_2(E)}{1a_1(E)a_2(E)}\}.$$
Since $`0a_k(E)<1`$ we can replace $`a_1(E)a_2(E)(1a_1(E)a_2(E))^1`$ either by $`a_1(E)(1a_1(E))^1`$ or by $`a_2(E)(1a_2(E))^1`$.
Now let us consider the operator $`H^{(n)}(\omega )`$ for finite $`n`$. Applying the inequality (25) we obtain
$`\left|\xi ^{(n)}(E;\omega )\xi _{\alpha _n(\omega )}(E)\xi _{\alpha _n(\omega )}(E)\xi ^{(n1)}(E;\omega )\right|`$
$``$ $`\mathrm{min}\{{\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{\pi }}{\displaystyle \frac{|R_{\alpha _n(\omega )}(E)|}{1|R_{\alpha _n(\omega )}(E)|}}\}+\mathrm{min}\{{\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{\pi }}{\displaystyle \frac{|R_{\alpha _n(\omega )}(E)|}{1|R_{\alpha _n(\omega )}(E)|}}\}`$
Repeating this procedure recursively we obtain
$$\left|\xi ^{(n)}(E;\omega )\underset{j=n}{\overset{n}{}}\xi _{\alpha _j(\omega )}(E)\right|\underset{j=n}{\overset{n}{}}\mathrm{min}\{\frac{1}{2},\frac{1}{\pi }\frac{|R_{\alpha _j(\omega )}(E)|}{1|R_{\alpha _j(\omega )}(E)|}\}.$$
From the existence of the spectral shift density (2) by the Birkhoff ergodic theorem it follows that
$$\left|\xi (E)𝔼\left\{\xi _{\alpha (\omega )}(E)\right\}\right|𝔼\left\{\mathrm{min}\{\frac{1}{2},\frac{1}{\pi }\frac{|R_{\alpha (\omega )}(E)|}{1|R_{\alpha (\omega )}(E)|}\}\right\}.$$
From the obvious inequality
$$𝔼\left\{\mathrm{min}\{\frac{1}{2},\frac{1}{\pi }\frac{|R_{\alpha (\omega )}(E)|}{1|R_{\alpha (\omega )}(E)|}\}\right\}\mathrm{min}\{\frac{1}{2},\frac{1}{\pi }𝔼\left\{\frac{|R_{\alpha (\omega )}(E)|}{1|R_{\alpha (\omega )}(E)|}\right\}\}$$
the bound (22) follows.
For large $`E`$ we have the following asymptotics uniformly in $`\alpha `$ on compact sets:
$`R_\alpha (E)`$ $`=`$ $`{\displaystyle \frac{\alpha }{2i\sqrt{E}}}{\displaystyle _{}}e^{2i\sqrt{E}t}f(t)𝑑t+O(E^1),`$
$`L_\alpha (E)`$ $`=`$ $`{\displaystyle \frac{\alpha }{2i\sqrt{E}}}{\displaystyle _{}}e^{2i\sqrt{E}t}f(t)𝑑t+O(E^1)`$
such that $`R_\alpha (E)=O(1/\sqrt{E})`$ and $`L_\alpha (E)=O(1/\sqrt{E})`$. If the single-site potential $`f`$ has $`p`$ derivatives in $`L^1()`$ then $`L_\alpha (E)=O(E^{(p+1)/2})`$ and $`R_\alpha (E)=O(E^{(p+1)/2})`$ as $`E\mathrm{}`$ . The estimate $`𝔼\{\xi _{\alpha (\omega )}(E)\}=O(1/\sqrt{E})`$ is Proposition 3 below. ∎
As an example we consider again the random Kronig-Penney model. The single-site spectral shift function is given in this case by
$$\xi _\alpha (E)=\frac{1}{\pi }\mathrm{arctan}\left(\frac{\alpha }{2\sqrt{E}}\right),E>0.$$
Therefore
$$𝔼\left\{\xi _{\alpha (\omega )}(E)\right\}=\frac{1}{\pi }_{}\mathrm{arctan}\left(\frac{\alpha }{2\sqrt{E}}\right)𝑑\kappa (\alpha )$$
and thus
$$\left|𝔼\left\{\xi _{\alpha (\omega )}(E)\right\}\right|\frac{|\alpha |}{2\pi \sqrt{E}}.$$
Using the explicit expression for the reflection amplitude one can easily show that
$$\frac{|\alpha |}{2\sqrt{E}}+\frac{\alpha ^2}{4E}𝔼\left\{\frac{|R_{\alpha (\omega )}(E)|}{1|R_{\alpha (\omega )}(E)|}\right\}\frac{|\alpha |}{2\sqrt{E}}+\frac{\alpha ^2}{2E}.$$
We complete this section with an estimate on $`𝔼\{\xi _{\alpha (\omega )}(E)\}`$ in the general case. We will prove
###### Proposition 3.
There is a constant $`c>0`$ independent of $`E`$, $`f`$, and $`\kappa `$ such that for all $`E>0`$
$$\left|𝔼\left\{\xi _{\alpha (\omega )}(E)\right\}\right|\frac{C}{2\sqrt{E}}𝔼\left\{|\alpha (\omega )|^{1/2}\right\}^2_{1/2}^{1/2}|f(x)|𝑑x.$$
Let $`l^{1/2}(L^1)`$ denote the Birman-Solomyak class of measurable functions $`V`$ for which
$$V_{l^{1/2}(L^1)}=\left[\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}\left(_{j1/2}^{j+1/2}|V(x)|𝑑x\right)^{1/2}\right]^2<\mathrm{}.$$
The claim of the proposition immediately follows from the following
###### Lemma 4.
Let $`Vl^{1/2}(L^1)`$. There is a constant $`c_1`$ independent of $`V`$ and $`E`$ such that
$$|\xi (E;H_0+V,H_0)|\frac{c_1}{2\sqrt{E}}V_{l^{1/2}(L^1)}$$
for all $`E>0`$.
###### Proof.
As proved in there is a constant $`c_2>0`$ independent of $`E`$ and $`V`$ such that
$$|\xi (E;H_0+V,H_0)|C_1V^{1/2}R_0(E+i0)|V|^{1/2}_{𝒥_1},$$
where $`V^{1/2}=\text{sign}V|V|^{1/2}`$, $`R_0(z)=(H_0z)^1`$, and $`_{𝒥_1}`$ denotes the trace class norm (see e.g. ). From the proof of Proposition 5.6 in it follows that
$$V^{1/2}R_0(E+i0)|V|^{1/2}_{𝒥_1}\frac{c_3}{\sqrt{E}}V_{l^{1/2}(L^1)}$$
for all $`E>0`$. ∎
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# 1 Introduction
## 1 Introduction
There is evidence that all M-theory or Type II string backgrounds of the form $`AdS_{p+2}\times X^{dp2}`$ in d-dimensions, where $`X^{dp2}`$ is an Einstein manifold, are dual to CFT’s in $`p+1`$ dimensions, living on the world-volume of p-branes . Many supergravity solutions associated with coset spaces $`X^{dp2}`$ are known and have been studied in the eighties. It is therefore interesting to identify the associated CFT’s and compare the KK spectrum with the spectrum of conformal operators. The dual CFT is realized on the world-volume of p-branes living at the singularities of $`C(X^{dp2})`$, the cone over $`X^{dp2}`$ . Unfortunately, there is no general method for determining the world-volume theory for branes in curved space-time, so one can only rely on geometrical intuition. Consider first the AdS<sub>5</sub> case. There are only two supersymmetric five-dimensional coset spaces, $`S^5`$ and $`T^{1,1}`$. $`S^5`$ is at the origin of the original Maldacena conjecture. The dual four-dimensional CFT dual to AdS$`{}_{5}{}^{}\times T^{1,1}`$ have been identified in . Many checks of this identification can be found in the literature . In the AdS<sub>4</sub> case, there is a richer zoo of seven-dimensional coset spaces, corresponding to supersymmetric backgrounds of M-theory . In , we proposed candidate dual CFT’s for the two $`𝒩=2`$ solutions, $`Q^{1,1,1}`$ and $`M^{1,1,1}`$ , using intuition from toric geometry. The KK spectrum and the properties of wrapped M5-branes associated to baryons nicely fit with the CFT expectations. A candidate dual for the $`𝒩=2`$ solution $`V^{5,2}`$, which does not admit toric description, has been proposed in . The purpose of this paper, which is the natural continuation of , is to discuss the $`𝒩=3`$ solution AdS$`{}_{4}{}^{}\times N^{0,1,0}`$.
$`N^{0,1,0}`$ can be written as $`\mathrm{SU}(3)/\mathrm{U}(1)`$ . It has isometry $`\mathrm{SU}(3)\times \mathrm{SU}(2)`$ and preserves $`𝒩=3`$ supersymmetry. Using geometrical arguments similar to those in , one is led to consider an $`𝒩=4`$ theory $`\mathrm{SU}(\mathrm{N})\times \mathrm{SU}(\mathrm{N})`$ with three hypermultiplets in the bi-fundamental representation of the two gauge groups. It was proposed in that the $`𝒩=3`$ CFT can be just obtained by adding an $`𝒩=3`$ preserving Chern-Simons term. We shall give evidence for this proposal by carefully comparing the observables in the CFT and the excitations of the supergravity background. The complete KK spectrum of M-theory on AdS$`{}_{4}{}^{}\times N^{0,1,0}`$ has been recently computed . Both KK and conformal field theory composite operators fall in representations of the superalgebra $`\mathrm{Osp}(3|4)`$ and can be conveniently described in terms of three-dimensional $`𝒩=3`$ superfields. In this paper, we first derive a general formalism for studying $`𝒩=3`$ superfields and the $`\mathrm{Osp}(3|4)`$ shortening conditions and we then apply it to the comparison between KK states and CFT composite operators. We shall exhibit the CFT supermultiplets of composite operators associated to all the short multiplets belonging to the KK spectrum <sup>1</sup><sup>1</sup>1One could also make an independent check of the dimension of supersingletons in the CFT by looking at the baryonic operators , which can be realized as wrapped M5-branes. Since such a calculation would simply be a repetition of known calculations that reveals no new feature we skip such additional check, which should be straightforward.. Indeed the analysis of the $`𝒩=3`$ solution reveals that all the general features which were common to the $`T^{1,1}`$, $`Q^{1,1,1}`$, $`M^{1,1,1}`$ and $`V^{5,2}`$ compactifications still hold true also for $`N^{0,1,0}`$. In particular, there are long multiplets with protected rational dimensions. We show that (in analogy with the other compactifications) many of them can be identified with CFT multiplets obtained by tensoring massless and short multiplets, as suggested in . We focus, in particular, on a very special long multiplet, which contains the volume of the internal manifold as one of the scalar components, and it is therefore universal for all compactification. In $`𝒩=2`$ compactification, the volume multiplet is a long vector multiplet. In $`𝒩=3`$ it becomes a long gravitino multiplet. In $`N^{0,1,0}`$, it has the right quantum numbers to be generated in a superHiggs mechanism, suggesting that the theory is a spontaneously broken phase of an $`𝒩=4`$ theory. This intriguing phenomenon will be investigated in a forthcoming publication .
The plan of this work is as follows. In section 2 we introduce the $`𝒩=3`$ superspace formalism and derive the shortening conditions of the $`\mathrm{Osp}(3|4)`$ irreducible representations in terms of differential constraints on primary conformal superfields. In section 3 we discuss the general structure of $`𝒩=3`$ three dimensional gauge theories using the component formalism and emphasizing the role of the Chern Simons interaction. In section 4 we identify the $`𝒩=3`$ gauge theory whose conformal fixed point realizes the AdS/CFT correspondence with the $`N^{0,1,0}`$ compactification of M-theory, while section 5 is devoted to test this correspondence. Finally section 6 briefly discusses, from a CFT point of view the long rational spin $`\frac{3}{2}`$ supermultiplet that suggests an interpretation in terms of superHiggs mechanism and that will be the focus of a forthcoming paper.
## 2 Three dimensional $`𝒩=3`$ superspace
In order to simplify the study of unitary irreducible representations of the $`\mathrm{Osp}(3|4)`$ superconformal algebra (see eq. (A.2) of ), we introduce a three dimensional $`𝒩=3`$ superspace formalism. This allows us to identify the short representations as particular constrained superfields. To this effect we introduce six Grassmann coordinates, $`\theta _\alpha ^i`$, transforming as three Majorana bispinors and as a triplet of the SO(3)<sub>R</sub> R-symmetry subalgebra:
$`[T^{ij},T^{kl}]=i(\delta ^{jk}T^{il}\delta ^{ik}T^{jl}\delta ^{jl}T^{ik}+\delta ^{il}T^{jk}),`$
$`T^{ij}=\theta _\alpha ^i{\displaystyle \frac{}{\theta _\alpha ^j}}\theta _\alpha ^j{\displaystyle \frac{}{\theta _\alpha ^i}}.`$ (2.1)
The other relevant generators have the following representation:
$`P_\mu =i_\mu ,`$
$`q^{\alpha i}={\displaystyle \frac{}{\theta _\alpha ^i}}+{\displaystyle \frac{1}{2}}/_\beta ^\alpha \theta ^{\beta i},`$
$`\{q^{\alpha i},q^{\beta j}\}=i\delta ^{ij}P/^{\alpha \beta },`$
$`[T^{ij},q^{\alpha k}]=i(\delta ^{jk}q^{\alpha i}\delta ^{ik}q^{\alpha j}).`$ (2.2)
We furthermore introduce the supercovariant derivatives:
$`𝒟^{\alpha i}={\displaystyle \frac{}{\theta _\alpha ^i}}{\displaystyle \frac{1}{2}}/_\beta ^\alpha \theta ^{\beta i},`$
$`\{𝒟^{\alpha i},𝒟^{\beta j}\}=i\delta ^{ij}P/^{\alpha \beta },\{q^{\alpha i},𝒟^{\beta j}\}=0,`$ (2.3)
in terms of which the shortening conditions can be expressed.
It is convenient to use the spherical irreducible basis of R-symmetry representations rather than the cartesian one, so that the Grassmann coordinates are renamed as in the following example:
$$\{\begin{array}{ccc}\theta ^+& =& \frac{1}{\sqrt{2}}(\theta ^1+i\theta ^2)\hfill \\ \theta ^{\mathrm{\hspace{0.17em}0}}& =& \theta ^3\hfill \\ \theta ^{}& =& \frac{1}{\sqrt{2}}(\theta ^1i\theta ^2)\hfill \end{array}$$
(2.4)
An $`𝒩=3`$ superfield $`\mathrm{\Theta }=\mathrm{\Theta }(x,\theta )`$ is a function of the bosonic $`x^\mu `$ and Grassmann coordinates $`\theta ^i`$, whose expansion in powers of $`\theta ^0`$ gives us the decomposition of the corresponding $`\mathrm{Osp}(3|4)`$ representation in $`𝒩=2`$ supermultiplets.
We are mainly interested in conformal primary superfields (see ), defined by:
$`\mathrm{\Theta }(x,\theta )=\mathrm{exp}\left[\text{i}x_\mu P^\mu +\theta ^iq^i\right]\mathrm{\Theta }(0),`$ (2.5)
where $`\mathrm{\Theta }(0)`$ is a primary field: $`[s_\alpha ^i,\mathrm{\Theta }(0)]=[K_\mu ,\mathrm{\Theta }(0)]=0`$. An irreducible representation of the superfield $`\mathrm{\Theta }(x,\theta )`$ is characterized by the Cartan labels of its highest weight state $`\mathrm{\Theta }(0)`$, namely its scaling dimension, its SO(1,2) Lorentz and SO(3)<sub>R</sub> R-symmetry quantum numbers. We denote the R-symmetry isospin with the suffix $`J`$ and the Lorentz character with a set of spinorial indices spanning the SO(1,2) irreducible representation.
### 2.1 Short $`\mathrm{Osp}(3|4)`$ representations as constrained superfields
In this section we analyze the differential constraints on the superfields that force their components to transform into short $`BPS`$ representations of the $`\mathrm{Osp}(3|4)`$ superalgebra. <sup>2</sup><sup>2</sup>2When this paper was nearly finished we learned of the recent work by Ferrara and Sokatchev that analyzes the differential constraints to be imposed, in harmonic superspace, on $`𝒩=8`$ superfields in order to describe short representations of the algebra $`\mathrm{Osp}(8|4)`$. Quite likely our results for $`\mathrm{Osp}(3|4)`$ can be described in that general formalism. A comparison is postponed to future investigations As in the best known case of $`𝒩=2`$ superfields, we find the existence of two kinds of constraints. The first one is a first order differential constraint given by:
$$𝒟_{\alpha _1}_{h.w.}\mathrm{\Theta }^{J(\alpha _1\mathrm{}\alpha _n)}(x,\theta )=0,$$
(2.6)
or
$$𝒟_\alpha _{h.w.}\mathrm{\Theta }^J(x,\theta )=0$$
(2.7)
for scalar superfields (without Lorentz indices). Here the tensor product refers to the SO(3)<sub>R</sub> isospin and “h.w” stands for highest weight. This means that only the SO(3)<sub>R</sub> highest weight part of the tensor product (2.6), between the isospin triplet $`𝒟`$ and the superfield $`\mathrm{\Theta }^J`$, is put to zero.
In complete analogy with the $`𝒩=2`$ case, we have another kind of constraint. It is a second order differential constraint, and it is allowed only for scalar (from the Lorentz viewpoint) superfields (of any isospin):
$$𝒟_\alpha _{h.w.}𝒟^\alpha _{h.w.}\mathrm{\Theta }^J=0.$$
(2.8)
The superconformal covariance of equations (2.6) and (2.8) poses some constraints on the conformal dimensions of the superfields.
#### 2.1.1 The $`𝒩=3`$ analogue of the chiral ring
Let us now analyze in more details the constraint (2.7) for the lowest Lorentz and isospin representations. The most interesting case to analyze is that of a scalar (from the Lorentz viewpoint) superfield. In this case the Lorentz character allows the existence of a ring structure which generalizes the chiral ring of $`𝒩=2`$ theories and seems to be a common feature shared by all the three dimensional superconformal field theories. The ring multiplicative operation is given by extracting the highest weight irreducible part from the ordinary (tensor) product of two short superfields of isospin $`J`$ and $`J^{}`$ respectively:
$$\mathrm{\Theta }^J\times \stackrel{~}{\mathrm{\Theta }}^J^{}=\left(\mathrm{\Theta }_{h.w.}\stackrel{~}{\mathrm{\Theta }}\right)^{J+J^{}}.$$
(2.9)
Indeed one can show that:
$$\begin{array}{c}𝒟^{J=1}_{h.w.}\mathrm{\Theta }^J=0\\ 𝒟^{J=1}_{h.w.}\stackrel{~}{\mathrm{\Theta }}^J^{}=0\end{array}\}𝒟^{J=1}_{h.w.}(\mathrm{\Theta }_{h.w.}\stackrel{~}{\mathrm{\Theta }})^{J+J^{}}=0.$$
(2.10)
The simplest case of short scalar superfield, apart from the trivial constant, is that of isospin $`J=1/2`$. In this case the shortening condition (2.7) reads:
$$\{\begin{array}{c}𝒟^+\mathrm{\Theta }^+=0,\\ \sqrt{2}𝒟^0\mathrm{\Theta }^++𝒟^+\mathrm{\Theta }^{}=0,\\ 𝒟^{}\mathrm{\Theta }^++\sqrt{2}𝒟^0\mathrm{\Theta }^{}=0,\\ 𝒟^{}\mathrm{\Theta }^{}=0.\end{array}$$
(2.11)
To make contact with $`𝒩=2`$ superspace formalism of it is useful to expand the most general form of $`\mathrm{\Theta }^{J=1/2}`$ in powers of $`\theta ^0`$. So we have:
$$\left(\begin{array}{c}\mathrm{\Theta }^+\\ \mathrm{\Theta }^{}\end{array}\right)=\left(\begin{array}{c}\mathrm{\Phi }_s\\ \mathrm{\Psi }_s^{}\end{array}\right)\frac{1}{\sqrt{2}}\left(\begin{array}{c}𝒟^+\mathrm{\Psi }_s^{}\\ 𝒟^{}\mathrm{\Phi }_s\end{array}\right)\theta ^0,$$
(2.12)
where $`\mathrm{\Phi }_s`$ and $`\mathrm{\Psi }_s`$ are two $`𝒩=2`$ supersingletons, namely they are two functions of $`x^\mu `$ and $`\theta ^\pm `$ fulfilling the constraints $`𝒟^+\mathrm{\Phi }_s=𝒟^{}𝒟^{}\mathrm{\Phi }_s=0`$. Hence we see that the direct generalization of the $`𝒩=2`$ supersingleton is the $`𝒩=3`$ short scalar superfield of minimum isospin. Let us now look at the case $`J=1`$, whose most general form is:
$$\left(\begin{array}{c}\mathrm{\Theta }^+\\ \mathrm{\Theta }^{\mathrm{\hspace{0.17em}0}}\\ \mathrm{\Theta }^{}\end{array}\right)=\left(\begin{array}{c}\mathrm{\Phi }\\ \mathrm{\Sigma }\\ \mathrm{\Psi }^{}\end{array}\right)\left(\begin{array}{c}𝒟^+\mathrm{\Sigma }\\ \frac{1}{2}(𝒟^{}\mathrm{\Phi }+𝒟^+\mathrm{\Psi }^{})\\ 𝒟^{}\mathrm{\Sigma }\end{array}\right)\theta ^0\frac{1}{8}\left(\begin{array}{c}𝒟^+𝒟^+\mathrm{\Psi }^{}\\ 2𝒟^+𝒟^{}\mathrm{\Sigma }\\ 𝒟^{}𝒟^{}\mathrm{\Phi }\end{array}\right)(\theta ^0\theta ^0),$$
(2.13)
where $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ are $`𝒩=2`$ chiral superfields and $`\mathrm{\Sigma }`$ is a linear superfield ($`𝒟^+𝒟^+\mathrm{\Sigma }=𝒟^{}𝒟^{}\mathrm{\Sigma }=0`$), which is a conserved (massless) vector current. Hence the superfield (2.13) represents the direct generalization of the $`𝒩=2`$ massless vector.
The $`𝒩=3`$ short scalar superfields of higher isospin can be obtained by multiplying smaller ones following the ring operation, i.e. by tensoring and taking the maximum isospin irreducible part. It is interesting to analyze the $`𝒩=2`$ field content, i.e. the single independent $`\theta ^0`$ components of such superfields. This gives an analytical version of the algebraic $`𝒩=3𝒩=2`$ decomposition of the short multiplets (see tables (5.1) and (5.2)). The first thing to note is that the shortening constraint (2.6) implies that the only independent components are the $`\theta ^0=0`$ restrictions of the $`𝒩=3`$ superfields. In the case of integer isospin we always obtain the same pattern:
$$\left(\begin{array}{c}\mathrm{\Phi }^1\mathrm{\Phi }^2\mathrm{}\mathrm{\Phi }^k\\ \mathrm{\Sigma }^1\mathrm{\Phi }^2\mathrm{}\mathrm{\Phi }^k+\mathrm{}+\mathrm{\Phi }^1\mathrm{}\mathrm{\Phi }^{(k1)}\mathrm{\Sigma }^k\\ \mathrm{}\\ \mathrm{\Sigma }^1\mathrm{\Psi }^2\mathrm{}\mathrm{\Psi }^k+\mathrm{}+\mathrm{\Psi }^1\mathrm{}\mathrm{\Psi }^{(k1)}\mathrm{\Sigma }^k\\ \mathrm{\Psi }^1\mathrm{\Psi }^2\mathrm{}\mathrm{\Psi }^k\end{array}\right)_{\theta ^0=0}^{J=k}\begin{array}{cc}& \mathrm{chiral}\\ & \mathrm{short}\mathrm{vector}\\ & 2k1\mathrm{long}\mathrm{vectors}\\ & \mathrm{short}\mathrm{vector}\\ & \mathrm{chiral}\end{array}$$
(2.14)
The half-integer isospin chiral superfields have a completely analogous structure. The only difference is that each field contains an odd number of $`𝒩=2`$ supersingletons. Thus the corresponding states are not observed in the Kaluza Klein spectrum of supergravity compactifications.
#### 2.1.2 The $`𝒩=3`$ short gravitinos
Let us now analyze the second order constraint (2.8), which yields the $`𝒩=3`$ short gravitinos. The lowest isospin case ($`J=0`$) corresponds to the massless gravitino superfield:
$$\mathrm{\Theta }=\mathrm{\Sigma }+G_\alpha \theta ^{0\alpha }+\frac{1}{4}(𝒟^+𝒟^{})\mathrm{\Sigma }(\theta ^0\theta ^0),$$
(2.15)
where $`G^\alpha `$ is an $`𝒩=2`$ massless gravitino ($`𝒟_\alpha ^\pm G^\alpha =0`$) and $`\mathrm{\Sigma }`$ a linear superfield, namely a massless vector.
Analogously we can derive the form of the most general $`J=1`$ short spinor superfield:
$`\left(\begin{array}{c}\mathrm{\Theta }^+\\ \mathrm{\Theta }^{\mathrm{\hspace{0.17em}0}}\\ \mathrm{\Theta }^{}\end{array}\right)=\left(\begin{array}{c}\mathrm{\Sigma }^+\\ \mathrm{\Sigma }^{\mathrm{\hspace{0.17em}0}}\\ \mathrm{\Sigma }^{}\end{array}\right)+\left(\begin{array}{c}G^+\theta ^0\\ G^{\mathrm{\hspace{0.17em}0}}\theta ^0\\ G^{}\theta ^0\end{array}\right)+\mathrm{derivative}\mathrm{terms},`$ (2.25)
which has six $`𝒩=2`$ independent components:
* two short gravitinos ($`𝒟^+G^+=𝒟^{}G^{}=0`$) ;
* one long gravitino, $`G^0`$ ;
* two short vectors: $`(𝒟^+𝒟^+)\mathrm{\Sigma }^+=(𝒟^{}𝒟^{})\mathrm{\Sigma }^{}=0`$ ;
* one long vector, $`\mathrm{\Sigma }^0`$ .
This $`𝒩=2`$ superfield content perfectly fits the algebraic decomposition of table (5.2). Short gravitinos of higher isospin can be obtained by composing the $`J=0`$ short gravitino with chiral superfields of any $`J`$. Obviously, even in this case, half-integer isospin gravitinos are not observed in the Kaluza Klein spectra, due to the presence of an odd number of supersingletons.
#### 2.1.3 The $`𝒩=3`$ short gravitons
The $`𝒩=3`$ short graviton multiplets are realized by spinor superfields fulfilling the first order constraint (2.6). Again, the massless case corresponds to the lowest ($`J=0`$) isospin superfield:
$$\mathrm{\Theta }^\alpha =G^\alpha +T^{(\alpha \beta )}\theta _\beta ^0\frac{1}{4}/_\beta ^\alpha G^\beta (\theta ^0\theta ^0),$$
(2.26)
where $`G^\alpha `$ is an $`𝒩=2`$ massless gravitino ($`𝒟_\alpha ^\pm G^\alpha =0`$) and $`T^{(\alpha \beta )}`$ is a massless graviton ($`𝒟_\alpha ^\pm T^{(\alpha \beta )}=0`$).
In an analogous way we can derive the form of the most general $`J=1`$ short spinor superfield:
$`\left(\begin{array}{c}\mathrm{\Theta }^{+\alpha }\\ \mathrm{\Theta }^{\mathrm{\hspace{0.17em}0}\alpha }\\ \mathrm{\Theta }^\alpha \end{array}\right)=\left(\begin{array}{c}G^{+\alpha }\\ G^{\mathrm{\hspace{0.17em}0}\alpha }\\ G^\alpha \end{array}\right)+\left(\begin{array}{c}T^{+(\alpha \beta )}\\ T^{\mathrm{\hspace{0.17em}0}(\alpha \beta )}\\ T^{(\alpha \beta )}\end{array}\right)\theta _\beta ^0+\mathrm{derivative}\mathrm{terms},`$ (2.36)
which has six $`𝒩=2`$ independent components (see table (5.2)):
* two short gravitons ($`𝒟_\alpha ^+T^{+(\alpha \beta )}=𝒟_\alpha ^{}T^{(\alpha \beta )}=0`$) ;
* one long graviton, $`T^{\mathrm{\hspace{0.17em}0}(\alpha \beta )}`$ ;
* two short gravitinos ($`𝒟^+G^+=𝒟^{}G^{}=0`$) ;
* one long gravitino, $`G^0`$ .
Short gravitons of higher isospin can be obtained by composing the $`J=0`$ massless graviton with chiral superfields of any $`J`$. Again, half-integer isospin gravitons, containing an odd number of supersingletons, are not observed in the Kaluza Klein spectra.
## 3 $`𝒩=3`$ gauge theory in three dimensions
In this section we discuss the structure of a three dimensional gauge theory with $`𝒩=3`$ supersymmetry. In paper we have already given the general form of an $`𝒩=2`$ three–dimensional gauge theory and the $`𝒩=3`$ case is just a particular case in that class since a theory with $`𝒩=3`$ SUSY, must a fortiori be an $`𝒩=2`$ theory. In we have also considered, within the $`𝒩=2`$ class, the case of $`𝒩=4`$ theories. These are obtained through dimensional reduction of an $`𝒩_4=2`$ theory in four–dimensions. Indeed since each $`D=4`$ Majorana spinor splits, under dimensional reduction on a circle $`𝕊^1`$, into two $`D=3`$ Majorana spinors, the number of three–dimensional supercharges is just twice the number of $`D=4`$ supercharges:
$$𝒩_3=2\times 𝒩_4$$
(3.1)
The $`𝒩_3=3`$ case corresponds to an intermediate situation. It is an $`𝒩_3=2`$ theory with the field content of an $`𝒩_3=4`$ one, but with additional $`𝒩_3=2`$ interactions that respect three out of the four supercharges obtained through dimensional reduction. Using an $`𝒩=2`$ superfield formalism and the notion of twisted chiral multiplets it was shown in that for abelian gauge theories these additional $`𝒩_3=3`$ interactions are
1. A Chern Simons term, with coefficient $`\alpha `$
2. A mass-term with coefficient $`\mu =\alpha `$ for the chiral field $`Y^I`$ in the adjoint of the color gauge group. By this latter we denote the complex field belonging, in four dimensions, to the $`𝒩_4=2`$ gauge vector multiplet.
In this section we want to retrieve the same result in the component formalism which is better suited to discus the relation between the world–volume gauge theory and the geometry of the transverse cone $`𝒞(X^7)`$. Then we dismiss superfields and turning to components we discuss the general form of a non abelian $`𝒩=3`$ gauge theory in three dimensions.
### 3.1 The field content and the interactions
Our strategy is that of writing the $`𝒩=3`$ gauge theory as aspecial case of an $`𝒩=2`$ theory, whose general form was derived in . For this latter the field content is given by:
$$\begin{array}{cccc}& & & \\ \text{multipl. type }/SO(1,2)\text{ spin}& 1& \frac{1}{2}& 0\\ & & & \\ & & & \\ & & & \\ \text{vector multipl.}& \underset{\text{gauge field}}{\underset{}{A_\mu ^I}}& \underset{\text{gauginos}}{\underset{}{(\lambda ^{+I},\lambda ^I)}}& \underset{\text{real scalar}}{\underset{}{M^I}}\\ & & & \\ & & & \\ \text{chiral multip.}& & \underset{\text{chiralinos}}{\underset{}{(\chi ^{+i},\chi ^i^{})}}& \underset{\text{complex scalars}}{\underset{}{z^i,\overline{z}^i^{}}}\end{array}$$
(3.2)
and without Fayet Iliopoulos terms, which do not exist in non abelian gauge theories with no U(1) factors, the complete Lagrangian has the following form:
$$^{𝒩=2}=^{kinetic}+^{fermionmass}+^{potential},$$
(3.3)
where
$`^{kinetic}`$ $`=`$ $`\{\eta _{ij^{}}_\mu z^i^\mu \overline{z}^j^{}\frac{1}{2}\eta _{ij^{}}(\chi ^j^{}/\chi ^{+i}+\chi ^{+i}/\chi ^j^{})`$ (3.4)
$`g_{IJ}F_{\mu \nu }^IF^{J\mu \nu }\alpha \left(g_{IJ}F_{\mu \nu }^IA_\rho ^J+f_{IJK}A_\mu ^IA_\nu ^JA_\rho ^K\right)ϵ^{\mu \nu \rho }`$
$`+\frac{1}{2}g_{IJ}_\mu M^I^\mu M^J\frac{1}{4}g_{IJ}(\lambda ^I/\lambda ^{+J}+\lambda ^{+I}/\lambda ^J)\}d^3x`$
$`^{fermionmass}`$ $`=`$ $`\{\frac{i}{2}(\chi ^{+i}_i_jW(z)\chi ^{+j}\chi ^i^{}_i^{}_j^{}\overline{W}(\overline{z})\chi ^j^{})`$ (3.5)
$`\frac{i}{2}f_{IJK}M^I\lambda ^J\lambda ^{+K}i\chi ^j^{}M^I(T_I)_{ij^{}}\chi ^{+i}`$
$`\left(\chi ^i^{}\lambda ^{+I}(T_I)_{i^{}j}z^j\chi ^{+i}\lambda ^I(T_I)_{ij^{}}\overline{z}^j^{}\right)`$
$`\frac{1}{2}\alpha g_{IJ}\lambda ^I\lambda ^{+J}\}d^3x`$
$`^{potential}`$ $`=`$ $`U(z,\overline{z})d^3x,`$ (3.6)
the scalar potential admitting the following general expression
$`U(z,\overline{z},M)`$ $`=`$ $`_iW(z)\eta ^{ij^{}}_j^{}\overline{W}(\overline{z})`$ (3.7)
$`+\frac{1}{2}g^{IJ}\left(\overline{z}^i^{}(T_I)_{i^{}j}z^j\right)\left(\overline{z}^k^{}(T_J)_{k^{}l}z^l\right)`$
$`+\overline{z}^i^{}M^I(T_I)_{i^{}j}\eta ^{jk^{}}M^J(T_J)_{k^{}l}z^l`$
$`2\alpha ^2g_{IJ}M^IM^J2\alpha M^I\left(\overline{z}^i^{}(T_I)_{i^{}j}z^j\right)`$
and the superpotential $`W(z)`$ being an arbitrary holomorphic function of the chiral scalars $`z^i`$. Our index notations and conventions are given in the appendices.
The $`𝒩=3`$ case is obtained when the following conditions are fulfilled:
* The spectrum of chiral multiplets and their representation assignments under the gauge and flavor groups are as follows
$$z^i=\{\begin{array}{cccc}\sqrt{2}Y^I\hfill & \text{adj}\left[𝒢_{gauge}\right]\hfill & \times & \text{id}\left[𝒢_{flavor}\right]\hfill \\ gu^a\hfill & 𝐑_g\left[𝒢_{gauge}\right]\hfill & \times & 𝐑_f\left[𝒢_{flavor}\right]\hfill \\ gv_a\hfill & \overline{𝐑}_g^1\left[𝒢_{gauge}\right]\hfill & \times & \overline{𝐑}_f^1\left[𝒢_{flavor}\right]\hfill \end{array}\eta ^{ik^{}}T_{k^{}j}^I=\{\begin{array}{c}if_{JK}^I\hfill \\ (T^I)_a^b\hfill \\ (\overline{T}^I)_a^b\hfill \end{array}$$
(3.8)
$`𝐑_g`$, and $`𝐑_f`$ being two complex representations of $`G_{gauge}`$ and $`G_{flavor}`$, respectively.
* The superpotential $`W(z)`$ has the following form:
$$W(Y,u,v)=g_{IJ}\left(2gY^Iv_aT_b^{J|a}u^b+\mathrm{\hspace{0.17em}2}\alpha Y^IY^J\right)$$
(3.9)
The reason why these two choices make the theory $`𝒩_3=3`$ invariant is simple: the first choice corresponds to assuming the field content of an $`𝒩_3=4`$ theory which is necessary since $`𝒩_3=3`$ and $`𝒩_3=4`$ supermultiplets are identical. The second choice introduces an interaction that preserves $`𝒩_3=3`$ supersymmetry but breaks (when $`\alpha 0`$) $`𝒩_3=4`$ supersymmetry. We can appreciate the last statement if we rewrite the Lagrangian in such a way that its invariance against the $`\mathrm{so}(3)_\mathrm{R}`$ R-symmetry becomes manifest. To this effect we begin by recalling that viewed from an $`𝒩_3=3`$ or $`𝒩_3=4`$ vantage point the chiral fields $`u^a,v_a`$ are the bosonic elements of a hypermultiplet and can be organized into a quaternion, according to the rule:
$$Q^a=\left(\begin{array}{cc}u^a& i\overline{v}^a\\ iv_a& \overline{u}_a\end{array}\right)q^{a|0}\text{ }\text{}+iq^{A|x}\sigma _x$$
(3.10)
In this way the transformation of the hypermultiplet $`u^a,v_a`$ under gauge or flavor generators can be rewritten as follows:
$`\delta ^I𝐐`$ $`=`$ $`i\widehat{T}^I𝐐`$
$`\delta ^I\left(\begin{array}{cc}u^a& i\overline{v}^a\\ iv_a& \overline{u}_a\end{array}\right)`$ $`=`$ $`i\left(\begin{array}{cc}T_b^{I|a}& \\ & \overline{T}_a^{Ib}\end{array}\right)\left(\begin{array}{cc}u^b& i\overline{v}^b\\ iv_b& \overline{u}_b\end{array}\right)`$ (3.17)
where the $`T_b^{I|a}`$ realize a representation of $`𝒢`$ in terms of $`n\times n`$ hermitian matrices. We define $`\overline{T}_a^{Ib}\left(T_b^{I|a}\right)^{}`$.
Under the SO(3)<sub>R</sub> R–symmetry the hypermultiplets transform as an SU(2) doublet, in the sense that for each $`𝒰`$ SU(2)$`{}_{R}{}^{}`$SO(3)<sub>R</sub> the quaternion varies as follows:
$$\delta _RQ^a=Q^a𝒰$$
(3.18)
On the other hand the auxiliary fields that appear in the gaugino’s supersymmetry transformation rules vary, under R–symmetry in the triplet representation of SO(3). Their on–shell values constitute the so called triholomorphic momentum map. This is a unimodular quaternion bilinear constructed by means of the gauge group generators. Explicitly one sets:
$$𝒫^I=\frac{1}{2}i\left(\overline{𝐐}\widehat{T}^I𝐐\right)=\left(\begin{array}{cc}𝒫_3^I& 𝒫_+^I\\ 𝒫_{}^I& 𝒫_3^I\end{array}\right)$$
(3.19)
where:
$`𝒫_3^I`$ $`=`$ $`\left(\overline{u}_aT_b^{I|a}u^b\overline{v}^a\overline{T}_a^{I|b}v_b\right)`$
$`𝒫_{}^I`$ $`=`$ $`2\delta _{ac}\overline{v}^cT_b^{I|a}u^b=2v_aT_b^{I|a}u^b`$
$`𝒫_+^I`$ $`=`$ $`\left(𝒫_{}^I\right)^{}=2\overline{v}^a\overline{T}_a^{I|b}\overline{u}_b`$ (3.20)
The first form of $`𝒫_{}^I`$ explicitly exhibits the SU(2) covariance in the sense that $`(u^a,\overline{v}^a)`$ is a doublet. The second expression will be interpreted later. Out of the triholomorphic momentum map we extract the three components of a real SO(3)<sub>R</sub> trivector. Explicitly we set:
$`𝒫_{\mathrm{}}^I`$ $``$ $`\{𝒫_1^I,𝒫_2^I,𝒫_3^I\}`$
$`𝒫_{}^I`$ $`=`$ $`i\left(𝒫_1^I+𝒫_2^I\right)`$
$`𝒫_+^I`$ $`=`$ $`i\left(𝒫_1^I𝒫_2^I\right)`$ (3.21)
There is another SO(3)<sub>R</sub> real trivector in the theory which is composed by the complex scalar field $`Y^I`$ in the adjoint representation of the gauge group together with the real scalar $`M^I`$ belonging to the $`𝒩=2`$ vector multiplet. Explicitly we set:
$$\varphi _{\mathrm{}}^I=\left(\begin{array}{c}\text{Im}Y^I\\ \text{Re}Y^I\\ \frac{1}{2}M^I\end{array}\right)$$
(3.22)
Inserting eq.(3.9) into the general $`𝒩_3=2`$ formula (3.7) and using the notations of eq.s(3.21,3.22) we can rewrite the final form of the $`𝒩=3`$ scalar potential in a way that exhibits manifest invariance under $`\mathrm{so}(3)`$ R-symmetry and is a sum of squares:
$`U`$ $`=`$ $`g_{IJ}\delta ^\mathrm{}m\left[2\sqrt{2}\alpha \varphi _{\mathrm{}}^I+\frac{1}{\sqrt{2}}g𝒫_{\mathrm{}}^I+𝒬_{\mathrm{}}^I\right]\left[2\sqrt{2}\alpha \varphi _m^J+\frac{1}{\sqrt{2}}g𝒫_m^J+𝒬_m^J\right]`$ (3.23)
$`+4g^2g_{IJ}\delta ^\mathrm{}m\varphi _{\mathrm{}}^I\varphi _m^J\left[\overline{u}_a(T^IT^J)_b^au^b+\overline{v}^a(\overline{T}^I\overline{T}^J)_a^bv_b\right]`$
where:
$$𝒬_{\mathrm{}}^I=\sqrt{2}ϵ_{\mathrm{}mn}\varphi _m^P\varphi _n^Qf_{PQ}^I$$
(3.24)
The classical vacua of the $`𝒩=3`$ theory are immediately determined from eq.(3.23). One has:
$`\varphi _{\mathrm{}}^I`$ $`=`$ $`0`$ (3.25)
$`𝒫_{\mathrm{}}^J(u,v)`$ $`=`$ $`0`$ (3.26)
Eq. (3.25) lifts the Coulomb branch of the theory setting to zero the vev.s of the scalar fields in the adjoint representation of the gauge group. Eq. (3.26), instead identifies the manifold of classical vacua with the HyperKähler quotient of the flat HyperKähler manifold spanned by the hypermultiplets $`u^a,\overline{v}^a`$ with respect to the triholomorphic action of the gauge group. The locus defined by (3.26) is the zero level set of the triholomorphic momentum map and it has to be further modded out by the action of U(1). When the generator $`T_b^{I|a}=\mathrm{i}\delta _b^a`$ is a U(1)–generator, eq.s (3.20) just reproduce the definition of the flag variety $`𝔽(1,2;3)`$ SU(3)/U(1)$`\times `$U(1) which is the base manifold of $`N^{0,1,0}`$ seen as a circle bundle (see ). This is what we explain in more details in the next section.
## 4 The $`𝒩=3`$ gauge theory corresponding to the $`N^{0,1,0}`$ compactification
Having clarified the structure of a generic $`𝒩=3`$ gauge theory let us consider the specific one associated with the $`N^{0,1,0}`$ seven–manifold. As explained in (see eq.(B.1) of that paper) the manifold $`N^{0,1,0}`$ is the circle bundle inside $`𝒪(1,1)`$ over the flag manifold $`𝔽(1,2;3)`$. In other words we have
$$N^{0,1,0}\stackrel{\pi }{}𝔽(1,2;3)$$
(4.1)
where, by definition,
$$𝔽(1,2;3)\frac{\mathrm{SU}(3)}{H_1\times H_2}$$
(4.2)
is the homogeneous space obtained by modding $`\mathrm{SU}(3)`$ with respect to its maximal torus:
$$H_1=\mathrm{exp}\left[i\theta _1\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 0\end{array}\right)\right];H_2=\mathrm{exp}\left[i\theta _2\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 2\end{array}\right)\right]$$
(4.3)
Furthermore as also explained in (see eq.(B.2)), the base manifold $`𝔽(1,2;3)`$ can be algebraically described as the following quadric
$$\underset{i=1}{\overset{3}{}}u^iv_i=0$$
(4.4)
in $`^2\times ^2`$, where $`u^i`$ and $`v_i`$ are the homogeneous coordinates of $`^2`$ and $`^2`$, respectively.
Hence a complete description of the metric cone $`𝒞\left(N^{0,1,0}\right)`$ can be given by writing the following equations in $`^3\times ^3`$:
$$𝒞\left(N^{0,1,0}\right)=\{\begin{array}{cccc}\hfill |u^i|^2|v_i|^2& =& 0\hfill & \text{fixes equal the radii of }^2\text{ and }^2\hfill \\ \hfill 2u^iv_i& =& 0\hfill & \text{cuts out the quadric locus}\hfill \\ \hfill (u^ie^{i\theta },v_ie^{i\theta })& & (u^i,v_i)\hfill & \text{identifies points of }\mathrm{U}(1)\text{ orbits}\hfill \end{array}$$
(4.5)
Eq.s (4.5) can be easily interpreted as the statement that the cone $`𝒞\left(N^{0,1,0}\right)`$ is the HyperKähler quotient of a flat three-dimensional quaternionic space with respect to the triholomorphic action of a $`\mathrm{U}(1)`$ group. Indeed the first two equations in (4.5) can be rewritten as the vanishing of the triholomorphic momentum map of a $`\mathrm{U}(1)`$ group. It suffices to identify:
$`𝒫_3`$ $`=`$ $`\left(|u^i|^2|v_i|^2\right)`$
$`𝒫_{}`$ $`=`$ $`2v_iu^i`$ (4.6)
Comparing with eq.s (3.25) we see that the cone $`𝒞(N^{0,1,0})`$ can be correctly interpreted as the space of classical vacua in an abelian $`𝒩=3`$ gauge theory with $`3`$ hypermultiplets in the fundamental representation of a flavor group $`\mathrm{SU}(3)`$.
This suggests that the $`𝒩=3`$ non–abelian gauge theory whose infrared conformal point is dual to supergravity compactified on AdS$`{}_{4}{}^{}\times N^{0,1,0}`$ has the following structure:
$$\begin{array}{cccc}\text{gauge group}& \hfill 𝒢_{gauge}& =& \mathrm{SU}(\mathrm{N})_1\times \mathrm{SU}(\mathrm{N})_2\hfill \\ & & & \\ \text{flavor group}& \hfill 𝒢_{flavor}& =& \mathrm{SU}(3)\hfill \\ & & & \\ \text{color representations of hypermultiplets}& \hfill \left[\begin{array}{c}u\\ v\end{array}\right]& & \left[\begin{array}{c}(𝐍_1,\overline{𝐍}_2)\\ (\overline{𝐍}_1,𝐍_2)\end{array}\right]\hfill \\ & & & \\ \text{flavor representations of hypermultiplets}& \hfill \left[\begin{array}{c}u\\ v\end{array}\right]& & \left[\begin{array}{c}(\mathrm{𝟑},\overline{\mathrm{𝟑}})\\ (\overline{\mathrm{𝟑}},\mathrm{𝟑})\end{array}\right]\hfill \end{array}$$
(4.7)
More explicitly and using an $`𝒩=2`$ notation we can say that the field content of our theory is given by the following chiral fields, that are all written as $`N\times N`$ matrices:
$$\begin{array}{cccc}Y_1& =& \left(Y_1\right)_{\mathrm{\Sigma }_1}^{\mathrm{\Lambda }_1}& \text{adjoint of }\mathrm{SU}(\mathrm{N})_1\\ Y_2& =& \left(Y_2\right)_{\mathrm{\Sigma }_2}^{\mathrm{\Lambda }_2}& \text{adjoint of }\mathrm{SU}(\mathrm{N})_2\text{ }\\ u^i& =& \left(u^i\right)_{\mathrm{\Sigma }_2}^{\mathrm{\Lambda }_1}& \text{in the }(\mathrm{𝟑},𝐍_1,\overline{𝐍}_2)\\ v_i& =& \left(v_i\right)_{\mathrm{\Sigma }_1}^{\mathrm{\Lambda }_2}& \text{in the }(\mathrm{𝟑},\overline{𝐍}_1,𝐍_2)\end{array}$$
(4.8)
and the superpotential can be written as follows:
$$W=2\left[g_1\text{Tr}\left(Y_1u^iv_i\right)+g_2\text{Tr}\left(Y_2v_iu_i\right)+\alpha _1\text{Tr}\left(Y_1Y_1\right)+\alpha _2\text{Tr}\left(Y_2Y_2\right)\right]$$
(4.9)
where $`g_{1,2},\alpha _{1,2}`$ are the gauge coupling constants and Chern Simons coefficients associated with the $`\mathrm{SU}(\mathrm{N})_{1,2}`$ simple gauge groups, respectively. Setting:
$`g_1`$ $`=`$ $`g_2=g`$
$`\alpha _1`$ $`=`$ $`\pm \alpha _2=\alpha `$ (4.10)
and integrating out the two fields $`Y_{1,2}`$ that have received a mass by the Chern Simons mechanism we obtain the effective quartic superpotential:
$$W^{eff}=\frac{1}{2}\frac{g^2}{\alpha }\left[\text{Tr}\left(v_iu^iv_ju^j\right)\pm \text{Tr}\left(u^iv_iu^ju_j\right)\right]$$
(4.11)
The vanishing relation one obtains from the above superpotential are the following ones:
$$u^iv_ju^j=\pm u^jv_ju^i;v_iu^jv_j=\pm v_ju^jv_i$$
(4.12)
Consider now the chiral conformal superfields one can write in this theory:
$$\mathrm{\Phi }_{j_1j_2\mathrm{}j_k}^{i_1i_2\mathrm{}i_k}\text{Tr}\left(u^{(i_1}v_{(j_1}u^{i_2}v_{j_2}\mathrm{}u^{i)_k}v_{j_k)}\right)$$
(4.13)
where the round brackets denote symmetrization on the indices. The above operators have $`k`$ indices in the fundamental representation of $`\mathrm{SU}(3)`$ and $`k`$ indices in the antifundamental one, but they are not yet assigned to the irreducible representation:
$$M_1=M_2=k$$
(4.14)
as it is predicted both by general geometric arguments and by the explicit evaluation of the Kaluza Klein spectrum of hypermultiplets . To be irreducible the operators (4.13) have to be traceless. This is what is implied by the vanishing relation (4.12) if we choose the minus sign in eq.(4.10).
Notice that for $`N^{0,1,0}`$ the form of the superpotential, which is dictated by the Chern-Simons term, is strongly reminiscent of the superpotential considered in . The CFT theory associated with $`N^{0,1,0}`$ has indeed many analogies with the simpler cousin $`T^{1,1}`$. There is however also a crucial difference. We recognize a general phenomenon that we already discussed in the $`M^{1,1,1}`$ and $`Q^{1,1,1}`$ compactifications . The moduli space of vacua of the abelian theory is isomorphic to the cone $`𝒞\left(N^{0,1,0}\right)`$. When the theory is promoted to an non-abelian one, there are naively conformal operators whose existence is in contradiction with geometric expectations and with the KK spectrum, in this case the hypermultiplets that do not satisfy relation (4.14). Differently from what happens for $`T^{1,1}`$ , the superpotential which can be added to the theory is not sufficient for eliminating these redundant non-abelian operators.
## 5 Tests of correspondence
In this section we present the basic checks of the correspondence between the $`𝒩=3`$ superconformal gauge theory just discussed and the $`N^{0,1,0}`$ compactification of M-theory. Here we identify the whole set of $`\mathrm{BPS}`$ composite operators dual to the short supermultiplets of the $`\mathrm{KK}`$ spectrum. In the next section we analyze the non-$`\mathrm{BPS}`$ composite operators dual to certain massive $`\mathrm{KK}`$ modes, which seem to be organized into the Higgs supermultiplet of a spontaneously broken $`𝒩=4`$ supergravity.
### 5.1 Comparison with the $`\mathrm{KK}`$ spectrum
Let us briefly summarize the $`\mathrm{KK}`$ spectrum of the AdS$`{}_{4}{}^{}\times N^{0,1,0}`$ compactification of $`11D`$ supergravity, organized into $`𝒩=3`$ supermultiplets . These are listed in table 5.1 and 5.2, where we give their decomposition in $`𝒩=2`$ supermultiplets and their flavor quantum numbers. The ultrashort multiplets are:
$$\begin{array}{ccc}& & \\ 𝒩=3\mathrm{multiplet}& \{\begin{array}{c}\mathrm{R}\mathrm{charge}\\ \mathrm{SU}(3)\mathrm{irrep}\end{array}& 𝒩=2\mathrm{multiplets}\\ & & \\ & & \\ \mathrm{massless}\mathrm{graviton}& \{\begin{array}{c}J=0\\ M_1=M_2=0\end{array}& \{\begin{array}{c}1\mathrm{massless}\mathrm{graviton}\\ 1\mathrm{massless}\mathrm{gravitino}\end{array}\\ & & \\ \mathrm{massless}\mathrm{vector}& \{\begin{array}{c}J=1\\ M_1=M_2=0\end{array}& \{\begin{array}{c}1\mathrm{massless}\mathrm{vector}\\ 2\mathrm{chiral}\mathrm{mult}.\end{array}\\ & & \\ \mathrm{massless}\mathrm{vector}& \{\begin{array}{c}J=1\\ M_1=M_2=1\end{array}& \{\begin{array}{c}1\mathrm{massless}\mathrm{vector}\\ 2\mathrm{chiral}\mathrm{mult}.\end{array}\end{array}$$
(5.1)
The short multiplets are:
$$\begin{array}{ccc}& & \\ 𝒩=3\mathrm{multiplet}& \{\begin{array}{c}\mathrm{R}\mathrm{charge}\\ \mathrm{SU}(3)\mathrm{irrep}\end{array}& 𝒩=2\mathrm{multiplets}\\ & & \\ & & \\ \mathrm{short}\mathrm{graviton}& \{\begin{array}{c}J=k1\\ M_1=M_2=k\end{array}& \{\begin{array}{cc}2& \mathrm{short}\mathrm{gravitons}\\ 2k1& \mathrm{long}\mathrm{gravitons}\\ 2& \mathrm{short}\mathrm{gravitinos}\\ 2k1& \mathrm{long}\mathrm{gravitinos}\end{array}\\ & & \\ \mathrm{short}\mathrm{gravitino}& \{\begin{array}{c}J=k+1,k0\\ M_1=k,M_2=k+3\end{array}& \{\begin{array}{cc}2& \mathrm{short}\mathrm{gravitinos}\\ 2k+1& \mathrm{long}\mathrm{gravitinos}\\ 2& \mathrm{short}\mathrm{vectors}\\ 2k+1& \mathrm{long}\mathrm{vectors}\end{array}\\ & & \\ \mathrm{short}\mathrm{vector}& \{\begin{array}{c}J=k,k2\\ M_1=M_2=k\end{array}& \{\begin{array}{cc}2& \mathrm{chiral}\mathrm{mult}.\\ 2& \mathrm{short}\mathrm{vectors}\\ 2k1& \mathrm{long}\mathrm{vectors}\end{array}\end{array}$$
(5.2)
#### 5.1.1 The fundamental supersingletons
In complete analogy with the $`𝒩=2`$ CFT’s analyzed in , the building blocks of all the superconformal primary fields are the supersingletons. In this case we have at our disposal the isospin doublet in the fundamental representation of the flavor group $`\mathrm{SU}(3)`$:
$$\mathrm{\Theta }^{iJ=1/2}=\left(\begin{array}{c}\mathrm{\Theta }^{+i}\\ \mathrm{\Theta }^i\end{array}\right)^{𝒩=3}=\left(\begin{array}{c}U^i\\ i\overline{V}^i\end{array}\right)^{𝒩=2}\frac{\sqrt{2}}{2}\theta ^0\left(\begin{array}{c}i𝒟^+\overline{V}^i\\ 𝒟^{}U^i\end{array}\right)^{𝒩=2},$$
(5.3)
and the conjugate doublet
$$\mathrm{\Theta }_j^{J=1/2}=\left(\begin{array}{c}\mathrm{\Theta }_j^+\\ \mathrm{\Theta }_j^{}\end{array}\right)^{𝒩=3}=\left(\begin{array}{c}iV_j\\ \overline{U}_j\end{array}\right)^{𝒩=2}\frac{\sqrt{2}}{2}\theta ^0\left(\begin{array}{c}𝒟^+\overline{U}_j\\ i𝒟^{}V_j\end{array}\right)^{𝒩=2}.$$
(5.4)
The $`𝒩=2`$ superfield components $`U^i`$ and $`V_i`$ are supersingletons:
$`U^i(x,\theta ^\pm )`$ $`=`$ $`u^i(x)+\theta ^+\chi _u^i(x)+\frac{1}{2}\theta ^+/u^i(x)\theta ^{},`$
$`V^i(x,\theta ^\pm )`$ $`=`$ $`v_i(x)+\theta ^{}\chi _{vi}^+(x)\frac{1}{2}\theta ^+/v_i(x)\theta ^{}.`$ (5.5)
The lowest components, the so-called Di’s, are the scalar fields $`u^i`$ and $`v_i`$ discussed in the previous section and realizing the homogeneous coordinates of $`^2\times ^2`$. Their color representations are given in (4.7) and is the same for the Rac’s $`\chi _u`$ and $`\chi _v`$.
#### 5.1.2 Field theory realization of the chiral ring
The generator of the chiral ring of our CFT is the highest weight part of the tensor product of (5.3) times its conjugate doublet, i.e. the $`\mathrm{SO}(3)_\mathrm{R}`$ triplet
$$\mathrm{\Theta }_j^{iJ=1}=Tr\left(\begin{array}{c}iU^iV_j\\ \frac{\sqrt{2}}{2}\left(U^i\overline{U}_j+\overline{V}^iV_j\right)\\ i\overline{V}^i\overline{U}_j\end{array}\right)^{𝒩=2}+𝒪(\theta ^0),$$
(5.6)
where the trace refers to the color indices. From the flavor viewpoint, we can extract the two irreducible pieces belonging to the symmetric tensor product of the $`\mathrm{𝟑}`$ and $`\overline{\mathrm{𝟑}}`$ of $`\mathrm{SU}(3)`$. They contain the two massless vectors in the Kaluza Klein spectrum (see table (5.1)):
* the adjoint,
$$\mathrm{\Sigma }_j^i\frac{\sqrt{2}}{2}Tr\left(U^i\overline{U}_j+\overline{V}^iV_j\right)\mathrm{flavor}\mathrm{trace}$$
(5.7)
corresponding to the conserved current of the global $`\mathrm{SU}(3)`$ flavor;
* the singlet,
$$\mathrm{\Sigma }\frac{\sqrt{2}}{2}Tr\left(U^i\overline{U}_i+\overline{V}^iV_i\right)$$
(5.8)
corresponding to the baryonic $`\mathrm{U}(1)`$ global symmetry.
By composing several massless vectors we obtain the whole chiral ring of superfields, containing $`𝒩=2`$ chiral fields and short vectors with the right flavor quantum numbers, as listed in table (5.1).
#### 5.1.3 Field theory realization of the short gravitinos
Let us come to the short gravitinos. Remember that we basically have at our disposal the $`𝒩=3`$ supersingleton of (5.3), which we will simply call $`\mathrm{\Theta }^i`$, and its conjugate $`\mathrm{\Theta }_j`$. Let us consider the following composite operator:
$$\mathrm{\Theta }^{(ijk)}=f^{lm(i}Tr\left[\mathrm{\Theta }^j\mathrm{\Theta }^{k)}\mathrm{\Theta }_l\mathrm{\Theta }_m\right],$$
(5.9)
where $`f^{ijk}`$ are the $`\mathrm{SU}(3)`$ structure constants and the round brackets mean symmetrization. From the isospin viewpoint, (5.9) is a triplet, while it transforms in the three time symmetric tensor product of the $`\mathrm{𝟑}`$ of $`\mathrm{SU}(3)`$, in agreement with the $`J=1`$ short gravitino of the $`N^{0,1,0}`$ Kaluza Klein spectrum (see table 5.2). By construction, the operator (5.9) is a short gravitino, namely it satisfies the second order differential constraint of eq. (2.8).
The $`𝒩=2`$ superfield content (see eq. 2.25) is given by:
$`\mathrm{\Sigma }^{+(ijk)}=f^{lm(i}U^jU^{k)}\left(V_l\overline{U}_mV_m\overline{U}_l\right);`$ (5.10)
$`\mathrm{\Sigma }^{\mathrm{\hspace{0.17em}0}(ijk)}=\sqrt{2}if^{lm(i}U^j\overline{V}^{k)}\left(V_l\overline{U}_mV_m\overline{U}_l\right);`$ (5.11)
$`\mathrm{\Sigma }^{(ijk)}=f^{lm(i}\overline{V}^j\overline{V}^{k)}\left(V_l\overline{U}_mV_m\overline{U}_l\right);`$ (5.12)
$`G_\alpha ^{+(ijk)}=f^{lm(i}U^jU^{k)}\left(\overline{U}_l𝒟_\alpha ^+\overline{U}_m\overline{U}_m𝒟_\alpha ^+\overline{U}_l\right);`$ (5.13)
$`G_\alpha ^{(ijk)}=f^{lm(i}\overline{V}^j\overline{V}^{k)}\left(V_l𝒟_\alpha ^{}V_mV_m𝒟_\alpha ^{}V_l\right).`$ (5.14)
The $`𝒩=3`$ short gravitinos of higher isospin are obtained by extracting the highest weight part from the product of operators in the chiral ring with (5.9):
$$\mathrm{\Theta }_{(j_1\mathrm{}j_{k1})}^{(i_1\mathrm{}i_{k1}klm)J=k}=Tr\left[\underset{k1\mathrm{objects}}{\underset{}{\mathrm{\Theta }_{j_1}^{i_1}_{h.w.}\mathrm{}_{h.w.}\mathrm{\Theta }_{j_{k1}}^{i_{k1}}}}_{h.w.}\mathrm{\Theta }^{(klm)}\right].$$
(5.15)
#### 5.1.4 Field theory realization of the short gravitons
Let us now consider the composite superfield:
$$\mathrm{\Theta }_\alpha ^{J=0}=Tr\left[\mathrm{\Theta }^i𝒟_\alpha \mathrm{\Theta }_i\mathrm{\Theta }_i𝒟_\alpha \mathrm{\Theta }^i\right],$$
(5.16)
where the scalar part is extracted from the isospin tensor product. It is straightforward to show that this superfield is a short graviton (2.26) by construction. From the $`𝒩=2`$ viewpoint, it is composed by:
* the massless graviton supermultiplet of the $`𝒩=2`$ subalgebra:
$`T_{(\alpha \beta )}=Tr[V/_{(\alpha \beta )}\overline{V}\overline{V}/_{(\alpha \beta )}VU/_{(\alpha \beta )}\overline{U}+\overline{U}/_{(\alpha \beta )}U`$
$`+2𝒟_{(\alpha }^{}U𝒟_{\beta )}^+\overline{U}+2𝒟_{(\alpha }^+\overline{V}𝒟_{\beta )}^{}V];`$ (5.17)
* the conserved current relative to the third supersymmetry charge, completing the $`𝒩=3`$ supersymmetry algebra:
$$G_\alpha =iTr\left[U𝒟_\alpha ^{}VV𝒟_\alpha ^+U+\overline{V}𝒟_\alpha ^+\overline{U}\overline{U}𝒟_\alpha ^+\overline{V}\right].$$
(5.18)
All together, $`T_{\alpha \beta }`$ and $`G_\alpha `$ constitute the supermultiplet containing the energy-momentum tensor, the $`𝒩=3`$ supersymmetry charges and the $`𝒩=3`$ R-symmetry currents.
Once again, the short gravitons of the $`CFT`$ are realized by composing (5.16) with the chiral ring operators and taking the highest weight part of isospin and flavor quantum numbers:
$$\mathrm{\Theta }_{\alpha (j_1\mathrm{}j_k)}^{(i_1\mathrm{}i_k)J=k}=Tr\left[\underset{k\mathrm{objects}}{\underset{}{\mathrm{\Theta }_{j_1}^{i_1}_{h.w.}\mathrm{}_{h.w.}\mathrm{\Theta }_{j_k}^{i_k}}}_{h.w.}\mathrm{\Theta }_\alpha ^{J=0}\right].$$
(5.19)
It is interesting to note that some of the $`𝒩=2`$ components of the short gravitons (5.19) of $`J1`$ (precisely, the second highest helicity states) are long gravitons with the following particular structure:
$$\mathrm{\Phi }\mathrm{conserved}\mathrm{vector}\mathrm{current}\times \mathrm{chiral}\mathrm{operator}\times \mathrm{stress}\mathrm{energy}\mathrm{tensor}.$$
(5.20)
These long $`𝒩=2`$ multiplets have nonetheless rational conformal dimensions, belonging to a short $`𝒩=3`$ graviton multiplet. Furthermore, they have the same structure of some long multiplets of rational conformal dimension identified in type IIB as well as in $`𝒩=2`$ (see eq. (6.64) of ) $`M`$-theory compactifications. This suggests that the existence of such rational multiplets in non-maximally supersymmetric AdS compactifications could be explained by the presence of a residual form of higher supersymmetry, possibly spontaneously broken. This explanation is confirmed by a second feature common to all the $`𝒩=3`$ AdS<sub>4</sub> compactifications of $`11D`$ supergravity: the presence of a superHiggs multiplet, that we discuss in the next section.
## 6 The Universal SuperHiggs multiplet
Finally we consider the CFT realization of a long gravitino multiplet that has integer conformal dimension:
$$E_0=3$$
(6.1)
and it is neutral with respect to the flavor group SU(3). It was found in the spectrum of the AdS$`{}_{4}{}^{}\times N^{0,1,0}`$ compactification but, as we shall argue in a forthcoming paper , it has a universal character, since it would appear with the same quantum numbers and the same conformal dimension (6.1) in any other Freund Rubin compactification of $`D=11`$ supergravity with $`𝒩=3`$ residual supersymmetry. In we shall discuss its interpretation as superHiggs multiplet in a partial supersymmetry breaking $`𝒩=4`$ to $`𝒩=3`$. Here we want to stress its universality also from the CFT point of view.
Consider the following scalar composite superfield:
$$𝒮=\mathrm{Tr}\left[\underset{J=0}{\underset{}{\mathrm{\Theta }_\mathrm{\Sigma }\mathrm{\Theta }_\mathrm{\Sigma }\mathrm{\Theta }_\mathrm{\Sigma }}}\right]=\mathrm{Tr}\left[\mathrm{\Theta }_\mathrm{\Sigma }^+\mathrm{\Theta }_\mathrm{\Sigma }^{\mathrm{\hspace{0.17em}0}}\mathrm{\Theta }_\mathrm{\Sigma }^{}\right],$$
(6.2)
where $`\mathrm{\Theta }_\mathrm{\Sigma }`$ is the field strength superfield, i.e. a real $`J=1`$ short superfield (see eq. 2.13) generalizing the linear multiplet of $`𝒩=2`$ gauge theories:
$`\mathrm{\Theta }_\mathrm{\Sigma }=\left(\begin{array}{c}Y\\ \mathrm{\Sigma }\\ Y^{}\end{array}\right)+𝒪(\theta ^0)=`$ (6.6)
$`\left(\begin{array}{c}Y+(\theta ^+\chi ^{})+(\theta ^+\theta ^+)H+\frac{1}{2}(\lambda ^+\theta ^0)+\frac{1}{2}(\theta ^0\theta ^+)P\hfill \\ \frac{1}{2}(\theta ^0\theta ^0)H^{}\frac{i}{2}(\theta ^0\gamma ^{\mu \nu }\theta ^+)F_{\mu \nu }\hfill \\ \\ M+\frac{1}{2}(\lambda ^+\theta ^{})+\frac{1}{2}(\lambda ^{}\theta ^+)+\frac{1}{2}(\theta ^+\theta ^{})P\frac{i}{2}(\theta ^{}\gamma ^{\mu \nu }\theta ^+)F_{\mu \nu }\hfill \\ +\frac{1}{2}(\theta ^0\chi ^{})\frac{1}{2}(\theta ^0\chi ^+)+(\theta ^0\theta ^+)H(\theta ^0\theta ^{})H^{}+\frac{1}{4}(\theta ^0\theta ^0)P\hfill \\ \\ Y^{}(\theta ^{}\chi ^+)(\theta ^{}\theta ^{})H^{}+\frac{1}{2}(\lambda ^{}\theta ^0)+\frac{1}{2}(\theta ^0\theta ^{})P\hfill \\ +\frac{1}{2}(\theta ^0\theta ^0)H\frac{i}{2}(\theta ^{}\gamma ^{\mu \nu }\theta ^0)F_{\mu \nu }\hfill \end{array}\right)+\begin{array}{c}\mathrm{derivative}\\ \mathrm{terms}.\end{array}`$
From eq. (6.2) it is possible to identify all the field components of the superHiggs multiplet, which turn out to be related, through the AdS/CFT correspondence, to certain Kaluza Klein modes of the $`N^{0,1,0}`$ compactification. Of particular interest, for its clear geometrical interpretation, is the scalar component of zero isospin and conformal dimension $`6`$. The corresponding supergravity state is given by the breathing mode, responsible for a uniform dilatation of the internal manifold $`X^7`$. It can be extracted by integrating the superfield $`𝒮`$ with respect to the Grassmann measure $`d^6\theta `$:
$$d^2\theta ^+d^2\theta ^{}d^2\theta ^{\mathrm{\hspace{0.17em}0}}𝒮=Tr\left[3iHH^{}P+\frac{1}{4}ϵ^{\lambda \mu \nu }ϵ^{\rho \sigma \tau }F_{\lambda \mu }F_{\nu \rho }F_{\sigma \tau }\right]+\mathrm{derivatives}.$$
(6.17)
The supergravity interpretation of this field as the volume mode of $`X^7`$ is a clear sign of the universality of the whole multiplet (it does not depend on any specific characteristic of the internal manifold). As we will show in this is true for all the components of the multiplet.
From the CFT viewpoint, the composite operator (6.17) is the $`𝒩=3`$ supersymmetrization of the following third power of the gauge field strength:
$$ϵ^{\lambda \mu \nu }ϵ^{\rho \sigma \tau }F_{\lambda \mu }F_{\nu \rho }F_{\sigma \tau },$$
whose dimension appears to be protected by some, so far unknown, non-renormalization theorem. Indeed a closely similar situation appears in type IIB AdS<sub>5</sub> compactifications, where the volume mode of the internal manifold is dual to the CFT operator $`F^4`$, of dimension 8, which is known to satisfy some non-renormalization theorem. This consideration suggests that the operator (6.17) could originate by the low energy expansion of an analogue of the Dirac Born Infeld for the $`M2`$-brane, as well as the operator $`F^4`$ comes from the $`\alpha ^{}`$ expansion of the DBI Lagrangian of the $`D3`$-brane. $`F^4`$ is indeed the operator that coupled to the background breathing mode on the D3-brane world-volume.
In this perspective, the universality of the (properly supersymmetrized) third power of $`F`$ could be understood: it should be traced back to the existence of a universal Lagrangian term for the $`M2`$-brane.
The explicit presence of $`F_{\mu \nu }`$ in the previous formulae deserves some comments. In three dimensions, the vector multiplet is not conformal and it does not make sense to consider it an elementary degree of freedom at the conformal point. The only singletons in three dimensions are hypermultiplets. Only hypermultiplets indeed appeared in the matching of the KK spectrum with the short multiplets of conformal operators that we discussed in the previous sections. The vector multiplet fields in the previous equations should be regarded as expressed in terms of the singletons at the conformal point, using the equations of motion, for example. Alternatively, we may consider the previous equations as operators in the three dimensional gauge theory that has the CFT as the IR limit. The previous discussion suggests that these operators become conformal operators at the fixed point.
## 7 Conclusions and perspectives
The identification and the study of conformal field theories dual to AdS supergravity compactifications is not a mere exercise of classification nor a simple test of AdS/CFT correspondence. As it is the case for the $`N^{0,1,0}`$ solution considered in this paper, a careful analysis of the properties of the theory, both on the CFT and on the supergravity side, may lead to surprising discoveries.
The most interesting lesson we have learned about non-maximally compactifications of $`M`$-theory regards the existence of some universal features which do not depend on the geometrical details of the compactification manifold, but only on the degree of supersymmetry of the solution.
From the supergravity viewpoint, we find that all the massless multiplets, related to symmetries of the theory, are always coupled to long shadow multiplets. Some of these can be interpreted as the massive (super)Higgs multiplets of some spontaneously broken (super)symmetry. This phenomenon is particularly interesting for the most general symmetries, such as the group of AdS space-time isometries and/or its supersymmetries. In the $`𝒩=3`$ case, for instance, the shadow multiplet of the massless graviton, related to the $`\mathrm{Osp}(3|4)`$ supergroup, is a massive gravitino multiplet with same quantum numbers of the first: it is a superHiggs multiplet. Hence every $`𝒩=3`$ solution of the form AdS$`{}_{4}{}^{}\times X^7`$, independently from $`X^7`$, turns out to be the broken phase of some not better specified $`𝒩=4`$ supergravity. The deepest implications of this fact are analyzed in .
Here we want to briefly discuss the consequences of the field theory counterpart of this phenomenon. The AdS/CFT prescriptions imply that the composite operators dual to the supergravity shadow multiplets have protected conformal dimensions. This fact is quit surprising because they are not organized in short multiplets, suggesting the existence of some non-trivial non-renormalization theorem, whose investigation is left to future speculations.
Another possible development is given by the $`M`$-brane interpretation of the CFT dual of the most universal shadow multiplet: the shadow of the stress-energy tensor, i.e. the breathing mode of the internal manifold. Its existence is independent even from the degree of supersymmetry of the theory, hence it must come from a universal term of the $`M2`$-brane action, not directly related to the background. The identification of such a term could shed new light on the microscopic structure of the $`M`$-theory.
##### Acknowledgements
We are grateful to Cesare Reina, Alessandro Tomasiello, Alessandro Zampa for many important discussions on the geometry of the $`N^{0,1,0}`$ manifold at the beginning of this work. We also acknowledge important exchanges of ideas with Sergio Ferrara and Riccardo D’Auria.
## Appendix A Conventions on spinors
In accordance with , spinor indices ($`\alpha ,\beta ,\gamma \mathrm{}`$) are contracted *from eight to two* and are raised and lowered with $`ϵ_{\alpha \beta }`$:
$`\psi ^\alpha ϵ^{\alpha \beta }\psi _\beta ϵ_{12}=ϵ^{21}=1`$
$`\psi _\alpha ϵ_{\alpha \beta }\psi ^\beta ϵ_{\alpha \gamma }ϵ^{\gamma \beta }=\delta _\alpha ^\beta `$
$`(\psi \varphi )\psi _\alpha \varphi ^\alpha =\psi ^\alpha \varphi _\alpha =\varphi _\alpha \psi ^\alpha (\varphi \psi ).`$ (A.1)
We choose the following representation of the $`SO(1,2)`$ Clifford algebra:
$`\{\begin{array}{ccc}\gamma ^{\mathrm{\hspace{0.17em}0}}& =& i\sigma ^2\\ \gamma ^1& =& \sigma ^3\\ \gamma ^2& =& \sigma ^1\end{array}\begin{array}{c}\eta _{\mu \nu }=\mathrm{diag}(++)\\ \gamma ^\mu \gamma _\beta ^{\mu |\alpha }\\ [\gamma ^\mu ,\gamma ^\nu ]=2ϵ^{\mu \nu \rho }\gamma _\rho ,\end{array}`$ (A.8)
hence the symmetry properties of the gamma matrices are:
$`\{\begin{array}{ccccc}\gamma ^{\mu |\alpha \beta }& & ϵ^{\beta \gamma }\gamma _\gamma ^{\mu |\alpha }& =& \gamma ^{\mu |\beta \alpha }\\ \gamma _{\alpha \beta }^\mu & & ϵ_{\alpha \gamma }\gamma _\beta ^{\mu |\gamma }& =& \gamma _{\beta \alpha }^\mu \end{array}`$ (A.11)
so that
$$(\psi \gamma ^\mu \varphi )=(\varphi \gamma ^\mu \psi ).$$
(A.12)
Complex conjugation acts as:
$$(\psi ^\alpha )^{}\overline{\psi }_\alpha ,$$
(A.13)
so that
$$(\psi \varphi )^{}=\overline{\varphi }\overline{\psi }=(\overline{\psi }\overline{\varphi }),$$
(A.14)
and, with our choice of gamma matrices,
$$(\psi \gamma ^\mu \varphi )^{}=(\overline{\varphi }\gamma ^\mu \overline{\psi })=(\overline{\psi }\gamma ^\mu \overline{\varphi }).$$
(A.15)
The spinorial derivatives act in the following way:
$$\frac{}{\theta _\alpha ^i}\theta _\beta ^j=\frac{}{\theta ^{i\beta }}\theta ^{j\alpha }=\delta _i^j\delta _\beta ^\alpha $$
(A.16)
and the supercovariant derivatives are:
$`𝒟^+\left({\displaystyle \frac{}{\theta ^{}}}+\frac{1}{2}/\theta ^+\right),`$
$`𝒟^{}\left({\displaystyle \frac{}{\theta ^+}}+\frac{1}{2}/\theta ^{}\right),`$
$`𝒟^{\mathrm{\hspace{0.17em}0}}\left({\displaystyle \frac{}{\theta ^0}}\frac{1}{2}/\theta ^{\mathrm{\hspace{0.17em}0}}\right).`$ (A.17)
## Appendix B Notes on the $`𝒩=2`$ superfields
Here we briefly review the differential constraint defining the $`𝒩=2`$ short superfield and their field decomposition, to fix the notations adopted in the paper.
* The chiral superfield
Identified by the constraint:
$$𝒟^+\mathrm{\Phi }(x,\theta ^\pm )=0.$$
(B.1)
In components is given by
$`\mathrm{\Phi }(x,\theta ^\pm )=z(x)+\theta ^+\chi ^{}(x)+(\theta ^+\theta ^+)H(x)`$
$`+\frac{1}{2}\theta ^+\gamma ^\mu \theta ^{}_\mu z(x)+\frac{1}{4}(\theta ^+\theta ^+)\theta ^{}\gamma ^\mu _\mu \chi ^{}(x)`$
$`+\frac{1}{16}(\theta ^+\theta ^+)(\theta ^{}\theta ^{})\mathrm{}z(x).`$ (B.2)
* The supersingleton
The $`𝒩=2`$ supersingleton is defined by
$$\{\begin{array}{ccc}𝒟^+\mathrm{\Phi }_s(x,\theta ^\pm )& =& 0\\ (𝒟^{}𝒟^{})\mathrm{\Phi }_s(x,\theta ^\pm )& =& 0.\end{array}$$
(B.3)
In components is given by
$`\mathrm{\Phi }_s(x,\theta ^\pm )=z(x)+\theta ^+\chi ^{}(x)+\frac{1}{2}\theta ^+\gamma ^\mu \theta ^{}_\mu z(x),`$ (B.4)
where $`z`$ and $`\chi ^{}`$ are on-shell massless fields:
$$\{\begin{array}{c}\mathrm{}z=0,\\ /\chi ^{}=0.\end{array}$$
(B.5)
* The short gravitino
The short gravitino, defined by:
$$𝒟_\alpha ^+G^{+\alpha }(x,\theta ^\pm )=0,$$
(B.6)
is given by
$`G^{+\alpha }(x,\theta ^\pm )=\lambda _L+A/^+\theta ^{}+A/^{}\theta ^++\varphi ^{}\theta ^+`$
$`+\lambda _T^+(\theta ^+\theta ^{})+\frac{1}{2}(\theta ^+\lambda _T^+)\theta ^{}+(\theta ^+\theta ^+)\lambda _T^{}`$
$`+(\theta ^+\gamma ^\mu \theta ^{})\psi _\mu +(\theta ^+\theta ^+)Z/\theta ^{}+\mathrm{derivative}\mathrm{terms}.`$ (B.7)
* The massless gravitino
The massless gravitino, defined by:
$$𝒟_\alpha ^+G^\alpha (x,\theta ^\pm )=𝒟_\alpha ^{}G^\alpha (x,\theta ^\pm )=0,$$
(B.8)
is given by
$`\mathrm{\Phi }^\alpha (x,\theta ^\pm )=\lambda _L+A/^+\theta ^{}+A/^{}\theta ^++(\theta ^+\gamma ^\mu \theta ^{})\psi _\mu +\mathrm{derivative}\mathrm{terms},`$ (B.9)
where the spinor $`\lambda _L`$ and the gravitino $`\psi _m`$ are massless:
$`/\lambda _L=ϵ^{\mu \nu \rho }\gamma _\mu _\nu \psi _\rho =0,`$ (B.10)
while the two vectors are in Lorentz gauge:
$$A^+=A^{}=0.$$
(B.11)
* The gauge potential superfield
Identified by the reality constraint
$$V^{}=V,$$
(B.12)
can be parametrized as:
$`V(x,\theta ^+,\theta ^{})=C(x)+\theta ^+\psi ^{}(x)+\theta ^{}\psi ^+(x)`$
$`+(\theta ^+\theta ^+)B(x)+(\theta ^{}\theta ^{})B^{}(x)`$
$`\frac{i}{2}\theta ^+\gamma ^\mu \theta ^{}A_\mu (x)+\frac{1}{2}(\theta ^+\theta ^{})M(x)`$
$`+\frac{1}{4}(\theta ^+\theta ^+)\theta ^{}\left[\lambda ^{}(x)+\gamma ^\mu _\mu \psi ^{}(x)\right]`$
$`+\frac{1}{4}(\theta ^{}\theta ^{})\theta ^+\left[\lambda ^+(x)+\gamma ^\mu _\mu \psi ^+(x)\right]`$
$`+\frac{1}{8}(\theta ^+\theta ^+)(\theta ^{}\theta ^{})\left[P(x)+\frac{1}{2}\mathrm{}C(x)\right].`$ (B.13)
The gauge transformation
$$VV+\mathrm{\Phi }+\mathrm{\Phi }^{},$$
(B.14)
corresponds to
$$\{\begin{array}{cc}CC+z+\overline{z}\hfill & PP\hfill \\ \psi ^\pm \psi ^\pm +\chi ^\pm \hfill & \lambda ^\pm \lambda ^\pm \hfill \\ BB+H\hfill & MM\hfill \\ A_\mu A_\mu +i(_\mu z_\mu \overline{z})\hfill & \end{array}$$
(B.15)
In Wess Zumino gauge, $`V`$ reduces to:
$`V(x,\theta ^+,\theta ^{})=\frac{i}{2}\theta ^+\gamma ^\mu \theta ^{}A_\mu (x)+\frac{1}{2}(\theta ^+\theta ^{})M(x)`$
$`+\frac{1}{4}(\theta ^+\theta ^+)\theta ^{}\lambda ^{}(x)+\frac{1}{4}(\theta ^{}\theta ^{})\theta ^+\lambda ^+(x)+\frac{1}{8}(\theta ^+\theta ^+)(\theta ^{}\theta ^{})P(x).`$ (B.16)
* The field strength
The gauge invariant super field strength is a real linear superfield:
$$𝒟^+𝒟^+\mathrm{\Sigma }=𝒟^{}𝒟^{}\mathrm{\Sigma }=0.$$
(B.17)
It is derived by the potential superfield $`V`$ by:
$`\mathrm{\Sigma }𝒟_\alpha ^+𝒟^\alpha V=𝒟_\alpha ^{}𝒟^{+\alpha }V=`$
$`M+\frac{1}{2}(\lambda ^+\theta ^{})+\frac{1}{2}(\lambda ^{}\theta ^+)+\frac{1}{2}(\theta ^{}\theta ^+)P\frac{i}{2}(\theta ^{}\gamma ^{\mu \nu }\theta ^+)F_{\mu \nu }`$
$`\frac{1}{8}(\theta ^{}\theta ^{})\theta ^+/\lambda ^+\frac{1}{8}(\theta ^+\theta ^+)\theta ^{}/\lambda ^{}+\frac{1}{16}(\theta ^+\theta ^+)(\theta ^{}\theta ^{})\mathrm{}M.`$ (B.18)
where
$$F_{\mu \nu }\frac{1}{2}\left(_\mu A_\nu _\nu A_\mu \right).$$
(B.19)
* The SYM and CS action
The abelian SYM action is:
$`4{\displaystyle d^3x\mathrm{\Sigma }^2}|_{(\theta ^+\theta ^+)(\theta ^{}\theta ^{})}=4{\displaystyle d^3xd^2\theta ^+d^2\theta ^{}\mathrm{\Sigma }^2}`$
$`={\displaystyle d^3x\left\{\frac{1}{4}\left(\lambda ^+/\lambda ^{}+\lambda ^{}/\lambda ^+\right)+\frac{1}{2}P^2\frac{1}{2}^\mu M_\mu MF_{\mu \nu }F^{\mu \nu }\right\}}.`$ (B.20)
The supersymmetric generalization of the Chern Simons term is:
$`4{\displaystyle d^3x\mathrm{\Sigma }V}|_{(\theta ^+\theta ^+)(\theta ^{}\theta ^{})}=4{\displaystyle d^3xd^2\theta ^+d^2\theta ^{}\mathrm{\Sigma }V}`$
$`={\displaystyle d^3x\left\{ϵ^{\mu \nu \rho }F_{\mu \nu }A_\rho \frac{1}{2}\lambda ^+\lambda ^{}PM\right\}}.`$ (B.21)
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# Folding Energetics in Thin-Film Diaphragms
## References and notes
1. Thin-film diaphragms have elicited much interest in recent years because of their applications in micro electro-mechanical systems; see almost any chapter in M. Madou, Fundamentals of Microfabrication (CRC Press, Boca Raton, Florida, 1997).
2. G. Gioia, G., and M. Ortiz, Adv. Appl. Mech. 33, 119 (1997).
3. A relatively high density ($`40`$kg/m<sup>3</sup>) is required for the foam to deform homogeneously. In fact, low-density solid foams display deformation patterns which are analogous to the folding patterns of anisotropically compressed diaphragms; see also .
4. A. Lobkovsky, S. Gentges, H. Li, D. Morse, and T. Witten, Science 270, 1482 (1995); E. Cerda, S. Chaieb, F. Melo, and L. Mahadevan, Nature 401, 46 (1999).
5. L. Modica, Arch. Rat. Mech. Anal. 98, 123 (1987).
6. R. V. Kohn, and S. Müller, Comm. Pure Appl. Math., 47, 405 (1994).
7. E. de Giorgi, Rendiconti di Matematica 8, 277 (1975).
8. P. Sternberg, Arch. Rat. Mech. Anal. 101, 209 (1988).
9. W. Jin, Singular Perturbation and the Energy of Folds, PhD Thesis, Courant Institute of Mathematical Sciences, New York University (1997).
10. J. M. Ball, and R. D. James, Arch. Ration. Mech. Anal. 100, 13 (1987).
11. R. V. Kohn, Continuum Mech. Thermodyn. 3, 193 (1991).
12. G. Gioia, A. DeSimone, and A. M. Cuitiño (to be published).
13. J. W. Cahn, Acta Metall. 9, 795 (1961).
14. A. G. Kachaturyan, Theory of Structural Transformations in Solids. (J. Willey & Sons, New York, 1983).
15. C. Orme, and B. G. Orr, Surf. Rev. Lett. 4, 71 (1997).
16. M. Ortiz, and E. A. Repetto, J. Mech. Phys. Solids 47, 397 (1999).
17. A. DeSimone, Arch. Ration. Mech. Anal. 125, 99 (1993).
18. J. Carr and R. Pego, Proc. Roy. Soc. London, A436, 569 (1992).
19. Y. Wang, G. Gioia, A. M. and Cuitiño (to be published).
20. M. Böltau, S. Walheim, J. Mlynek, G. Krausch, and U. Steiner, Nature 391, 877 (1998).
21. R. V. Kohn and S. Müller, Phil. Mag. A66, 697 (1992).
22. A. S. Argon, V. Gupta, H. S. Landis, and J. A. Cornie, J. Mater. Sci. 24, 1207 (1989).
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# 1 Sinh-Gordon model
## 1 Sinh-Gordon model
It seems that from the recent developments of the string theory there are some persistent requests about better understanding of the two-dimensional sigma-models with non-compact (in particular singular) target spaces. Physical properties of such models are expected to be quite different from these of the better studied compact sigma-models. It is therefore a challenge for the two-dimensional integrable field theory community to reveal the corresponding peculiarities and new features.
The two-dimensional sinh-Gordon model (ShG) is defined by the (Euclidean) action
$$A_{\mathrm{ShG}}=\left[\frac{1}{4\pi }(_a\varphi )^2+2\mu \mathrm{cosh}(2b\varphi )\right]d^2x$$
(1.1)
I believe that this model can be considered to say as “a model” of (suitably perturbed) non-compact sigma-models, in the sense that its properties are considerably different from those of the perturbed rational conformal field theories, these differences being sometimes quite similar to these between non-compact and compact sigma-models.
In (1.1) $`\mu `$ is a dimensional (of dimension $`\mu [\mathrm{mass}]^{2+2b^2}`$) coupling constant which determines the scale in the model and $`b`$ is the dimensionless ShG parameter. For the beginning I suppose it to be real and non-negative. The case $`b=0`$ turns to be somewhat singular for the subsequent considerations and here the study is restricted to positive values $`0<b<\mathrm{}`$ only. It is also convenient to use different parameters
$$Q=b+1/b$$
(1.2)
and
$$p=\frac{b}{Q}=\frac{b^2}{1+b^2}$$
(1.3)
instead of $`b`$.
The perturbing operators $`\mathrm{exp}(\pm 2b\varphi )`$ in (1.1) have negative dimension $`\mathrm{\Delta }=b^2`$. To make the coupling $`\mu `$ a strict sense we need to fix a normalization of these operators. Here they are implied to be normal ordered w.r.t. the massless (unperturbed) vacuum in the way that in unperturbed theory
$$e^{2b\phi (x_1)}\mathrm{}e^{2b\phi (x_n)}e^{2b\phi (y_1)}\mathrm{}e^{2b\phi (y_n)}_{\mu =0}=\frac{_{i,j}|x_iy_j|^{4b^2}}{_{i>j}(|x_ix_j||y_iy_j|)^{4b^2}}$$
(1.4)
The model is known to be integrable and can be solved for many important characteristics. In particular its factorized scattering theory is known since long . The spectrum consists of only one neutral particle $`B(\theta )`$ subject to a factorized scattering with two-particle amplitude
$$S(\theta )=\frac{\mathrm{sinh}\theta i\mathrm{sin}\pi p}{\mathrm{sinh}\theta +i\mathrm{sin}\pi p}$$
(1.5)
With the normalization (1.4) the mass $`m`$ of this particle is related to the scale parameter $`\mu `$ as
$$\pi \mu \frac{\mathrm{\Gamma }(b^2)}{\mathrm{\Gamma }(1b^2)}=\left[mZ(p)\right]^{2+2b^2}$$
(1.6)
where
$$Z(p)=\frac{1}{8\sqrt{\pi }}p^p(1p)^{1p}\mathrm{\Gamma }\left(\frac{1p}{2}\right)\mathrm{\Gamma }\left(\frac{p}{2}\right)$$
(1.7)
Notice that the scattering theory is invariant under the (week-strong coupling) duality transformation $`b1/b`$ which brings $`p1p`$. This means that the physical content of the model remains unchanged up to the overall mass scale. Since the combination (1.7) is again invariant under $`p1p`$, the mass scale also remains unchanged if the coupling constant $`\mu `$ is simultaneously substituted by the “dual” coupling constant $`\stackrel{~}{\mu }`$ related to $`\mu `$ as follows
$$\left(\pi \mu \frac{\mathrm{\Gamma }(b^2)}{\mathrm{\Gamma }(1b^2)}\right)^{1/b}=\left(\pi \stackrel{~}{\mu }\frac{\mathrm{\Gamma }(1/b^2)}{\mathrm{\Gamma }(11/b^2)}\right)^b$$
(1.8)
Therefore the sinh-Gordon model is completely invariant under duality $`b1/b`$; $`\mu \stackrel{~}{\mu }`$. Due to this symmetry it is sufficient to consider only the region $`0<b^21`$ or $`0<p1/2`$.
The infinite volume bulk vacuum energy of the model is also known exactly . In terms of the particle mass $`m`$ it is given by the following apparently self-dual expression
$$=\frac{m^2}{8\mathrm{sinh}\pi p}$$
(1.9)
## 2 TBA equation
Contrary to the on-mass-shell data of ShG quoted above, the off-mass-shell characteristics such as the correlation functions (with the exception of the vacuum expectation values of some local fields and their matrix elements between the asymptotic states ) are not known exactly.
Some progress can be made for the finite size effects where the problem is reduced to a non-linear integral equation known as the thermodynamic Bethe ansatz (TBA) one. Namely, consider the ground state energy $`E(R)`$ of the finite size ShG model placed on a circle of finite circumference $`R`$. In the TBA framework it appears as
$$E(R)=\frac{m}{2\pi }\mathrm{cosh}\theta \mathrm{log}\left(1+e^{\epsilon (\theta )}\right)d\theta $$
(2.1)
It is convenient to introduce also the $`R`$-dependent effective central charge $`c_{\text{eff}}(R)`$ as
$$c_{\text{eff}}(R)=\frac{6R}{\pi }E(R)$$
(2.2)
In (2.1) $`\epsilon (\theta )`$ is the solution of the TBA equation
$$mR\mathrm{cosh}\theta =\epsilon +\phi \mathrm{log}\left(1+e^{\epsilon (\theta )}\right)$$
(2.3)
where $``$ denotes the $`\theta `$-convolution. The kernel $`\phi (\theta )`$ is related to the ShG scattering data (1.5)
$$\phi (\theta )=\frac{i}{2\pi }\frac{d}{d\theta }\mathrm{log}S(\theta )=\frac{1}{2\pi }\frac{4\mathrm{sin}\pi p\mathrm{cosh}\theta }{\mathrm{cosh}2\theta \mathrm{cos}2\pi p}$$
(2.4)
The Fourier transform of the kernel reads
$$\phi (\omega )=e^{i\omega \theta }\phi (\theta )𝑑\theta =\frac{\mathrm{cosh}{\displaystyle \frac{a\pi \omega }{2}}}{\mathrm{cosh}{\displaystyle \frac{\pi \omega }{2}}}$$
(2.5)
where the parameter
$$a=12p$$
(2.6)
is simply reflected $`aa`$ under the dualily $`p1p`$ and therefore can be taken non-negative $`0a<1`$.
The following conclusion are readily made from the structure of the integral equation (2.3).
1. Function is even $`\epsilon (\theta )=\epsilon (\theta )`$ and analytic in the strip $`|Im\theta |<\pi /2\pi a/2`$. At $`Re\theta \mathrm{}`$ in this strip it has the asymptotic $`\epsilon (\theta )mRe^\theta /2`$. Therefore the function
$$Y(\theta )=\mathrm{exp}(\epsilon (\theta ))$$
(2.7)
is analyic and non-zero in this strip and at $`Re\theta \mathrm{}`$ behaves as
$$Y(\theta )\mathrm{exp}(mRe^\theta /2)$$
(2.8)
The asymptotic at $`Re\theta \mathrm{}`$ is related to (2.8) by the symmetry $`Y(\theta )=Y(\theta )`$. Let us define another even function
$$X(\theta )=\mathrm{exp}\left[\frac{mR}{2\mathrm{sin}\pi p}\mathrm{cosh}\theta +\frac{\mathrm{log}(1+Y(\theta ^{}))}{\mathrm{cosh}(\theta \theta ^{})}\frac{d\theta ^{}}{2\pi }\right]$$
(2.9)
which is obviously analytic and non-zero in the strip $`|Im\theta |<\pi /2`$ and at $`Re\theta \mathrm{}`$ in this strip
$$X(\theta )\mathrm{exp}\left(\frac{mR}{4\mathrm{sin}\pi p}\mathrm{exp}\theta \right)$$
(2.10)
From the TBA equation it follows that
$$X\left(\theta +ia\pi /2\right)X\left(\theta ia\pi /2\right)=Y(\theta )$$
(2.11)
2. In fact on the real axis $`Y(\theta )`$ is real and positive and therefore we expect a strip $`|Im\theta |<ϵ`$ with some finite $`ϵ>0`$ where $`1+Y(\theta )0`$. Therefore the analyticity condition for $`X(\theta )`$ can be extended to the strip $`|Im\theta |<\pi /2+ϵ`$. This is enough to prove the relation
$$X\left(\theta +i\pi /2\right)X\left(\theta i\pi /2\right)=1+Y(\theta )$$
(2.12)
The functional equation
$$X\left(\theta +i\pi /2\right)X\left(\theta i\pi /2\right)=1+X\left(\theta +ia\pi /2\right)X\left(\theta ia\pi /2\right)$$
(2.13)
follows. This relation allows to extend the original analyticity strip $`|Im\theta |<\pi /2+ϵ`$ to the strip $`|Im\theta |<3\pi /2`$ and, as we’ll see before long, to the whole comlex plane of $`\theta `$ so that $`X(\theta )`$ is an entire function of $`\theta `$. Notice that from (2.13) it follows that the asymptotic (2.10) holds in the larger strip $`|Im\theta |<\pi `$. The asymptotic outside this strip is more complicated.
3. As a consequence, $`Y(\theta )`$ is also entire function of $`\theta `$ and satisfyes the following functional relation
$$Y\left(\theta +i\pi /2\right)Y\left(\theta i\pi /2\right)=\left(1+Y\left(\theta +ia\pi /2\right)\right)\left(1+Y\left(\theta ia\pi /2\right)\right)$$
(2.14)
The last equation is very similar to the functional relations appearing in the TBA study of the integrable perturbed rational conformal field theories (the so called $`Y`$-systems). Typically such $`Y`$-systems imply a periodicity of the $`Y`$-functions in $`\theta `$ with some imaginary period related to the scale dimension $`\mathrm{\Delta }`$ of the perturbing operator (see e.g., ). This periodicity in order entails special “perturbative” structure of the short distance $`R0`$ behavior of the ground state energy $`E(R)`$. Namely, up to one exceptional term, it is a regular expansion in powers of $`R^{22\mathrm{\Delta }}`$
$$E(R)=_{\mathrm{vac}}R\frac{\pi }{6R}\underset{n=0}{\overset{\mathrm{}}{}}c_nR^{(22\mathrm{\Delta })n};R0$$
(2.15)
where $`_{\mathrm{vac}}`$ is the infinite volume vacuum energy of the model. Unlike this typical situation, the $`Y`$-system (2.14) or the $`X`$-system (2.13) does not imply any apparent periodic structure of $`Y(\theta )`$ in $`\theta `$. As a manifestation of this peculiarity, the $`R0`$ behavior of $`E(R)`$ is different from (2.15) and includes softer logarithmic corrections
$$E(R)=R\frac{\pi }{6R}\left(1\frac{3\pi ^2p(1p)}{2\mathrm{log}^2R}+O\left(\mathrm{log}^3R\right)\right);R0$$
(2.16)
The purpose of the next two sections is to reveal two hidden periodic structures (with different periods) of the $`Y`$-system (2.14).
## 3 Discrete Liouville equation
In this section I discuss the following two-dimensional non-linear finite-difference equation for the function $`X(u,v)`$
$$X(u+1,v)X(u1,v)=1+X(u,v+1)X(u,v1)$$
(3.1)
which is apparently resemblent of the functional $`X`$-system (2.13). In the next section we’ll see how some of the results for (3.1) can be specialized to our TBA problem. Equation (3.1) is a particular case of Hirota difference equation . The constructions of this section can be found in (see also ) where more general difference system is analysed. They appeared also in a quite close context in .
Equation (3.1) can be considered as a discretisation of the hyperbolic Liouville equation
$$_u^2\phi _v^2\phi =e^{2\phi }$$
(3.2)
Indeed, let $`X(u,v)=\mathrm{exp}(\phi (u,v))`$ and let $`\phi (u,v)`$ be large and negative. Then eq.(3.1) is approximated by
$$\phi (u+1,v)+\phi (u1,v)\phi (u,v+1)\phi (u,v1)=\mathrm{exp}(\phi (u,v+1)+\phi (u,v1))$$
(3.3)
In the long-wave limit this is reduced to (3.2). As we’ll see below, eq.(3.1) is in many respects very similar to the Liouville equation. It seems quite natural to call it the discrete Liouville equation .
Let me remind well known construction of a local solution to eq.(3.2).
1a. It is convenient to use the light cone variables $`x^+=u+v`$ and $`x^{}=uv`$, so that $`_+=(_u+_v)/2`$; $`_{}=(_u_v)/2`$ and (3.2) reads
$$4_+_{}\phi =e^{2\phi }$$
(3.4)
Let $`\phi `$ be a local solution of (3.4). Define
$`t`$ $`=(_+\phi )^2+_+^2\phi `$ (3.5)
$`\stackrel{~}{t}`$ $`=(_{}\phi )^2_{}^2\phi `$
As a consequence of eq.(3.4) we have
$$_{}t=_+\stackrel{~}{t}=0$$
(3.6)
so that $`t(x^+)`$ and $`\stackrel{~}{t}(x^{})`$ are respectively functions of only $`x^+`$ and $`x^{}`$.
2a. Field $`X=\mathrm{exp}(\phi )`$ satisfies two linear differential equations
$$(_+^2+t(x^+))X(u,v)=0;(_{}^2+\stackrel{~}{t}(x^{}))X(u,v)=0$$
(3.7)
3a. Let $`Q_\pm (x)`$ and $`\stackrel{~}{Q}_\pm (x)`$ be linearly independent solutions to the ordinary differential equations
$`(_x^2+t(x))Q_\pm (x)`$ $`=0`$ (3.8)
$`(_x^2+\stackrel{~}{t}(x))\stackrel{~}{Q}_\pm (x)`$ $`=0`$
normalized in the way that
$`_xQ_+(x)Q_{}(x)Q_+(x)_xQ_{}(x)`$ $`=1/2`$ (3.9)
$`_x\stackrel{~}{Q}_+(x)\stackrel{~}{Q}_{}(x)\stackrel{~}{Q}_+(x)_x\stackrel{~}{Q}_{}(x)`$ $`=1/2`$
A local solution of (3.4) can be constructed as
$$\mathrm{exp}(\phi (u,v))=Q_+(x^+)\stackrel{~}{Q}_+(x^{})+Q_{}(x^+)\stackrel{~}{Q}_{}(x^{})$$
(3.10)
4a. Introduce the functions
$$F(x)=\frac{Q_+(x)}{Q_{}(x)};G(x)=\frac{\stackrel{~}{Q}_{}(x)}{\stackrel{~}{Q}_+(x)}$$
(3.11)
The solution (3.10) can be rewritten as
$$\mathrm{exp}2\phi (u,v)=\frac{4F^{}(x^+)G^{}(x^{})}{(F(x^+)G(x^{}))^2}$$
(3.12)
Notice also that in terms of $`F`$ and $`G`$
$$t(x)=2\{F(x),x\};\stackrel{~}{t}(x)=2\{G(x),x\}$$
(3.13)
where
$$\{f(x),x\}=\frac{f^{\prime \prime \prime }}{f^{}}\frac{3}{2}\left(\frac{f^{\prime \prime }}{f^{}}\right)^2$$
(3.14)
is the Schwarz derivative.
Now let us turn to the discrete Liouville equation (3.1). Defining, similarly to eq.(2.11),
$`Y(u,v)`$ $`=X(u,v+1)X(u,v1)`$ (3.15)
$`1+Y(u,v)`$ $`=X(u+1,v)X(u1,v)`$
we arrive at the finite difference equation analogous to (2.14)
$$Y(u+1,v)Y(u1,v)(1+Y(u,v+1))(1+Y(u,v1))$$
(3.16)
1b. Introduce
$`T(u,v)`$ $`={\displaystyle \frac{X(u+1,v+1)+X(u1,v1)}{X(u,v)}}`$ (3.17)
$`\stackrel{~}{T}(u,v)`$ $`={\displaystyle \frac{X(u+1,v1)+X(u1,v+1)}{X(u,v)}}`$
As a consequence of eq.(3.1) we obtain
$`T(u+1,v1)`$ $`=T(u,v)`$ (3.18)
$`\stackrel{~}{T}(u+1,v+1)`$ $`=\stackrel{~}{T}(u,v)`$
so that $`T=T(u+v)`$ and $`\stackrel{~}{T}=\stackrel{~}{T}(uv)`$. In the continuous limit these objects are related to (3.5) as $`T(u)=24t(u)+\mathrm{}`$; $`\stackrel{~}{T}(u)=24\stackrel{~}{t}(u)+\mathrm{}`$ .
2b. Eqs.(3.17) can be rewritten as (similarly to (3.7))
$`X(u+1,v+1)+X(u1,v1)`$ $`=T(u+v)X(u,v)`$ (3.19)
$`X(u+1,v1)+X(u1,v+1)`$ $`=\stackrel{~}{T}(uv)X(u,v)`$
3b. Let $`Q_\pm (u)`$ and $`\stackrel{~}{Q}_\pm (u)`$ be linearly independent solutions of the second order finite difference equations
$`Q_\pm (u+2)+Q_\pm (u2)`$ $`=T(u)Q_\pm (u)`$ (3.20)
$`\stackrel{~}{Q}_\pm (u+2)+\stackrel{~}{Q}_\pm (u2)`$ $`=\stackrel{~}{T}(u)\stackrel{~}{Q}_\pm (u)`$
normalized by the “quantum Wronskians”
$`Q_+(u+1)Q_{}(u1)Q_+(u1)Q_{}(u+1)`$ $`=1`$ (3.21)
$`\stackrel{~}{Q}_+(u+1)\stackrel{~}{Q}_{}(u1)\stackrel{~}{Q}_+(u1)\stackrel{~}{Q}_{}(u+1)`$ $`=1`$
Then it is verified that
$$X(u,v)=Q_+(u+v)\stackrel{~}{Q}_+(uv)+Q_{}(u+v)\stackrel{~}{Q}_{}(uv)$$
(3.22)
is a local solution of the discrete Liouville equaiton (3.1).
4b. Introduce the functions
$$F(u)=\frac{Q_+(u)}{Q_{}(u)};G(u)=\frac{\stackrel{~}{Q}_{}(u)}{\stackrel{~}{Q}_+(u)}$$
(3.23)
which can be used to present the local solution (3.22) in the form
$`Y(u,v)`$ $`={\displaystyle \frac{(F(u+v+1)G(uv1))(F(u+v1)G(uv+1))}{(F(u+v+1)F(u+v1))(G(uv+1)G(uv1))}}`$ (3.24)
$`1+Y(u,v)`$ $`={\displaystyle \frac{(F(u+v+1)G(uv+1))(F(u+v1)G(uv1))}{(F(u+v+1)F(u+v1))(G(uv+1)G(uv1))}}`$
Let me mention also the discrete analogue of the Schwarzian derivative (3.13)
$$T(u+1)T(u1)=\frac{(F(u+3)F(u1))(F(u+1)F(u3))}{(F(u+3)F(u+1))(F(u1)F(u3))}$$
(3.25)
## 4 Application to TBA
The above constructions for the discrete Liouville equation can be translated for the ShG $`X`$-system (2.13) we are interested in. Let us require the following periodicity condition for $`X(u,v)`$ in (3.1)
$$X(u+a,v+1)=X(u,v)$$
(4.1)
with some parameter $`a`$ (at this point we start to diverge from the lines of ). With this periodicity eq.(3.1) reads
$$X(u+1,v)X(u1,v)=1+X(u+a,v)X(ua,v)$$
(4.2)
Here $`v`$ can be considered as a parameter. Suppressing this redundant dependence and rescaling $`u`$ as
$$\theta =i\pi u/2$$
(4.3)
we are back to the ShG $`X`$ system in the form (2.13).
From $`X(\theta )`$ the two functions $`T(\theta )`$ and $`\stackrel{~}{T}(\theta )`$ are readily restored as
$`T(\theta )`$ $`={\displaystyle \frac{X\left(\theta +i\pi (1a)/2\right)+X\left(\theta i\pi (1a)/2\right)}{X(\theta )}}`$ (4.4)
$`\stackrel{~}{T}(\theta )`$ $`={\displaystyle \frac{X\left(\theta +i\pi (1+a)/2\right)+X\left(\theta i\pi (1+a)/2\right)}{X(\theta )}}`$
The “holomorphic” property (3.18) is translated to the following periodicity of these functions
$`T\left(\theta +i\pi (1+a)/2\right)`$ $`=T(\theta )`$ (4.5)
$`\stackrel{~}{T}\left(\theta +i\pi (1a)/2\right)`$ $`=\stackrel{~}{T}(\theta )`$
Notice that the period $`i\pi /(1+b^2)`$ of $`T`$ corresponds to the negative dimension $`\mathrm{\Delta }=b^2`$ of the perturbing operator in (1.1). As it can be anticipated from the self-duality of ShG, the second period $`i\pi b^2/(1+b^2)`$ of $`\stackrel{~}{T}`$ is related to the dimension $`\stackrel{~}{\mathrm{\Delta }}=b^2`$ of the “dual” exponentials $`\mathrm{exp}(\pm 2\varphi /b)`$.
As it is discussed in sect.2 $`X(\theta )`$ is analytic and non-zero in the strip $`|Im\theta |<\pi /2`$ and analytic in the larger strip $`|Im\theta |<3\pi /2`$. Therefore $`T(\theta )`$ and $`\stackrel{~}{T}(\theta )`$ are analytic in the strip $`|Im\theta |<\pi /2`$ and by periodicity (4.5) are entire functions of $`\theta `$. It follows from (4.4) that $`X(\theta )`$ is an entire function of $`\theta `$ too.
The asymptotics at $`Re\theta \mathrm{}`$ follow from (2.10) and (4.4)
$`T(\theta )`$ $`\mathrm{exp}\left({\displaystyle \frac{mR\mathrm{exp}(\theta i\pi (1p)/2)}{4\mathrm{cos}(\pi p/2)}}\right)\text{in the strip }0<Im\theta <\pi (1+a)/2`$ (4.6)
$`\stackrel{~}{T}(\theta )`$ $`\mathrm{exp}\left({\displaystyle \frac{mR\mathrm{exp}(\theta i\pi p/2)}{4\mathrm{sin}(\pi p/2)}}\right)\text{in the strip }0<Im\theta <\pi (1a)/2`$
The real axis $`Im\theta =0`$ is a Stokes line and here
$`T(\theta )`$ $`2\mathrm{exp}\left({\displaystyle \frac{mR}{4}}\mathrm{tan}(\pi p/2)e^\theta \right)\mathrm{cos}\left({\displaystyle \frac{mR}{4}}\mathrm{exp}\theta \right)`$ (4.7)
$`\stackrel{~}{T}(\theta )`$ $`2\mathrm{exp}\left({\displaystyle \frac{mR}{4}}\mathrm{cot}(\pi p/2)e^\theta \right)\mathrm{cos}\left({\displaystyle \frac{mR}{4}}\mathrm{exp}\theta \right)`$
The $`Re\theta \mathrm{}`$ asymptotic in the whole plane of $`\theta `$ is restored from the periodicity (4.5). Following (4.7) both $`T`$ and $`\stackrel{~}{T}`$ have infinite number of zeroes on the real axis located at $`\theta =\pm \theta _n`$, $`n=1,2,\mathrm{},\mathrm{}`$ with $`\theta _n\mathrm{log}(2\pi n/mR)+O(1/n)`$ at $`n\mathrm{}`$. The half-period shifted functions $`T(\theta +i\pi (1p)/2)`$ and $`\stackrel{~}{T}(\theta +i\pi p/2)`$ are also real at real $`\theta `$ and at $`\theta \mathrm{}`$ behave as
$`T(\theta +i\pi (1p)/2)`$ $`\mathrm{exp}\left({\displaystyle \frac{mR}{4\mathrm{cos}(\pi p/2)}}e^\theta \right)`$ (4.8)
$`\stackrel{~}{T}(\theta +i\pi p/2)`$ $`\mathrm{exp}\left({\displaystyle \frac{mR}{4\mathrm{sin}(\pi p/2)}}e^\theta \right)`$
Analytic properties of $`T`$ and $`\stackrel{~}{T}`$ allow the following convergent expansions
$`T(\theta )`$ $`={\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}T_n\mathrm{exp}(2nQb\theta )`$ (4.9)
$`\stackrel{~}{T}(\theta )`$ $`={\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}\stackrel{~}{T}_n\mathrm{exp}(2nQ\theta /b)`$
with real $`T_n`$ and $`\stackrel{~}{T}_n`$. From the symmetry of the original TBA equations $`T`$ and $`\stackrel{~}{T}`$ both even functions of $`\theta `$ so that in our present case $`T_n=T_n`$ and $`\stackrel{~}{T}_n=\stackrel{~}{T}_n`$. One can read-off the leading large $`n`$ behavior of the coefficients $`T_n`$ and $`\stackrel{~}{T}_n`$ from the asymptotics (4.6)
$`T_n`$ $`()^n\sqrt{{\displaystyle \frac{bQ}{\pi }}}n^{2Qbn1/2}\left(e{\displaystyle \frac{mR}{8\pi bQ\mathrm{cos}(\pi p/2)}}\right)^{2Qbn}`$ (4.10)
$`\stackrel{~}{T}_n`$ $`()^n\sqrt{{\displaystyle \frac{Q}{\pi b}}}n^{2Qn/b1/2}\left(e{\displaystyle \frac{bmR}{8\pi Q\mathrm{sin}(\pi p/2)}}\right)^{2Qn/b}`$
So far the constructions were explicitely based on the integral equation (2.3). The rest of the section is more speculative. Suppose that for $`T(\theta )`$ and $`\stackrel{~}{T}(\theta )`$ constructed as in (4.4) we can find $`Q_\pm (\theta )`$ and $`\stackrel{~}{Q}_\pm (\theta )`$ which solve
$`Q_\pm \left(\theta +i\pi \right)+Q_\pm (\theta i\pi )`$ $`=T(\theta )Q_\pm (\theta )`$ (4.11)
$`\stackrel{~}{Q}_\pm \left(\theta +i\pi \right)+\stackrel{~}{Q}_\pm (\theta i\pi )`$ $`=T(\theta )\stackrel{~}{Q}_\pm (\theta )`$
and are the “Bloch waves” with respect to the periods of $`T`$ and $`\stackrel{~}{T}`$ respectively
$`Q_\pm \left(\theta +i\pi (1+a)/2\right)`$ $`=\mathrm{exp}\left(\pm 2i\pi P/Q\right)Q_\pm (\theta )`$ (4.12)
$`\stackrel{~}{Q}_\pm \left(\theta +i\pi (1a)/2\right)`$ $`=\mathrm{exp}\left(\pm 2i\pi P/Q\right)\stackrel{~}{Q}_\pm (\theta )`$
with some Floquet index $`P`$. Let them be normalized by the quantum Wronskians
$`Q_+(\theta +i\pi /2)Q_{}(\theta i\pi /2)Q_+(\theta i\pi /2)Q_{}(\theta +i\pi /2)`$ $`=1`$ (4.13)
$`\stackrel{~}{Q}_+(\theta +i\pi /2)\stackrel{~}{Q}_{}(\theta i\pi /2)\stackrel{~}{Q}_+(\theta i\pi /2)\stackrel{~}{Q}_{}(\theta +i\pi /2)`$ $`=1`$
Then formally
$$X(\theta )=Q_+(\theta )\stackrel{~}{Q}_+(\theta )+Q_{}(\theta )\stackrel{~}{Q}_{}(\theta )$$
(4.14)
solves eqs.(4.4) as well as the $`X`$-system (2.13).
Unfortunately at present I know no effective means to construct these $`Q`$-functions. Moreover, there are serious doubts that the objects satisfying both (4.11) and (4.12) can be constructed in any sense, at least at rational values of $`b^2`$. I hope to say something more definite on this point in close future.
## 5 Large $`Re\theta `$ asymptotics
Let me comment a little more about the $`Re\theta \mathrm{}`$ asymptotics (with $`Im\theta `$ fixed) of the function $`X(\theta )`$ in the whole complex plane. In principle it can be restored from the asymptotic (2.10) in the strip $`|Im\theta |\pi /2`$ using the functional relation (2.13) or, more conveniently, the relations (4.4) together with the asymptotics (4.6). The asymptotics is always of the form
$$X(\theta )\mathrm{exp}\left(A(Im\theta )\mathrm{exp}(Re\theta )\right)$$
(5.1)
with some complex function $`A(\eta )`$ of real variable $`\eta =Im\theta `$. Apparently $`ReA(\eta )`$ controls the rate of growth of the absolute value of $`X`$. At $`|\eta |<\pi /2`$ we have
$$A(\eta )=\frac{mR}{4\mathrm{sin}\pi p}e^{i\eta }$$
(5.2)
It follows from (4.4) that $`A(\eta )`$ satisfies two functional relations
$`A(\eta +\pi (1a)/2))`$ $`=\{{\displaystyle \genfrac{}{}{0.0pt}{}{A(\eta )+\sigma (\eta )\text{if }Re(A(\eta )+\sigma (\eta ))>ReA(\eta +\pi (1a)/2)}{A(\eta +\pi (1a)/2)\text{otherwise}}}.`$ (5.3)
$`A(\eta +\pi (1+a)/2))`$ $`=\{{\displaystyle \genfrac{}{}{0.0pt}{}{A(\eta )+\stackrel{~}{\sigma }(\eta )\text{if }Re(A(\eta )+\stackrel{~}{\sigma }(\eta ))>ReA(\eta +\pi (1+a)/2)}{A(\eta +\pi (1+a)/2)\text{otherwise}}}.`$
where the functions $`\sigma (\eta )`$ and $`\stackrel{~}{\sigma }(\eta )`$ control the asymptotics of $`T(\theta )`$ and $`\stackrel{~}{T}(\theta )`$. They are defined as
$`\sigma (\eta )`$ $`={\displaystyle \frac{mR}{4\mathrm{cos}(\pi p/2)}}\mathrm{exp}(i\eta i\pi (1p)/2)\text{at }0\eta <(1+a)/2`$ (5.4)
$`\stackrel{~}{\sigma }(\eta )`$ $`={\displaystyle \frac{mR}{4\mathrm{sin}(\pi p/2)}}\mathrm{exp}(i\eta i\pi p/2)\text{at }0\eta <(1a)/2`$
and continued outside these regions periodically as $`\stackrel{~}{\sigma }(\eta +\pi (1a)/2)=\stackrel{~}{\sigma }(\eta )`$ and $`\sigma (\eta +\pi (1+a)/2)=\sigma (\eta )`$. Both $`\sigma (\eta )`$ and $`\stackrel{~}{\sigma }(\eta )`$ jump by $`imR/2`$ at $`\eta =\pi (1+a)n/2`$, $`nZ`$ and $`\eta =\pi (1+a)n/2`$, $`nZ`$ respectively. This corresponds to the limiting density of zeroes prescribed by (4.7).
A common solution to (5.3) exists. Generally the solution is discontinuous at all values $`\eta =\pm \pi (m(1+a)/2+n(1a)/2)`$ with arbitrary positive integers $`n`$ and $`m`$. At each such point the imaginary part $`ImA`$ jumps down by $`mR/2`$ indicating an asymptotic line of accumulation of zeroes of $`X(\theta )`$, the asymptotic density being the same as that of the functions $`T`$ and $`\stackrel{~}{T}`$ (4.7). The real part of $`A`$ at these points is continuous itself but has discontinuities in the first derivative. The first Stokes line appears at $`\eta =\pi `$.
At large $`\eta `$ the structure is qualitetively different dependent on the arithmetic nature of parameter $`b^2`$. If it is a rational number the periods of $`T`$ and $`\stackrel{~}{T}`$ are commensurable. Some of the discontinuities merge forming multiple jumps in the imaginary part. The solution $`A(\theta )`$ bears a regular “quasiperiodic” structure with the common period of $`T`$ and $`\stackrel{~}{T}`$ . For irrational $`b^2`$ the periods are incommensurable and as $`\eta \mathrm{}`$ the singularities are more and more dense, the solution having quite irregular behavior.
In fig.1 real and imaginary parts of $`A(\eta )`$ are plotted for the simplest (self-dual) case $`b=1`$. Both periods are equal to $`\pi /2`$. The discontinuities are located at $`\eta =\pm \pi (n+1)/2`$, $`n=1,2,\mathrm{}`$ where the imaginary part jumps by $`mRn/2`$. Contrary to the location of zeroes of $`T`$, zeroes of $`X`$ cannot lie exactly on the lines $`\mathrm{Im}\theta =`$const, at least for the first line $`\eta =\pi `$. Indeed, in the case $`b^2=1`$ it follows from the functional relation that at real $`\theta `$
$$\left|X(\theta +i\pi )\right|^2=X^2(\theta )+T^2(\theta )$$
(5.5)
which is strictly positive.
Another rational situation $`b^2=1/2`$ corresponds to the periods $`\pi /3`$ and $`2\pi /3`$. Function $`A(\eta )`$ is plotted in fig.2. The structure is again quite regular, the discontinuities occuring at $`\eta =\pi +n\pi /3`$, $`n=1,2,\mathrm{}`$, the first two discontinuities in $`ImA`$ being $`mR/2`$, next two are twice of this amount, then next two trice, etc.
With irrational $`b^2`$ the picture is far less regular. To illustrate what happens when $`b^2`$ slightly deviates from a simple rational number, in fig.3 we plot $`A(\eta )`$ for $`b^2=0.8086\mathrm{}`$ which is reasonably close to $`1`$. Comparing with fig.1 we see that the first discontinuity at $`\eta =\pi `$ remains basically the same while the second (double) discontinuity at $`\eta =3\pi /2`$ splits in two simple ones, the third (triple) splits in three simple jumps etc. At some point these splitted groups come to overlap and the picture turns irregular.
## 6 Staircase situation
The staircase model is a formal analytic continuation of ShG to complex values of the parameter $`b`$ such that $`b^1=b^{}`$. Although the physical content of this continuation is not completely clear from the field theory point of view, the TBA equation (2.3) remains completely sensible and this continuation of (2.3) can be studied on its own footing. The effective central charge (2.2) is still real and develops sometimes quite intriguing patterns (see ).
In the complex $`b`$ has been parameterised as follows
$$b^2=\frac{1+2i\theta _0/\pi }{12i\theta _0/\pi }$$
(6.1)
with real $`0\theta _0<\mathrm{}`$. Parameter $`p`$ of (1.3) now reads
$$p=\frac{1}{2}+\frac{i\theta _0}{\pi }$$
(6.2)
while $`1p=p^{}`$ and $`a`$ defined in (2.6) is purely imaginary $`a=2i\theta _0/\pi `$. The TBA kernel (2.4) is real and reads
$$\phi (\theta )=\frac{1}{2\pi }\left(\frac{1}{\mathrm{cosh}(\theta +\theta _0)}+\frac{1}{\mathrm{cosh}(\theta \theta _0)}\right)$$
(6.3)
with the Fourier transform
$$\phi (\omega )=\frac{\mathrm{cos}(\omega \theta _0)}{\mathrm{cosh}(\pi \omega /2)}$$
(6.4)
After these substitution the integral equation (2.3) determines real-analytic functions $`\epsilon (\theta )`$, $`Y(\theta )`$ and, through (2.9), a real-analytic and symmetric $`X(\theta )`$ with the asymptotic behavior at $`Re\theta \mathrm{}`$ in the strip $`\pi /2<Im\theta <\pi /2`$
$$X(\theta )\mathrm{exp}\left(\frac{mR}{4\mathrm{cosh}\theta _0}\mathrm{exp}\theta \right)$$
(6.5)
The functional equation (2.13) reads now
$$X\left(\theta +i\pi /2\right)X\left(\theta i\pi /2\right)=1+X\left(\theta +\theta _0\right)X\left(\theta \theta _0\right)$$
(6.6)
All the considerations of sect.4 can be repeated literally. Function $`X(\theta )`$ is still entire as well as $`T(\theta )`$ and $`\stackrel{~}{T}(\theta )`$ defined in eq.(4.4) and the asymptotic (6.5) can be extended to the strip $`\pi <Im\theta <\pi `$. The difference is that the periods of $`T\left(\theta \right)`$ and $`\stackrel{~}{T}(\theta )`$
$`T(\theta +\tau )`$ $`=T\left(\theta \right);\tau =i\pi (1+a)/2=i\pi /2+\theta _0`$ (6.7)
$`\stackrel{~}{T}(\theta +\stackrel{~}{\tau })`$ $`=\stackrel{~}{T}(\theta );\stackrel{~}{\tau }=i\pi (1a)/2=i\pi /2\theta _0`$
are now complex $`\tau =\stackrel{~}{\tau }^{}`$. Functions $`T(\theta )=T\left(\theta \right)`$ and $`\stackrel{~}{T}(\theta )=\stackrel{~}{T}(\theta )`$ are still symmetric but no more real analytic. Instead
$$T^{}(\theta )=\stackrel{~}{T}(\theta ^{})$$
(6.8)
Expansions similar to (4.9)
$`T(\theta )`$ $`={\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}T_n\mathrm{exp}(2i\pi n\theta /\tau )`$ (6.9)
$`\stackrel{~}{T}(\theta )`$ $`={\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}\stackrel{~}{T}_n\mathrm{exp}(2i\pi n\theta /\stackrel{~}{\tau })`$
are convergent and $`T_n=T_n`$; $`\stackrel{~}{T}_n=\stackrel{~}{T}_n`$. Instead of being real as in the real $`b`$ case, these coefficients are complex conjugate $`T_n^{}=\stackrel{~}{T}_n`$.
In the strip $`0<Im\theta <\pi /2`$ the following $`Re\theta \mathrm{}`$ asymptotics holds for $`T`$ and $`\stackrel{~}{T}`$
$`T(\theta )`$ $`X(\theta +\stackrel{~}{\tau })/X(\theta )\mathrm{exp}\left({\displaystyle \frac{mR(1i\mathrm{exp}(\theta _0))}{4\mathrm{cosh}\theta _0}}e^\theta \right)`$ (6.10)
$`\stackrel{~}{T}(\theta )`$ $`X(\theta +\tau )/X(\theta )\mathrm{exp}\left({\displaystyle \frac{mR(1i\mathrm{exp}(\theta _0))}{4\mathrm{cosh}\theta _0}}e^\theta \right)`$
The asymptotic changes at the Stokes line along the real axis where
$`T(\theta )`$ $`2\mathrm{exp}\left({\displaystyle \frac{mR(1+i\mathrm{sinh}\theta _0)}{4\mathrm{cosh}\theta _0}}e^\theta \right)\mathrm{cos}\left({\displaystyle \frac{mR}{4}}e^\theta \right)`$ (6.11)
$`\stackrel{~}{T}(\theta )`$ $`2\mathrm{exp}\left({\displaystyle \frac{mR(1i\mathrm{sinh}\theta _0)}{4\mathrm{cosh}\theta _0}}e^\theta \right)\mathrm{cos}\left({\displaystyle \frac{mR}{4}}e^\theta \right)`$
and we observe again an infinite sequence of zeroes accumulating at infinity, the density being the same as in the real $`b`$ case of sect.5. In the strips $`\pi n/2<`$Im$`\theta <\pi (n+1)/2`$ (with arbitrary integer $`n`$) the Re$`\theta \mathrm{}`$ asymptotics follow from the periodicity (6.7). In particular, along the lines Im$`\theta =in\pi /2`$ (with any integer $`n`$) we’ll have
$`T(\theta )`$ $`2\mathrm{exp}\left({\displaystyle \frac{mR(1+i\mathrm{sinh}\theta _0)}{4\mathrm{cosh}\theta _0}}e^{Re\theta n\theta _0}\right)\mathrm{cos}\left({\displaystyle \frac{mR}{4}}e^{Re\theta n\theta _0}\right)`$ (6.12)
$`\stackrel{~}{T}(\theta )`$ $`2\mathrm{exp}\left({\displaystyle \frac{mR(1i\mathrm{sinh}\theta _0)}{4\mathrm{cosh}\theta _0}}e^{Re\theta +n\theta _0}\right)\mathrm{cos}\left({\displaystyle \frac{mR}{4}}e^{Re\theta +n\theta _0}\right)`$
Functions $`\sigma (\eta )`$ and $`\stackrel{~}{\sigma }(\eta )`$ which control the asymptotics of $`T`$ and $`\stackrel{~}{T}`$ at $`Re\theta \mathrm{}`$
$`T(\theta )`$ $`\mathrm{exp}\left(\sigma (Im\theta )e^{Re\theta }\right)`$ (6.13)
$`\stackrel{~}{T}(\theta )`$ $`\mathrm{exp}\left(\stackrel{~}{\sigma }(Im\theta )e^{Re\theta }\right)`$
are plotted in fig.4 for the case $`\theta _0=1`$. It’s enough to present them for $`\eta 0`$ since $`\sigma (\eta )=\stackrel{~}{\sigma }^{}(\eta )`$. The imaginary part of $`\sigma (\eta )`$ jumps at the points $`\eta =n\pi /2`$ by the amount $`mR\mathrm{exp}(n\theta _0)/2`$, in accord with the density of zeroes predicted by (6.12). Unlike the previously considered case of real $`b`$, in the staircase situation the zeroes of $`T`$ and $`\stackrel{~}{T}`$ are not located exactly at the lines $`Im\theta =in\pi /2`$ but slightly shifted in the imaginary direction (we’ll observe this deviation numerically in the next section) and approach these lines asymptoticlally as $`Re\theta \mathrm{}`$.
Notice also that e.g. $`T(\theta )`$ is a single valued function of $`\xi =\mathrm{exp}(2i\pi \theta /\tau )`$. In the complex plane of this variable the asymptotic lines of accumulation of zeroes $`Im\theta =0`$; $`Re\theta \pm \mathrm{}`$ are parts of the spiral $`\left|\xi \right|=\mathrm{exp}(\pi \mathrm{arg}\xi /2\theta _0)`$ near which zeroes become dense at $`\left|\xi \right|\mathrm{}`$ or $`\left|\xi \right|0`$, the density growing as $`\left|\xi \right|^{\pm (1/4+\theta _0^2/\pi ^2)}`$ respectively. Therefore the large (or small)$`\left|\xi \right|`$ asymptotics of $`T(\xi )`$ at fixed $`\mathrm{arg}\xi `$ is rather complicated.
Large $`Re\theta `$ asymtotic of $`X(\theta )`$ at fixed $`\eta =Im\theta `$ is controlled by the function $`A(\eta )`$ (see eq.(5.1)). An example corresponding to the case $`\theta _0=1`$ is presented in fig.5. At the points $`\eta =\pm \pi (n+1)/2`$, $`n=1,2,\mathrm{}`$ imaginary part of $`A(\eta )`$ has discontinuities equal to $`mR\mathrm{sinh}(n\theta _0)/(2\mathrm{sinh}\theta _0)`$. These amounts determine the asymptotic density of zeroes of $`X(\theta )`$ along these lines.
## 7 Numerics
Integral equation (2.3) can be easily solved numerically e.g., by iterations. The iterations happen to be well convergent (the convergence is somewhat slower if $`R`$ approaches to $`0`$ or the parameter $`p`$ is taken very small). In the strip $`|Im\theta |<\pi /2`$ function $`X(\theta )`$ can be computed using the integral representation (2.9). This allows to continue $`X(\theta )`$ to the whole complex plane iterating the relation (2.13) (in fact at large Im$`\theta `$ it is more convenient to evaluate first $`T(\theta )`$ inside its period and then use (4.4)). In the rest of this section we will use the logarithmic scale parameter $`x=\mathrm{log}(mR/2)`$ instead of $`R`$.
1. Self-dual point $`𝐛^2=\mathrm{𝟏}`$.
In fig.6 several examples of function $`X(\theta )`$ on the real axis of $`\theta `$ are plotted for different values of $`x`$. Function is typically bell-shaped. As $`x`$ becomes large negative the width of the support of the bell as well as its height grow proportionally to $`x`$. No plateau typical for perturbed rational CFT’s is developed.
Few samples of $`T(\theta )`$ (which is the same as $`\stackrel{~}{T}(\theta )`$ at the self-dual point) are presented in fig.7. At $`x`$ negative and large enough, $`T`$ develops a plateau in the “central region” $`x<\theta <x`$ of the height which approaches slowly to $`2`$ as $`x`$ grows. We’ll comment more about this approach below. Outside the central region it starts to oscillate with growing amplitude and friquency. Approach to the asymptotic (4.7) is very fast.
Due the the symmetry $`T(\theta )=T(\theta )`$ this function is real on the imaginary axis too. A couple of examples are plotted in fig.8. At $`x`$ essentially negative, when the plateau is well developped in the central region, the mean value $`T_0`$ (see eq.(4.9)) is very close to the plateau height, the oscillations around (determined mainly by $`T_1`$) being very small ($`T_1R^4`$, see eq.(8.9) below).
Function $`T(\theta )`$ is real also at the half-period line $`Im\theta =\pi /4`$. Again there is a plateau region (at large negative $`x`$) of the same height as on the real axis. Then $`T(\theta +i\pi /4)`$ remains positive and grows following the asymptotic (4.8) (see fig.9). Numerical comuputations in the whole period strip $`0`$Im$`\theta <\pi /2`$ show no sign of other zeros then those on the real axis.
In fig.10 few examples of $`\left|X(\theta )\right|`$ in the complex plane at different values of $`Re\theta `$ are plotted vs. $`Im\theta `$ (for $`x=0`$). The specific values $`Re\theta =1.22`$ and $`Re\theta =2.25`$ are chosen close to the positions of the first two zeros of $`T(\theta )`$ on the real axis. The deeps near $`Im\theta =\pi `$ and $`Im\theta =3\pi /2`$ indicate a presence of zeros of $`X(\theta )`$ nearby.
More precise positions of zeros of $`X(\theta )`$ near the line $`Im\theta =\pi `$ (for the same value $`x=0`$) are examplified in fig.11. In fact all these zeros are inside the strip $`|Im\theta |<\pi `$. Only the first zero deviates noticably from the line $`Im\theta =\pi `$. The imaginary parts of next zeros are already very close to $`\pi `$ and tend to this value very fast.
2. Rational points. As an example of a rational point we take the simplest case $`b^2=1/2`$. The periods of $`T(\theta )`$ and $`\stackrel{~}{T}(\theta )`$ are commensurable and equal to $`2i\pi /3`$ and $`i\pi /3`$ respectively. In fact in this case there is no need to study separately $`T`$ and $`\stackrel{~}{T}`$ since, as it is readily derived from (2.13), they are bound up by the relation
$$\stackrel{~}{T}(\theta )=T(\theta )T(\theta +i\pi /3)2$$
(7.1)
It should be noted that similar finite degree functional relations between $`T`$ and $`\stackrel{~}{T}`$ exist for any rational $`b^2`$. For example, at $`b^2=1/3`$ the periods of $`T`$ and $`\stackrel{~}{T}`$ are $`3i\pi /4`$ and $`i\pi /4`$ respectively and
$$\stackrel{~}{T}(\theta )=T(\theta )T(\theta +i\pi /4)T(\theta +i\pi /2)T(\theta )T(\theta +i\pi /4)T(\theta +i\pi /2)$$
(7.2)
Numerical patterns of $`T\left(\theta \right)`$, $`\stackrel{~}{T}(\theta )`$ and $`X(\theta )`$ are essentially the same as for $`b^2=1`$. I’d only show few plots of $`|X(\theta )|`$ in the complex plane along the lines Im$`\theta =\pi `$, $`4\pi /3`$, $`5\pi /3`$, $`2\pi `$, $`7\pi /3`$ etc. to illustrate the mechanism of multiplication of the zeros density in the asymptotics $`Re\theta \mathrm{}`$ as required by the prediction of fig.2. In fig.12 $`|X(\theta )|`$ is plotted at $`Re\theta `$ in the vicinity of the real position of the first zero in $`T(\theta )`$ at $`\theta =1.2241\mathrm{}`$ (the case $`x=0`$ is taken as an example). At $`Im\theta =\pi `$ and $`4\pi /3`$ the picture indicates simple zeros located closely to this point in $`Re\theta `$ and slightly displaced in the imaginary direction. At $`Im\theta =5\pi /3`$ and $`2\pi `$ the zeros are splitted in two closely located ones again near the same position in the real direction. For $`Im\theta =7\pi /3`$ and $`8\pi /3`$ there are triplets of close zeros, etc. In fig.13 the same is examplified near the next zero of $`T(\theta )`$ at $`\theta =2.2527\mathrm{}`$. It is seen already that the scale of splitting becomes very small with $`Re\theta `$ growing and such zero multiplets look like multiple zeros if the numerical resolution is not enough.
3. General real $`𝐛^2`$. In general the periods of $`T(\theta )`$ and $`\stackrel{~}{T}(\theta )`$ are incommensurable. As it was mentioned in sect.5, it leads in particular to quite complicated $`Re\theta \mathrm{}`$ asymptotics at sufficiently large $`Im\theta `$. While the analytic structure of $`T(\theta )`$ and $`\stackrel{~}{T}(\theta )`$ remains essentially as described above (in particular, I verified for many examples that all zeroes of $`T`$ and $`\stackrel{~}{T}`$ are on the real axis), the structure of zeros in $`X(\theta )`$ becomes, as $`Im\theta `$ comes essentially large, rather chaotic. I hope to turn again to this point in future studies. Let me mention only the following observation conserning the small $`R`$ (or large negative $`x`$) picture. If $`x1`$, in the central region $`x<\theta <x`$ function $`X(\theta )`$ matches extremely well the following expression
$$X(\theta )=\frac{\mathrm{cos}(2QP\theta )}{\left[\mathrm{sinh}(2\pi bP)\mathrm{sinh}(2\pi P/b)\right]^{1/2}}$$
(7.3)
where $`P`$ is an $`R`$-dependent parameter. Roughly it can be estimated from the requirement that $`X(\theta )0`$ at $`\theta =\pm x`$, i.e., $`P=\pi /(4Qx)+O(x^2)`$. Notice that substituting of this approximation to the expression of the effective central charge
$$c_{\text{eff}}=124P^2$$
(7.4)
(see ref. for the motivations) we arrive just at the leading UV logarithmic correction (2.16). It is easy to verify that expression (7.3) *satisfies exactly* the functional $`X`$-system (2.13). This means in particular that it remains a valid approximation (at large $`x`$) of $`X(\theta )`$ in the whole complex strip $`x<Re\theta <x`$.
In fact along the lines of and a far better estimate of $`P`$ can be found which takes into account all logarithmic in $`R`$ corrections to (2.16). In this framework $`P`$ is determined as the first root of the transcendental equation
$$4QP\mathrm{log}\left(R/2\pi \right)+i\mathrm{log}(S_L(P))=\pi $$
(7.5)
where $`S_L(P)`$ is the so-called Liouville reflection amplitude (for the arguments see )
$$S_L(P)=\left(\pi \mu \frac{\mathrm{\Gamma }(b^2)}{\mathrm{\Gamma }(1b^2)}\right)^{2iP/b}\frac{\mathrm{\Gamma }(1+2ibP)\mathrm{\Gamma }(1+2iP/b)}{\mathrm{\Gamma }(12ibP)\mathrm{\Gamma }(12iP/b)}$$
(7.6)
For expample, in fig.14 the shape of $`X(\theta )`$ is compared with the approximation (7.3) for $`x=6`$ and two values of the parameter $`b^2=0.8086\mathrm{}`$ and $`b^2=0.1201\mathrm{}`$ According to (7.6) they correspond to $`P=0.05332\mathrm{}`$ and $`P=0.04052\mathrm{}`$ respectively.
In view of (7.3) the plateau heights of $`T(\theta )`$ and $`\stackrel{~}{T}(\theta )`$ in the central region $`x<\theta <x`$ and $`x\mathrm{}`$ can be estimated as
$`T_0`$ $`=2\mathrm{cosh}(2\pi bP)`$ (7.7)
$`\stackrel{~}{T}_0`$ $`=2\mathrm{cosh}(2\pi P/b)`$
with the same $`P`$ determined by (7.6).
4. Complex (staircase) values of $`𝐛^2`$. The staircase version of TBA equation (2.3) (with the kernel (6.3)) is solved numericaly in the same way as in the ShG case. The structure of the solution has some interesting differences from the case of real $`b`$. The peculiarities are more manifested if parameter $`\theta _0`$ in (6.3) is taken sufficiently large and the deep UV region $`x1`$ is considered.. At $`\theta _01`$ parameter $`b`$ is close to imaginary unity and
$$Q=\frac{\pi }{\sqrt{\theta _0^2+\pi ^2/4}}$$
(7.8)
is small.
In fig.15 two pictures of $`X(\theta )`$ on the real axis are plotted for $`x=46`$ and $`x=91`$, both at $`\theta _0=20`$. The characteristic staircase behavior (in fact very similar to that of $`Y(\theta )`$ observed in ref.) is very apparent. It is seen also that expression (7.3) with the parameter $`P`$ determined as the solution to the staircase version of (7.5) (in the present case $`P=0.06680\mathrm{}`$ and $`P=0.04176\mathrm{}`$ respectively) still follows in the central region the average behavior of the solution. The deviations (or in other words the corrections to the approximation (7.3)) are now oscillating and much bigger then in the ShG case. At sufficiently large $`x`$ (many amounts of $`\theta _0`$) the ascending part (starting from $`\theta =x\theta _0`$) of the staircase consists of a succession of almost flat steps of constant width $`\theta _0`$, the heights being very well fitted by the expression
$$X_n=\frac{\mathrm{sin}(2QPn/\theta _0)}{\mathrm{sin}(2QP/\theta _0)}$$
(7.9)
with $`n=0,1,2,\mathrm{}`$. It should be noted that for the reasons to be explained just below, in the case of complex $`b`$ expression (7.3) does not give an approximation in the whole strip $`x<Re\theta <x`$. Although (7.3) is still an exact solution to the staircase $`X`$-system (6.6), its validity is restricted to a certain parallelogram in the complex $`\theta `$-plane.
To see this in fig.16 I present the location of zeros of the function $`X(\theta )`$ (at $`\theta _0=20`$ and $`x=46`$) in the upper half-plane. The picture is obviously reflected to the lower half-plane by the symmetry of $`X(\theta )`$. We see several strings of zeros accumulating along the lines $`Im\theta =\pi n/2`$, $`n=2,3,4,\mathrm{}`$ which start at the values $`Re\theta _2=47.1500\mathrm{}`$, $`Re\theta _3=27.12267\mathrm{}`$ $`Re\theta _2\theta _0`$, $`Re\theta _4=7.11742\mathrm{}`$ $`Re\theta _3\theta _0`$ etc. There are opposite sets of strings symmetric with respect of the imaginary axis. As $`n`$ grows they are shifted in steps of $`\theta _0`$ and $`\theta _0`$ respectively and meet each other at a certain value of $`n`$ ($`n=5`$ in the present example). After that, continuous lines of dense zeroes extending from $`\mathrm{}`$ to $`\mathrm{}`$ are formed. In general at large $`\theta _0`$ the first such line is already dense enough to produce the effect of a “cut” where the behavior of $`X(\theta )`$ changes drastically. In our example, across the first “cut” at $`Im\theta =5\pi /2`$ the absolute value of $`X(\theta )`$ jumps by many ($`10^5`$) orders of magnitude. This is examplified in fig.17 where $`X(\theta )`$ is plotted along the imaginary axis. Notice that before $`Im\theta =5\pi /2`$ function $`X(\theta )`$ is almost constant. As $`\theta _0`$ grows this effect becomes more and more dramatic and in the limit $`\theta _0\mathrm{}`$ the lines of zeros become real cuts. In the forthcoming paper I’m going to comment more on this effect which plays a crucial role in the analytic connection between the staircase behavior at finite $`\theta _0`$ and the “sin-Gordon” solutions corresponding to purely imaginary $`b=i\beta `$.
In fig.18 location of zeros of the periodic function $`T(\theta )`$ is shown at the “moderate” value $`\theta _0=1`$ and different $`x`$. The picture is presented here just to demonstrate two observations: (a) At staircase values of $`b`$ zeros of $`T(\theta )`$ are not exactly on the real axis but displaced slightly to the complex plane, the displacement becoming negligible very quickly with the number of zero. (b) At $`x\mathrm{}`$ the picture of zeros is “frozen” in the sense that at $`x`$ sufficiently large the pattern of zeros is simply shifted by the amount $`\mathrm{\Delta }x`$ in the $`\theta `$ plane as one changes $`xx\mathrm{\Delta }x`$. As it is usual in the TBA practice, it is convenient to study these frosen limiting patterns substituting the original “massive” TBA equation (2.3) by the massless version of it. This is one of the motivations for the subsequent study (to be published ).
## 8 Concluding remarks
* In the present study I didn’t touch at all the important question about the $`R`$ (or $`x`$) dependence of the effective central charge (2.2) determined through the finite-size ground state energy (2.1). The UV behavior at $`x\mathrm{}`$ is especially interesting since the analytic structure of $`c_{\text{eff}}(R)`$ is quite unusual (like in (2.16)). The Liouville quantization condition (7.5) together with (7.4) proves to be a very good approximation to the UV effective central charge behavior. However, while it takes into account all the “soft” (logarithmic in $`R`$) contrubutions to the asymptotic, there are certainly power-like corrections in $`R`$. The most important of them (at least at real $`b`$) is the contribution of the ground state energy (1.9). Approximation (7.4) is essentially impruved if it is taken into account
$$c_{\text{eff}}=124P^2+3(mR)^2/(4\pi \mathrm{sin}\pi p)$$
(8.1)
Usually after the ground state energy contibution is subtracted the reminder is a series in the in the perturbative powers of $`R`$ like in (2.15). In our present case with Lagrangian (1.1) naively one could expect a series like
$$c_{\text{eff}}3(mR)^2/(4\pi \mathrm{sin}\pi p)=124P^2+\underset{n=1}{\overset{\mathrm{}}{}}c_n(P)\left(\mu R^{2+2b^2}\right)^{2n}$$
(8.2)
where the powers of $`R`$ are predicted on the dimensional arguments and the coefficients $`c_n(P)`$ are something like the Coulomb gaz perturbative integrals corresponding to expansion in $`\mu `$ around the vacuum of momentum $`P`$ (therefore they keep some smooth $`R`$ dependence). Although the leading correction $`n=1`$ of order $`R^{4+4b^2}`$ is likely in agreement with the numerical data (at least at sufficiently small $`b^2`$), the whole structure (8.2) certainly contradicts duality. Correct expansion must contain also the “dual” corrections with powers $`R^{4+4/b^2}`$. Of course many people immediately propose to add the dual interaction to the Lagrangian and consider an action like
$$A_{\mathrm{trivial}}=\left[\frac{1}{4\pi }(_a\varphi )^2+2\mu \mathrm{cosh}(2b\varphi )+2\stackrel{~}{\mu }\mathrm{cosh}(2b^1\varphi )\right]d^2x$$
(8.3)
with $`\stackrel{~}{\mu }`$ taken from (1.8) or sometimes introduced as an independent coupling. Simultaneous expansion in both couplings would supposedly lead to a self-dual series
$$c_{\text{eff}}3(mR)^2/(4\pi \mathrm{sin}\pi p)=124P^2+\underset{m,n=1}{\overset{\mathrm{}}{}}c_{m,n}(P)\mu ^{2m}\stackrel{~}{\mu }^{2n}R^{4Q(mb+n/b)}$$
(8.4)
Coefficients $`c_{m,n}(P)`$ are computed as the mixed Coulomb gaz integrals which incude both kinds of charges produced by the expansion of action (8.3). Although the general structure of (8.4) looks very likely, to my sense this scenario (which I call trivial) with naive addition of dual interactions (as in (8.3)) is not exactly in the spirit of duality. However at present it does not contradict any data and in fact should be verified. A check of this trivial scenario (as well as any other one) requires very subtle measurements of the subleading power-like corrections to the effective central charge as well as tedious calculations of mixed perturbative integrals. I understand that this is a *quantitative* work which hardly can be replaced by general speculations.
* Of course the definition of $`P`$ as the solution to the quantization condition (7.5) is not completely unambiguous. The power-like in $`R`$ corrections (which are exponentially small in $`1/P`$) can be arbitrarily redistributed between the expression for the observable effective central charge (8.4) and the formulation of the quantization condition. In other words one can add exponentially small in $`1/P`$ corrections to (7.6) and consider this as a new definition of $`P`$. The problem is that at present the parameter $`P`$ is not precisely observable, i.e., it cannot be directly measured in the TBA calculations (apart from the abovementioned *definition* through the observable finite-size effective central charge).
* At this point we arrive at the most intriguing question touched only slightly in the present study. This is about the possibility to construct the solutions $`Q_\pm (\theta )`$ and $`\stackrel{~}{Q}_\pm (\theta )`$ of eqs.(4.11) with the properties (4.12) and (4.13) such that the solution to (2.3) can be built as the combination (4.14). Was this be possible we’d have another unambiguous definition of the parameter $`P`$ as the Floquet index in (4.12). However, there are serious doubts that such solutions exist in any sense, at least for rational values of $`b^2`$. To clarify this point, in the next publication I’ll consider the massless version of the ShG TBA equation where the parameter $`P`$ in introduced from the very beginning instead of the scale parameter $`R`$. In this context the solutions (4.12) can be found at least as formal series for irrational values of $`b^2`$ in the way that the construction (4.14) can be given an exact sense.
* In connection with the periodic structures encoded in the functions $`T(\theta )`$ and $`\stackrel{~}{T}(\theta )`$, it seems quite interesting to understand better the $`R`$ dependence of the coefficients $`T_n`$ and $`\stackrel{~}{T}_n`$ in the expansions (4.9) or (6.9). I checked numerically the $`R0`$ asymptotic of $`T_0`$ and $`\stackrel{~}{T}_0`$. The leading asymptotics of $`T_0`$ and $`\stackrel{~}{T}_0`$ (remember that for definiteness in this study I always take $`b1`$, in particular the data discussed in this item were calculated at $`b=0.3466`$) are very well fitted by the expressions (7.7) with $`P`$ the solution to (7.5). The correction to (7.7) for $`T_0`$ can be set in the form
$$T_0=2\mathrm{cosh}(2\pi bP)+(mR/2\pi )^{4bQ}T_0^{(1)}(P)$$
(8.5)
where $`T_0^{(1)}(P)`$ is a smoothly varying function of $`P`$ with certain UV limit $`T_0^{(1)}(0)`$. As for $`\stackrel{~}{T}_0`$ the correction is better fitted as
$$\stackrel{~}{T}_02\mathrm{cosh}(2\pi P/b)A(P)(mR/2\pi )^\alpha $$
(8.6)
with $`\alpha `$ again numerically close to $`4bQ`$. The common TBA experience would expect from the periodic structure of $`\stackrel{~}{T}(\theta )`$ a much smaller correction $`(mR/2\pi )^{4Q/b}`$. A probable explanation is that there are some power-like corrections to the Liouville quantization condition (7.5) so that the “correct” value of $`P_{\text{correct}}`$ is off from $`P`$ (calculated from (7.5)) by an amount of order $`(mR/2\pi )^{4bQ}`$
$$P=P_{\text{correct}}+P_{1,0}(mR/2\pi )^{4bQ}+\mathrm{}$$
(8.7)
With this $`P_{\text{correct}}`$ an asymptotic
$$\stackrel{~}{T}_0=2\mathrm{cosh}(2\pi P_{\text{correct}}/b)+(mR/2\pi )^{4Q/b}\stackrel{~}{T}_0^{(1)}(P_{\text{correct}})+\mathrm{}$$
(8.8)
must hold. If the power $`4Q/b4bQ`$ (like in the present experiment with $`4Q/b=37.31\mathrm{}`$ $`4bQ=4.480\mathrm{}`$) we can even try to relate the coefficient $`A`$ in (8.6) to the leading power correction to (7.5). As the corrections in (8.8) are much more suppressed at $`R0`$ this seems reasonable.
* It is easy to verify that the leading asymptotic of the higher coefficients $`T_n`$ and $`\stackrel{~}{T}_n`$ in the expansion (4.9) are of the form
$`T_n`$ $`=(mR/2\pi )^{2bQ|n|}T_n^{(0)}(P)+\mathrm{}`$ (8.9)
$`\stackrel{~}{T}_n`$ $`=(mR/2\pi )^{2Q|n|/b}\stackrel{~}{T}_n^{(0)}(P)+\mathrm{}`$
with $`T_n^{(0)}(P)`$ and $`\stackrel{~}{T}_n^{(0)}(P)`$ regular at $`P=0`$. I analysed quantitatively the functions $`T_1^{(0)}(P)`$ and $`T_1^{(0)}(P)`$. Motivated by the constructions of refs. let us introduce slightly rescaled functions
$`t_1(P)`$ $`=T_1^{(0)}(P)\left(Z(p)\right)^{2bQ}`$ (8.10)
$`\stackrel{~}{t}_1(P)`$ $`=\stackrel{~}{T}_1^{(0)}(P)\left(Z(p)\right)^{2Q/b}`$
with $`Z(p)`$ defined in (1.7). In table 1 the values of $`t_1(P)`$ and $`\stackrel{~}{t}_1(P)`$ are compared with the following analytic expressions borrowed from refs. where these coefficients enter the explicit constructions of the “sin-Gordon” (i.e., related to purely imaginary values of $`b`$) analogs of $`T(\theta )`$
$`t_1^{\text{CFT}}(P)`$ $`={\displaystyle \frac{4\pi ^2\mathrm{\Gamma }(1+2b^2)}{\mathrm{\Gamma }^2(b^2)\mathrm{\Gamma }(1+b^2+2ibP)\mathrm{\Gamma }(1+b^22ibP)}}`$ (8.11)
$`\stackrel{~}{t}_1^{\text{CFT}}(P)`$ $`={\displaystyle \frac{4\pi ^2\mathrm{\Gamma }(1+2b^2)}{\mathrm{\Gamma }^2(b^2)\mathrm{\Gamma }(1+b^2+2iP/b)\mathrm{\Gamma }(1+b^22iP/b)}}`$
The TBA numbers are measured at $`b=0.3465545.\mathrm{}`$ and different values of $`x`$. The same for the staircase example $`\theta _0=10`$ is presented in table2. The convergence is noticably slower due to much weeker suppression of the higher power corrections (Re$`(4+4b^2)=0.1926\mathrm{}`$ in this case). The numbers quoted make it clear that an analytic continuation of the constructions of refs. for the ShG or staircase values of the parameter is quite relevant. The corresponding “CFT integrable structures” will be shown to have a precise relation to the solutions of the massless versions of ShG and staircase TBA equations where the parameter $`P=P_{\text{correct}}`$ is fixed by construction.
* I’d like to mention a quite intriguing recent article where the author arrived at the function (2.9) in a rather different context. It appears as the exact wave function of the finite-size sinh-Gordon model in special “$`\gamma `$-coordinates” which are the sinh-Gordon version of the Flaschka-McLaughlin variables (see for the details). This gives again a motivation to continue the study of the analytic structures related to the ShG TBA equation.
Acknowledgments. I acknowledge very useful discussions with C.Ahn, V.Bazhanov, V.Fateev, C.Rim and mostly with A.Zamolodchikov. I thank also P.Wiegmann who brought my attention to papers and introduced to me the world of incommensurable periods. Discussions with G.Mussardo about the “double sin-Gordon” were also quite relevant. My special gratitude is to S.Lukyanov who communicated me his article before publication and encouraged my work by exciting discussions.
The work is supported in part by EU under the contract ERBFMRX CT 960012.
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